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warning/0507/cond-mat0507102.html | ar5iv | text | # Extended hydrodynamics from Enskog’s equation The bidimensional casepublished in Physica A, 354, 77-87, 2005
## 1 Introduction
Today Enskog’s original kinetic theory is known as the standard Enskog theory (SET) because after the pioneer work of van Beijeren and Ernst there are several new versions of Enskog’s theory collectively called revised Enskog’s theory (RET) . Among the latter there are versions that have been extended to describe condensed matter . To Navier-Stokes level both SET and RET lead to the same results , whether or not an external force is present.
In the present article we make use of extended hydrodynamic equations for the bidimensional case . They are more complete than a linear approximation but still they are the result of an approximation scheme that we explain elsewhere. Using a strategy as in Ref. and approximations defined in Ref. we obtain in Sec. 5 the same hydrodynamic equations for SET and RET.
In this article we apply our extended hydrodynamics to a one dimensional steady heat conductive state. There is much work on this as, for example, the experimental results in or the theoretical ones in . Recently Kim and Hayakawa studied this problem for hard core and Maxwellian particles using Boltzmann’s equation combined with Chapman-Enskog’s method. They tried a test and criticized the analysis of the nonequilibrium steady state thermodynamics (SST) proposed by Sasa and Tasaki . In the last reference the authors state that if there is gas in a one dimensional heat conductive configuration in contact, through a porous wall, with an equilibrium gas state, then a pressure difference must appear in the direction of the heat flow. We analyze this double system making use of the extended hydrodynamic equations derived from Enskog’s equation using Grad’s moment expansion method . Our conclusions differ from those in .
The organization of the present article is as follows. In Sec. 2 the configuration of these systems is drawn schematically, in Sec. 3 the condition for the two systems to be in contact via the central porous plate is introduced: the upper and lower plates are normal plates; the central plate has many small pores through which the gas can pass. In Sec. 4 we give the basic equations used in this paper. Comments are in Sec. 5. Finally, our discussion and conclusions are written in Sec. 6 and 7, respectively.
## 2 Definition of the system
Sasa and Tasaki proposed an interesting system consisting of a nonequilibrium steady state subsystem in contact with a subsystem in equilibrium as explained below. This system has three plates as shown in Fig. 1 and there is gas between them. The upper and lower plates (plates 1 and 3) are normal plates. The central plate (plate 2) has pores through which gas can pass.
Following Sasa and Tasaki, we consider the system consisting of three infinite parallel plates 1, 2 and 3 separated by a distance $`L`$. The $`Y`$ axis is defined perpendicular to them while an $`X`$ axis is placed on plate 2. The pores in plate 2 are distributed homogeneously. Plates 1 and 2 have fixed temperature $`T_1`$ while plate 3 has a different (fixed) temperature $`T_2`$.
After a sufficiently long time, by effusion, some of the gas passes through plate 2 and the gas between plates 1 and 2 reaches an equilibrium state. The system between the plates 2 and 3 reaches a nonequilibrium steady state with translation symmetry along the $`X`$ axis.
We assume that the typical distances between pores is very small and that the diameter of the pores is also very small, so that the ratio between such lengths and the mean free path is much smaller than unity. Having no external force and no hydrodynamic velocity there is no heat flux parallel to the plates. The system is in a static configuration.
## 3 The contact condition
In general, there is a difference between the temperatures of the plates and the temperatures of the gas in contact with them, this is a well-known effect called thermal slip. However, for simplicity sake, we assume that the temperature of plate 2 and the gas in contact on both sides of it are equal, namely we are neglecting the Knudsen layer.
The velocity and the peculiar velocity of the gas will be denoted by $`𝐜`$ and $`𝐂`$, respectively. The condition that there is no mass flux through plate 2 is
$$_{\mathrm{}}^{\mathrm{}}𝑑c_x_0^{\mathrm{}}𝑑c_yc_yf_{\mathrm{equil}.}+_{\mathrm{}}^{\mathrm{}}𝑑c_x_{\mathrm{}}^0𝑑c_yc_yf_{y=0}=0,$$
(1)
where $`f_{\mathrm{equil}.}=n_{\mathrm{𝖾𝗊}}\left[\frac{1}{2\pi T_1}\right]\mathrm{exp}\left[\frac{C^2}{2T_1}\right]`$ is the equilibrium distribution function associated to the gas between plates 1 and 2 and $`n_{\mathrm{𝖾𝗊}}`$ is the uniform number density in this same region while $`f_{y=0}`$ is the nonequilibrium distribution function between plates 2 and 3 evaluated at $`y=0`$. Next, it is necessary to see how the two distributions satisfy condition (1).
Imposing condition (1) yields
$$n_{\mathrm{𝖾𝗊}}=\frac{1}{2}\left[n(0)+P_{yy}\right]\delta \frac{\chi }{2\mathrm{K}\mathrm{n}}\left[n(0)+P_{yy}\right]n(0)+\delta ^2\frac{\chi ^2}{2\mathrm{K}\mathrm{n}^2}\left[n(0)+P_{yy}\right]n(0)^2.$$
(2)
In addition, the total mass conservation law for the system is
$$n_{\mathrm{𝖾𝗊}}+_0^1n(y)𝑑y=2.$$
(3)
Above we are using dimensionless fields and dimensionless variables in general. The fields $`n`$ (number density), $`P_{ij}`$ (pressure tensor), $`\stackrel{}{Q}`$ (net heat flux vector), and $`T`$ (temperature) generally depend on the coordinate $`y`$, where $`p_{ij}`$ and $`𝐪^k`$ are the symmetric and traceless part of the pressure tensor and the kinetic part of the heat flux vector respectively. These hydrodynamic fields are defined according to the following sum rules:
$`{\displaystyle f𝑑𝐜}`$ $`=`$ $`n(y),`$ (4)
$`{\displaystyle 𝐂(y)f𝑑𝐜}`$ $`=`$ $`0,`$ (5)
$`{\displaystyle \frac{1}{2}C(y)^2f𝑑𝐜}`$ $`=`$ $`n(y)T(y),`$ (6)
$`{\displaystyle C_i(y)C_j(y)f𝑑𝐜}`$ $`=`$ $`n(y)T(y)\delta _{ij}+p_{ij}(y),`$ (7)
$`{\displaystyle \frac{1}{2}C(y)C^2(y)f𝑑𝐜}`$ $`=`$ $`q_y^\mathrm{k}(y).`$ (8)
We also use the following dimensionless numbers
$$\mathrm{K}\mathrm{n}=\frac{8\sqrt{2}}{\pi }\frac{\mathrm{}}{L},\text{Knudsen number},\delta =\frac{\sigma }{L}=\mathrm{K}\mathrm{n}\rho _0$$
where $`\sigma `$ is the particle’s diameter, $`\mathrm{}`$ their mean free at equilibrium and $`\rho _0`$ is the mean area density.
## 4 Balance equations
The basic concrete equations solved here are the following.
$``$ In the case of the linearized Boltzmann-Grad method (LBG): $`P_{yy}(y)P_{yy}=\text{constant}`$, $`P_{xy}(y)P_{xy}=\text{constant}`$, $`Q_y(y)Q_y=\text{constant}`$,
$$n(y)T(y)=P_{xx}=P_{yy},\frac{dT(y)}{dy}+\frac{2Q_y}{\mathrm{K}\mathrm{n}\sqrt{\pi \mathrm{T}(\mathrm{y})}}=0.$$
(9)
$``$ In the case of the Enskog-Grad method (EG): $`P_{yy}(y)P_{yy}=\text{constant}`$, $`Q_y(y)Q_y=\text{constant}`$, $`P_{xy}(y)=p_{xy}(y)=0`$,
$`P_{yy}`$ $`=`$ $`\left[1+{\displaystyle \frac{\delta }{\mathrm{K}\mathrm{n}}}\chi n(y)\right]p_{xx}\left(y\right)+\left[1+2{\displaystyle \frac{\delta }{\mathrm{K}\mathrm{n}}}\chi n(y)\right]n\left(y\right)T\left(y\right),`$ (10)
$`Q_y`$ $`=`$ $`\left[1+{\displaystyle \frac{3}{2}}\chi n(y){\displaystyle \frac{\delta }{\mathrm{K}\mathrm{n}}}\right]q_y^k(y)\delta ^2{\displaystyle \frac{2}{\sqrt{\pi }\mathrm{K}\mathrm{n}}}\chi n(y)^2\sqrt{T(y)}{\displaystyle \frac{dT(y)}{dy}},`$ (11)
$`{\displaystyle \frac{1}{2}}{\displaystyle \frac{dq_y^k(y)}{dy}}`$ $`=`$ $`{\displaystyle \frac{8}{\sqrt{\pi }\mathrm{K}\mathrm{n}}}\chi \sqrt{T\left(y\right)}\left[n\left(y\right)p_{xx}\left(y\right){\displaystyle \frac{q_y^k(y)^2}{128T\left(y\right)^2}}\right]`$
$`+{\displaystyle \frac{\delta }{4\mathrm{K}\mathrm{n}}}\chi \left[5q_y^k(y){\displaystyle \frac{dn(y)}{dy}}+3n(y){\displaystyle \frac{dq_y^k(y)}{dy}}\right]`$
$`{\displaystyle \frac{\delta ^2}{\sqrt{\pi }\mathrm{K}\mathrm{n}}}\chi n(y)\sqrt{T(y)}\left[2{\displaystyle \frac{dn(y)}{dy}}{\displaystyle \frac{dT(y)}{dy}}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{n(y)}{T(y)}}\left({\displaystyle \frac{dn(y)}{dy}}\right)^2+n(y){\displaystyle \frac{d^2T(y)}{dy}}\right],`$
and
$`T(y){\displaystyle \frac{dp_{xx}(y)}{dy}}p_{xx}(y){\displaystyle \frac{dT(y)}{dy}}n(y)T(y){\displaystyle \frac{dT(y)}{dy}}+{\displaystyle \frac{T(y)}{n(y)}}p_{xx}(y){\displaystyle \frac{dn(y)}{dy}}`$ (13)
$`{\displaystyle \frac{p_{xx}(y)}{n(y)}}{\displaystyle \frac{dp_{xx}(y)}{dy}}`$
$`=`$ $`{\displaystyle \frac{4}{\sqrt{\pi T(y)}\mathrm{K}\mathrm{n}}}\left[n(y)T(y)p_{xx}(y)\right]q_y^k(y)`$
$`+{\displaystyle \frac{\delta }{\mathrm{K}\mathrm{n}}}\chi [{\displaystyle \frac{7}{2}}n(y)T\left(y\right){\displaystyle \frac{dp_{xx}(y)}{dy}}+{\displaystyle \frac{7}{2}}T\left(y\right)p_{xx}(y){\displaystyle \frac{dn(y)}{dy}}+2p_{xx}\left(y\right){\displaystyle \frac{dT(y)}{dy}}`$
$`8n(y)T(y)^2{\displaystyle \frac{dn(y)}{dy}}7n\left(y\right)^2T(y){\displaystyle \frac{dT(y)}{dy}}].`$
The substitution of $`\delta =0`$ and $`\chi =1`$ in the EG equations yields the equations corresponding to the nonlinearized Boltzmann-Grad method (NLBG).
## 5 The pressure difference between the equilibrium and nonequilibrium sides
All the results we describe in what follows were obtained using perturbation methods choosing $`T_2>T_1`$ and using $`ϵ=(T_2T_1)/T_1`$ as the perturbation parameter. We solve the system of equations and their boundary (contact) conditions up to $`ϵ^6`$. We choose $`\delta =0.001`$. In such case $`\mathrm{K}\mathrm{n}`$ is in inverse proportion to $`\rho _0`$. We choose Henderson’s expression as the concrete expression for $`\chi `$:
$$\chi =\frac{1\frac{7}{16}\rho _0}{(1\rho _0)^2}.$$
(14)
We calculate the pressure in both sides of plate 2. Using $`P_{yy}`$ for the nonquilibrium steady state side and the pressure $`P_{\mathrm{𝖾𝗊}}`$ which is estimated by the state equation for the equilibrium side
$$P_{\mathrm{𝖾𝗊}}\left[1+\delta \frac{2}{\mathrm{K}\mathrm{n}}\chi n_{\mathrm{𝖾𝗊}}\right]n_{\mathrm{𝖾𝗊}},$$
(15)
the pressure difference $`\mathrm{\Delta }P`$ is defined by
$$\mathrm{\Delta }P=P_{yy}P_{\mathrm{𝖾𝗊}}.$$
(16)
Note that $`n(0)=n_{\mathrm{𝖾𝗊}}`$ and $`\mathrm{\Delta }P=0`$ to first order in $`ϵ`$. Hence, we rewrite our results in the following way
$$n(0)=n_{\mathrm{𝖾𝗊}}\left[1+\lambda _n\frac{Q_y^2}{n_{\mathrm{𝖾𝗊}}^2}\right],$$
(17)
$$\mathrm{\Delta }P=\lambda _{\mathrm{\Delta }P}\frac{Q_y^2}{n_{\mathrm{𝖾𝗊}}}$$
(18)
where $`\lambda _n`$ and $`\lambda _{\mathrm{\Delta }P}`$ are constants. Furthermore, it is possible to rewrite $`P_{yy}`$:
$$P_{yy}=n(0)\left[1+\delta \frac{2}{\mathrm{K}\mathrm{n}}\chi n(0)\right]\left[1+\lambda _p^{yy}\frac{Q_y^2}{n(0)^2}\right].$$
(19)
Tables 1 and 2 give the values of these constants for $`ϵ=0.05`$ and 0.1, respectively. The value and sign of $`\mathrm{\Delta }P`$ depend on $`ϵ`$ and $`\mathrm{K}\mathrm{n}`$. Table 3 gives the value of $`\lambda _{\mathrm{\Delta }P}`$ obtained by first order EG, that is, up to $`\delta `$ for $`ϵ=0.05`$ and 0.1, respectively. In this case, the pressure difference also exists and its value and sign depend on $`ϵ`$ and $`\mathrm{K}\mathrm{n}`$, too.
In the case of LBG, since $`p_{xx}=0`$ then $`\lambda _n=\lambda _{\mathrm{\Delta }P}=\lambda _p^{yy}=0`$. There is no pressure difference in this case.
On the other hand, for the case of NLBG, the substitution of $`\delta =0`$ in Eqs. (2), (10), (11), (4), (15), (16), (17), (18) and (19) leads to $`\lambda _n=1/256`$, $`\lambda _{\mathrm{\Delta }P}=1/256<0`$ and $`\lambda _p^{yy}=1/128`$ where for $`\lambda _n`$ and $`\lambda _{\mathrm{\Delta }P}`$ it is correct to consider only up to second order in $`Q_y`$. Hence the osmotic pressure difference does exist and its value is constant and negative.
Furthermore, for the case of EG, $`\lambda _p^{yy}/\lambda _n=2\delta \frac{4}{\mathrm{K}\mathrm{n}}n_{\mathrm{𝖾𝗊}}2`$.
We analyze the pressure difference $`\mathrm{\Delta }P`$ from another point of view. It is sufficient to calculate $`\mathrm{\Delta }P`$ up to $`ϵ^2`$. It is given by
$`\mathrm{\Delta }P`$ $`=`$ $`ϵ^2{\displaystyle \frac{\pi \mathrm{K}\mathrm{n}^2}{4096}}\left[215\rho _0^2{\displaystyle \frac{52}{\chi }}\rho _0{\displaystyle \frac{1}{\chi ^2}}\right]`$
$`=`$ $`ϵ^2{\displaystyle \frac{\pi \mathrm{K}\mathrm{n}^2}{4096(7\rho _016)^2}}[15335\rho _0^469024\rho _0^3+81344\rho _0^29216\rho _01024].`$
It is seen that the sign of $`\mathrm{\Delta }P`$ changes from negative to positive approximately at $`\rho _0=0.2`$, whereas it is always negative in the NLBG and to first order in the EG’s case.
Besides the system far from equilibrium, we are also interested in a region extremely close to the equilibrium condition. Therefore, we analyze the case without the strong nonlinear term, namely, we eliminate the terms involving $`q_y^k(y)^2`$ and $`p_{xx}(y)q_y^k(y)`$ in Eqs. (4) and (13). In this case, up to $`\delta `$,
$`P_{yy}`$ $`=`$ $`n(0)\delta {\displaystyle \frac{\sqrt{\pi }Q_y}{16n(0)}}\left({\displaystyle \frac{dn(y)}{dy}}\right)_{y=0},`$
$`\mathrm{\Delta }P`$ $`=`$ $`\delta {\displaystyle \frac{\sqrt{\pi }Q_y}{32n(0)}}\left({\displaystyle \frac{dn(y)}{dy}}\right)_{y=0}.`$ (21)
Assuming that the derivative $`dn(y)/dy`$ of the density at plate 2 has the same sign as $`Q_y`$ (this is normaly correct), $`\mathrm{\Delta }P`$ is always negative. Evaluating up to $`\delta ^2`$ and $`ϵ^2`$ yields
$`\mathrm{\Delta }P`$ $`=`$ $`ϵ^2\delta {\displaystyle \frac{\pi \mathrm{K}\mathrm{n}}{128\chi }}\left[9\chi \rho _02\right]=ϵ^2\delta {\displaystyle \frac{\pi \mathrm{K}\mathrm{n}}{128\chi }}\left[{\displaystyle \frac{9\rho _0(167\rho _0)}{16(1\rho _0)^2}}2\right],`$
$`\lambda _n`$ $`=`$ $`{\displaystyle \frac{\chi }{16}}\left[1{\displaystyle \frac{9}{2}}\chi \rho _0\right],{\displaystyle \frac{\lambda _p^{yy}}{\lambda _n}}=2+2\chi \rho _0+9(\chi \rho _0)^2.`$ (22)
As $`0<\rho _0<1`$, $`\mathrm{\Delta }P`$ is positive but $`\frac{\lambda _p^{yy}}{\lambda _n}2`$.
Furthermore, we calculate $`\lambda _n`$, $`\lambda _{\mathrm{\Delta }P}`$ and $`\frac{\lambda _p^{yy}}{\lambda _n}`$ up to $`\delta ^2`$ and $`ϵ^6`$. Tables 4 and 5 give the values of these constants for $`ϵ=0.05`$ and 0.1, respectively. The value and sign of $`\mathrm{\Delta }P`$ depends on $`ϵ`$ and $`\mathrm{K}\mathrm{n}`$, too.
## 6 Discussion
In Ref. the authors argue that there is a pressure difference at plate 2, namely the pressure in one side of the plate is different to that on the other side. They call this new pressure which acts on the central plate the “flux induced osmosis” (FIO). We consider the existence of FIO proposed by identifying $`\mathrm{\Delta }P`$ as the pressure difference defined in Sec. 5.
In Ref. the following criteria are stated:
1. $`\mathrm{\Delta }P>0`$ regardless of the sign of $`Q_y`$.
2. $`P_{yy}`$ is a function of the nonequilibrium quantities: $`T_1`$, the nonequilibrium steady heat flow $`Q_y`$, and it is related to the equilibrium quantity $`P_{\mathrm{𝖾𝗊}}`$ as long as the nonequilibrium and equilibrium temperature at both sides of plate 2 coincide,
$$\frac{n(0)}{n_{\mathrm{𝖾𝗊}}}=\left(\frac{P_{yy}}{P_{\mathrm{𝖾𝗊}}}\right)_{T_1,Q_y}$$
(23)
where $`T_1`$ is the thermodynamic temperature of plate 2.
In Ref. , for a system of hard spheres and of maxwellian particles—which obey Boltzmann’s equation and which obey the BGK equation —using the Chapman-Enskog method, it is shown that criterion 1 in Ref. is valid but criterion 2 is not valid.
On the other hand, for the hard disk’s system, from our present scheme based on Enskog’s equation we obtain that
1. Criterion 1 in Ref. is not obeyed in the case of NLBG: $`\mathrm{\Delta }P`$ is always negative independent of the sign of $`Q_y`$.
2. In the case of LBG, (23) is valid. However, substitution of (17) and (19) into (23) leads to $`\lambda _p^{yy}/\lambda _n=2`$. This is correct only in the case of NLBG. Hence, criterion 2 in Ref. is not satisfied either.
In the formulation of SST it is assumed that the number of particles and the size of the system is infinite but that the number density of the system is finite . This condition implies that the terms $`O(\delta ^2)`$ in the collision terms can be neglected. Table 3 still indicates that under such condition the osmotic pressure difference is not always positive. Especially when the system is extremely close to equilibrium, Eq. (5) implies that $`\mathrm{\Delta }P`$ is always negative.
However when the system is extremely close to equilibrium, Eq. (5) implies that $`\mathrm{\Delta }P`$ is always positive. This result coincides with the hard sphere case in Ref. .
Furthermore, the condition under which there is no heat flux at the porous wall is
$$_{\mathrm{}}^{\mathrm{}}𝑑c_x_0^{\mathrm{}}𝑑c_y\frac{c_y}{2}C^2f_{\mathrm{equil}.}+_{\mathrm{}}^{\mathrm{}}𝑑c_x_{\mathrm{}}^0𝑑c_y\frac{c_y}{2}C^2f_{y=0}=0.$$
(24)
Using Eq. (1) yields
$$q_y^k(0)=\sqrt{\frac{2}{\pi }}[n_{\mathrm{𝖾𝗊}}n(0)].$$
(25)
The above condition is not satisfied without introducing a difference between the temperature of the gas and the porous wall except in the LBG case. Therefore, except in the LBG case, it must be difficult to maintain the equilibrium state between plates 1 and 2 even if the heat conductivity of plate 2 is extremely high. Here, as a rough simplification, let us introduce the temperature of the gas in contact with plate 2, $`T_gT_1`$. Substituting $`\delta =0`$ in Eqs.(2), (10), (11), (4), (15), (16), (17), (18) and (19), using the no-mass-flux condition given by an equation similar to (1), yields,
$$\mathrm{\Delta }P=n(0)\left[1\frac{\sqrt{T_g}}{2}\frac{1}{2\sqrt{T_g}}\right]+\left[\frac{1}{\sqrt{T_g}}2\right]\frac{Q_y^2}{256n(0)}.$$
(26)
As $`T_1<T_g<T_2`$ it follows that $`1<T_g<1+ϵ`$, and putting $`q_y^k(y)^2=Q_y^2=0`$ in the right side of Eq. (4), not only for NLBG but also for LBG, $`\mathrm{\Delta }P<0`$. The behavior of the pressure difference changes qualitatively even if we restrict the analysis to Boltzmann’s regime. This implies that the estimation of $`\mathrm{\Delta }P`$ is a very delicate problem. Even if one can prepare the walls which satisfy Eq. (23) and estimate the pressure difference in Enskog’s regime, it is difficult to know what physical meaning lies behind such case.
## 7 Conclusions
We have analyzed a simple nonequilibrium steady state system inspired by . Our study refers only to a hard disk system and analyze in great detail its behavior using our extended hydrodynamic equations using various approximations. Since we obtain that the osmotic pressure difference is negative in many cases for which Eq. (23) is not satisfied we cannot agree with .
We have assumed that the pores in plate 2 are small enough and we have not considered the problem about reflections on the wall at all. As we point out in the last part of Sec. 6, the boundary (contact) condition is very delicate. We recognize that a more sophisticated analysis is necessary.
However, in cases when strong nonlinearities can be neglected and the system is quite close to equilibrium—so that higher order terms in $`ϵ`$ do not contribute—then the osmotic pressure is positive. This implies the possibility of the existence of FIO.
In addition, in the full-paper , the authors point out that Eq. (23) is directly related to the condition at the wall and this condition is essential to construct the formalism of SST in a complete form which gets a new nonequilibrium extensive quantity which determines the degree of nonequilibrium. Therefore, to clarify the problem, the measure of the pressure difference is done only for the case of a wall obeying Eq. (23). Within this context, they still recognize the results of as implying the existence FIO. Hence we also think that it is worth estimating the pressure difference starting from Eq. (23). If one were to analyze the problem in such a way then the boundary (contact) condition would have to be reconsidered to solve the kinetic equation. In other words, one would have to evaluate $`\mathrm{\Delta }P`$ under the rather complex conditions required by kinetic theory that would lead to satisfy Eq. (23). This has not been done.
The SST formalism is quite interesting and the present study has only put to test the possible existence of FIO.
Finally, we briefly comment about EIT: Extended irreversible thermodynamics . For an ideal gas, Refs. studied a problem quite similar to the one in the present article. In Ref. the authors estimated the pressure difference without considering a special wall. They assumed that the direction of the heat flow is parallel to the interface and their results are very interesting. Furthermore, in the authors estimated the difference between the pressures which are parallel and perpendicular to the heat flow.
For a hard disk, appling LBG to the systems of , it is easily possible to get similar equations to (9) and to see that the pressure difference predicted by is positive and that the difference predicted by is zero. In both cases, of course, it is totally unnecessary to use conditions (1) and (3).
Acknowledgements
The author wishes to express his sincere gratitude to Prof. P. Cordero who has guided him into the investigation using kinetic theory, given him the basic theme related to the present investigation and conducted many important discussions. The author also wishes to thank Prof. Dino Risso and Prof. Rodrigo Soto for many helpfull hints. The scholarship from Mecesup UCh 0008 project is gratefully acknowledged. |
warning/0507/quant-ph0507129.html | ar5iv | text | # Simple scheme for two-qubit Grover search in cavity QED
## Abstract
Following the proposal by F. Yamaguchi et al.\[Phys. Rev. A 66, 010302 (R) (2002)\], we present an alternative way to implement the two-qubit Grover search algorithm in cavity QED. Compared with F. Yamaguchi et al.’s proposal, with a strong resonant classical field added, our method is insensitive to both the cavity decay and thermal field, and doesn’t require that the cavity remain in the vacuum state throughout the procedure. Moreover, the qubit definitions are the same for both atoms, which makes the experiment easier. The strictly numerical simulation shows that our proposal is good enough to demonstrate a two-qubit Grover’s search with high fidelity.
The Grover search algorithm1 is an efficient quantum algorithm to look for one item in an unsorted datebase of size $`N`$. While the most efficient classical algorithm which examines items one by one needs on average $`N/2`$ queries, the Grover’s quantum algorithm uses only $`O(\sqrt{N})`$ queries to accomplish the same task. The efficiency of this algorithm has been tested experimentally in few-qubit cases by NMR 2 and by optics 3 .
Grover’s search can be briefly described as follows: For items represented by the computational states $`|X`$ with $`X=0,1\mathrm{}N1,`$ in a quantum register with $`n`$ qubits, we have $`N=2^n`$ possible states. The search starts from a superposition state $`|\mathrm{\Psi }_0=\frac{1}{\sqrt{N}}\underset{X=0}{\overset{N1}{}}|X`$, where each item has an equal probability to get picked. One search step (i.e., a query) includes two key operations 4 : (i) Inverting the amplitute of the target item; (ii) Performing a diffusion transform $`D`$, i.e., inversion about the average state $`|\mathrm{\Psi }_0`$ with $`D_{ij}=2/N`$ for $`ij`$ and $`D_{ii}=1+2/N`$ . If the target item is $`|\tau ,`$ then the operation in case (i) results in a conditional phase gate $`I_\tau =`$ $`I2|\tau \tau |,`$ where $`I`$ is the $`N\times N`$ identity matrix. After O($`\sqrt{N}`$) queries, the amplitude of the identified target would be amplified while amplitude of non-target items are shrunk to be negligible. Thus we get the target item with high probability.
Various schemes have been proposed for implementing several quantum algorithms in cavity QED, for example, Grover search algorithm5 , quantum discrete Fourier transform6 , Deutsch-Jozsa algorithm7 , quantum dense coding8 and so on. Although cavity QED is one of the qualified candidates for quantum information processing(QIP) and attracts much attention, decoherence of the cavity field remains to be a big obstacle for QIP in cavity QED. Recently, Zheng and Guo proposed two atoms interacting with a nonresonant cavity, in which the two atoms can be entangled without information transfer between atoms and the cavity9 . However, it requires the cavity to be initially in the vacuum state. Osnaghi et al.10 have experimentally demonstrated this scheme. Yamaguchi et al. extended Zheng and Guo’s proposal to a scheme to realize the two-qubit Grover search algorithm in cavity QED. In order to achieve the quantum phase gate, they have to choose different levels for qubit encoding for the two atoms5 . In Refs. 8 ; 11 , with a strong resonant classical field added, the photon-number-dependent Stark shift can be canceled. It does not require that the cavity be initially in the vacuum state, and the scheme is insensitive to both the cavity decay and the thermal field.
In this paper, we modify the proposal in Ref. 5 by adding a strong resonant classical field. Comparing with Ref. 5 , we find our method have following merits: (i) Initial vacuum cavity field is not needed, and our method is insensitive to the thermal field besides cavity decay; (ii) Qubit definitions are the same for the two atoms, which makes the experiment easier; (iii) The two-qubit gate can be easily achieved by an appropriate Rabi frequency $`\mathrm{\Omega }`$ (defined below); (iiii) Except for two NOT gates on atom 2 for labeling target state $`|g_1|e_2`$ or $`|e_1|g_2`$, all the other operations are simultaneously imposed on the two atoms, which makes the implementation more compact.
We consider two identical two-level atoms simultaneously interacting with a single-mode cavity field and driven by a classical field. The Hamiltonian (assuming $`\mathrm{}=1`$) in the rotating-wave approximation reads8 ; 11
$$H=\frac{1}{2}\underset{j=1}{\overset{2}{}}\omega _0\sigma _{z,j}+\omega _aa^+a+\underset{j=1}{\overset{2}{}}[g(a^+\sigma _j^{}+a\sigma _j^+)+\mathrm{\Omega }(\sigma _j^+e^{i\omega t}+\sigma _j^{}e^{i\omega t})]$$
(1)
where $`\sigma _{z,j}=|e_je_j||g_jg_j|`$, $`\sigma _j^+=|e_jg_j|`$, $`\sigma _j^{}=|g_je_j|`$, with $`|e_j`$ ( $`|g_j`$ ) being the excited ( ground ) state of the $`j`$th atom. $`\omega _0`$, $`\omega _a`$, $`\omega `$ are the frequency for atomic transition, cavity mode and classical field respectively. $`a^+`$, $`a`$ are the creation and annihilation operators for the cavity mode. $`g`$ is the atom-cavity coupling strength and $`\mathrm{\Omega }`$ is the Rabi frequency of the classical field. Assuming $`\omega _0=\omega `$ and $`\delta =\omega _0\omega _a,`$ we have following Hamiltonian in the interaction picture 8 ; 11
$$H_I=\underset{j=1}{\overset{2}{}}[\mathrm{\Omega }(\sigma _j^++\sigma _j^{})+g(e^{i\delta t}a^+\sigma _j^{}+e^{i\delta t}a\sigma _j^+)]$$
(2)
When $`\mathrm{\Omega }`$ $`\delta ,`$ $`g`$ and $`\delta g`$, we can get the evolution operator of the system in the interaction picture8 ; 11
$$U_I(t)=e^{iH_0t}e^{iH_et}$$
(3)
with
$$H_0=\underset{j=1}{\overset{2}{}}\mathrm{\Omega }(\sigma _j^++\sigma _j^{})$$
(4)
$$H_e=\lambda [\frac{1}{2}\underset{j=1}{\overset{2}{}}(|e_je_j|+|g_jg_j|)+(\sigma _1^+\sigma _2^++\sigma _1^+\sigma _2^{}+H.c.)]$$
(5)
where $`\lambda =g^2/2\delta .`$ It’s obvious from Eq. (5), the phonon-number dependent Stark shift has been canceled by an additional strong resonant classical field. If we define $`J_x=\frac{1}{2}\underset{j=1}{\overset{2}{}}(\sigma _j^++\sigma _j^{}),`$ Eqs. (4) and (5) reduce to $`H_0=2\mathrm{\Omega }J_x`$, $`H_e=2\lambda J_x^2`$ respectively and Eq. (3) becomes
$`U_I(t)`$ $`=e^{i2\mathrm{\Omega }tJ_x}e^{i2\lambda tJ_x^2}`$ (6)
$`=e^{ib(J_x^2+hJ_x)}`$
with $`b=2\lambda t`$ and$`h=\frac{\mathrm{\Omega }}{\lambda }`$. In the subspace spanned by $`|e_1|e_2,`$ $`|e_1|g_2,`$ $`|g_1|e_2,`$ $`|g_1|g_2,`$ we define the two-qubit Hadamard gate
$`H^2`$ $`=\underset{i=1}{\overset{2}{\mathrm{\Pi }}}H_i`$ (7)
$`=({\displaystyle \frac{1}{\sqrt{2}}})^2\left[\begin{array}{cc}1& 1\\ 1& 1\end{array}\right]\left[\begin{array}{cc}1& 1\\ 1& 1\end{array}\right]`$ (12)
$`={\displaystyle \frac{1}{2}}\left[\begin{array}{cccc}1& 1& 1& 1\\ 1& 1& 1& 1\\ 1& 1& 1& 1\\ 1& 1& 1& 1\end{array}\right]`$ (17)
where $`H_i`$ is the Hadamard gate acting on the ith atom, transforming states as $`|g_i\frac{1}{\sqrt{2}}(|g_i+|e_i)`$, $`|e_i\frac{1}{\sqrt{2}}(|g_i|e_i)`$. The $`U_I(t)`$ can be expressed in the same basis as:
$$U_I(t)=\left[\begin{array}{cccc}\frac{1}{2}+\frac{1}{2}\mathrm{cos}(bh)e^{ib}& \frac{i}{2}\mathrm{sin}(bh)e^{ib}& \frac{i}{2}\mathrm{sin}(bh)e^{ib}& \frac{1}{2}+\frac{1}{2}\mathrm{cos}(bh)e^{ib}\\ \frac{i}{2}\mathrm{sin}(bh)e^{ib}& \frac{1}{2}+\frac{1}{2}\mathrm{cos}(bh)e^{ib}& \frac{1}{2}+\frac{1}{2}\mathrm{cos}(bh)e^{ib}& \frac{i}{2}\mathrm{sin}(bh)e^{ib}\\ \frac{i}{2}\mathrm{sin}(bh)e^{ib}& \frac{1}{2}+\frac{1}{2}\mathrm{cos}(bh)e^{ib}& \frac{1}{2}+\frac{1}{2}\mathrm{cos}(bh)e^{ib}& \frac{i}{2}\mathrm{sin}(bh)e^{ib}\\ \frac{1}{2}+\frac{1}{2}\mathrm{cos}(bh)e^{ib}& \frac{i}{2}\mathrm{sin}(bh)e^{ib}& \frac{i}{2}\mathrm{sin}(bh)e^{ib}& \frac{1}{2}+\frac{1}{2}\mathrm{cos}(bh)e^{ib}\end{array}\right]$$
(18)
If we choose $`b=\frac{\pi }{2}`$, $`bh=\frac{\pi }{2}+2m\pi `$ ($`m`$ is an interger), i.e., $`\lambda t=\frac{\pi }{4},\frac{\mathrm{\Omega }}{\lambda }=4m+1,`$ we can get
$$U_I(t_D)=\left[\begin{array}{cccc}\frac{1}{2}& \frac{1}{2}& \frac{1}{2}& \frac{1}{2}\\ \frac{1}{2}& \frac{1}{2}& \frac{1}{2}& \frac{1}{2}\\ \frac{1}{2}& \frac{1}{2}& \frac{1}{2}& \frac{1}{2}\\ \frac{1}{2}& \frac{1}{2}& \frac{1}{2}& \frac{1}{2}\end{array}\right]=D$$
(19)
where $`t_D=\frac{\pi }{4\lambda }`$12 . So by choosing an appropriate value of $`\mathrm{\Omega },`$ we can generate a two-qubit diffusion transform $`D`$ (different by $`1`$ prefactor). The two-qubit conditional phase gate to label different target states will also be generated in a natrural way. It’s easy to have
$$H^2U_I(t)H^2=\left[\begin{array}{cccc}e^{ib(h1)}& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& e^{ib(h+1)}\end{array}\right]$$
(20)
If we choose $`b=\frac{\pi }{2}`$, $`h=4m+1`$ ($`m`$ is an interger), i.e., $`t_1=\frac{\pi }{4\lambda },\frac{\mathrm{\Omega }}{\lambda }=4m+1,`$ we have
$$H^2U_I(t_1)H^2=\left[\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right]=I_{g_1g_2}$$
(21)
which is to label the target state $`|g_1|g_2`$. Similarily, target state $`|e_1|e_2`$ can be labeled by setting $`b=\frac{\pi }{2}`$, $`h=4m+3`$ ($`m`$ is an interger), i.e., $`t_2=\frac{\pi }{4\lambda }`$, $`\frac{\mathrm{\Omega }}{\lambda }=4m+3`$, which yields
$$H^2U_I(t_2)H^2=\left[\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right]=I_{e_1e_2}$$
(22)
As for the target state $`|g_1|e_2`$ or $`|e_1|g_2`$ , they can be achieved by slight modification of the above operations as follows13 ,
$$\sigma _{x,2}I_{g_1g_2}\sigma _{x,2}=\left[\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right]=I_{g_1e_2}$$
(23)
$$\sigma _{x,2}I_{e_1e_2}\sigma _{x,2}=\left[\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right]=I_{e_1g_2}$$
(24)
where $`\sigma _{x,2}=\left[\begin{array}{cc}1& 0\\ 0& 1\end{array}\right]\left[\begin{array}{cc}0& 1\\ 1& 0\end{array}\right]`$ is the NOT gate acting on atom 2. Therefore, Given a specific $`\lambda ,`$ by choosing an appropriate $`\mathrm{\Omega }`$, we can generate all the two-qubit operations necessary in the two-qubit Grover’s algorithm.
To carry out our scheme, we consider a cavity, i.e., a Fabry-Perot resonator with a single mode and a standing-wave pattern along the cavity axis 5 , as shown in Fig. 1. Two atoms, first simultaneously prepared in box B into high-lying circular Rydberg state denoted by $`|g_1|g_2,`$ after $`H^2`$ operation, are in the initial average state. Then they undergo the operations in Fig. 1 from the left to the right. For searching $`|g_1|g_2`$ or $`|e_1|e_2`$, our implementation is straightforward because the atoms interact with the cavity and the classical field simultaneously. While to search $`|e_1|g_2`$ or $`|g_1|e_2`$, since NOT gates are only performed on atom 2, we have to employ an inhomogeneous field to distinguish the two atoms. This can be done by the same trick as in Ref. 5 , i.e., introducing an inhomogeneous electric field in the two regions respectively where $`\sigma _{x,2}`$ takes action. Finally the atoms are separately read out by the state-selective field-ionization detectors D<sub>1</sub> and D$`_2.`$
Now we briefly discuss the experimental possibility of our proposal. Considering two Rydberg atoms with principal quantum numbers 50 and 51 and the radiative time $`T_r=3\times 10^2`$ s, we assume that the atom-cavity coupling stength $`g`$ is $`25\times 2\pi `$ kHz 7 ; 14 , $`\delta =`$ $`20\times g,`$ and $`\mathrm{\Omega }20\times \delta .`$ Direct calculation shows that the time for the two $`U_I`$ operations is $`4.0\times 10^4`$ s. As the time for single-qubit operations is negligible, the implementation time in the cavity in Fig. 1 is much shorter than the radiative time $`T_r`$. Moreover, the atomic state evolution is independent of the cavity field state. So our proposal is realizable with presently available cavity QED techniques. Furthermore, it should be pointed out that the Rabi frequency $`\mathrm{\Omega }`$ during the two-qubit gate is about $`10\times 2\pi `$ MHz and should be slightly adjusted to satisfy the condition $`\frac{\mathrm{\Omega }}{\lambda }=4m+1`$ or $`\frac{\mathrm{\Omega }}{\lambda }=4m+3`$ mentioned above.
To check the validity of our scheme more strictly, we numerically simulate the time evolution of the system for finding the target state $`|g_1|g_2`$. As shown in Fig. 2 (a), if we assume that the cavity is initially in a Fock state $`|n,`$ the probability of finding the target state $`|g_1|g_2`$ slightly decreases with the increase of the photon number. Even for $`n=10,`$ however, the fidelity can still be $`99.2\%,`$ which means the whole process is almost independent of the cavity state. Fig. 2 (b) presents the influence of the imperfect operations on the fidelity. For simplicity, we assume the initial cavity state to be $`|5`$ and the same imperfection in each pulse. We see from the plot that even for $`7\%`$ pulse error, the fidelity is still larger than $`90\%.`$
In summary, we have proposed a simple scheme for implementing two-qubit Grover search algorithm in cavity QED. By adding a strong resonant classical field during the two-qubit operation, we cancelled the phonon-number dependent Stark shift. Thus our scheme is immune to both cavity decay and the thermal field. In addition, different from Ref. 5 , in our proposal, two levels $`|g`$ and $`|e`$ of each atoms are employed to encode the qubits (i.e., qubit definitions are the same for both atoms) and all the operations except the two NOT gates acting on atom 2 are imposed on the two atoms simultaneouly, which may make the experimental implementation easier.
We are grateful for warmhearted help from Chaohong Li, Yong Li, Qiongtao Xie and Chenxi Yue. This work is partly supported by National Natural Science Foundation of China under Grant Nos. 10474118 and 10274093, and partly by the National Fundamental Research Program of China under Grant No. 2001CB309309.
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warning/0507/math-ph0507045.html | ar5iv | text | # Geometry of quantum systems: density states and entanglement
## 1 Introduction
Dirac’s approach to Quantum Mechanics uses a Hilbert space as a fundamental object to start with, motivating the linear structure with the superposition principle necessary to describe phenomena like those of interference . Born’s probabilistic interpretation requires the use of a Hermitian inner product to deal with normalized states, therefore the physical identification of states in the Hilbert space leads to the requirement that (pure) states of a quantum mechanical system are described by elements of the complex projective space (one-dimensional subspaces of a separable complex Hilbert space $``$). By means of the Hermitian structure on $``$ it is possible to define a binary product on the pure states $`P`$ . The physical interpretation of this binary operation is given in terms of probability transition from one state to another. On this space $`P`$, bijective maps which preserve the transition probability are necessarily projection of unitary or anti-unitary transformations on the original Hilbert space, this statement is the main content of Wigner’s theorem .
More likely, due to this ”equivalence” between the two descriptions (on $``$ and on $`P`$), physicists have barely paid any attention to the geometrization of Quantum Mechanics, i.e. to introduce a ”tensorial description” in such a way that non-linear coordinate transformations could be performed, notably exception obviously do exist and we provide a partial list of references . The recent great interest in the foundational aspects of Quantum Mechanics motivated by the use of entanglement as a resource for quantum information and quantum computing has boosted a more deep study of many fundamental aspects, for instance the possibility to have a binary composition of pure states without the use of the Hilbert space linear structure , the possibility to have a non-linear Quantum Mechanics , more generally the possibility to have non-linear transformations among states.
The possibility of non-linear transformations may turn out to be quite useful in the classification problem of separability and entanglement because these properties are not preserved by taking linear combinations. Moreover, an appropriate description of atomic phenomena involving polarization, spin orientation and angular correlations, requires that we go beyond pure states in the description of quantum systems. This larger family of states was introduced by von Neumann as dual objects with respect to the quantum observables, they constitute the set of density states and an early, physically motivated, review was written by U. Fano .
Again, for these states a proper mathematical setting is provided by the dual space of the Lie algebra of the observables, with respect to the coadjoint action of the unitary group. Density states emerge as elements of the coadjoint orbits passing trough some special elements in the dual of the Cartan subalgebra. The mathematical context of coadjoint orbits is quite well known to those physicists involved with geometric quantization and the field was widely studied in the seventies by Kostant, Kirillov, and Souriau .
Each coadjoint orbit bears a natural differential structure. Observe, however, that the spectrum of the state does not change along the orbit of the unitary action. From the point of view of quantum evolution it corresponds to the situation of an isolated system, when all interactions with the environment are negligible, so there is no dissipation and the evolution is unitary. In many cases this is only a very exceptional situation, very rarely adequately corresponding to the physical reality. On the other hand, it is *a priori* not clear that the whole set of density states, i.e. a union of coadjoint orbits of the unitary action of different dimensionality, possesses a natural differential structure. Exhibiting such a structure in terms of local coordinates and/or *via* a general geometric construction of a smooth stratification of density states is thus of great interest when investigating dissipative systems.
Density states form a convex subset of the set of Hermitian operators on $``$. Some properties of these convex body attracted recently an attention . It is thus legitimate to ask about ”the shape” of the set of density matrices, in particular about the smoothness properties of its boundary. In the simplest case of the two-dimensional $``$, the density matrices form the three-dimensional unit ball with a smooth boundary - the two-dimensional unit sphere comprising all pure states. But this situation is exceptional - in higher dimensions the boundary does not consists exclusively of pure states, it is in addition not smooth.
The space of density states carries additional structures with respect to those available on the space of coadjoint orbits of general Lie groups because they are related to the unitary group and therefore additional structures are available. Moreover the need to consider composite quantum systems, tensor products of the spaces associated with a choice of subsystems making up the whole system, will bring up novel problems which will require further investigations.
All these various considerations have convinced us that a review of these mathematical aspects along with the identification of the novel emerging problems may be useful to those people interested in the application of quantum mechanics to quantum information and are not at home with the geometrical background required. A recent book by Chruściński and Jamiołkowski deals with geometrical aspects of quantum mechanics, these authors however are primarily concerned with the application of these methods to describe the geometric phase . At this point one should also point to the paper in which, in connection with geometric phase and parallel transport along mixed states, the geometry of the manifold of density matrices as a stratified space, was discussed along slightly different lines than in the present paper (cf. Section 3 below).
In the present paper the Hilbert space $``$ will be assumed to be of finite dimension $`n`$ in order to make the differential geometry expressible in local coordinates classical. The reader will understand that passing to an infinite-dimensional $``$ (i.e. a differential geometry of a Banach (or a Hilbert) manifold) is straightforward, to this aim we will try to use coordinate-free expressions, which serve in both cases, as much as possible. The paper is organized as follows :
In section 2 we start with presenting the Kähler structure on the Hilbert projective space $`P`$ of pure states obtained from the standard Hermitian product on $``$ via the momentum map associated with the Hamiltonian action of the group $`U()`$ of unitary transformations of $``$. In this picture the pure states form just an orbit in the dual space $`u^{}()`$ of the unitary Lie algebra $`u()`$ of the group $`U()`$. Because of the nondegeneracy of the canonical invariant scalar product on $`u()`$ we have a canonical identification of $`u^{}()`$ with $`u()`$ which makes the geometry of $`u^{}()`$ very rich. We decided to interpret $`u^{}()`$ as the space of Hermitian operators on $``$ which makes possible to understand the density states as a subset of $`u^{}()`$.
Consequently, in sections 3 and 4 we present the density states as a convex body $`𝒟()`$ in $`u^{}()`$ which is a family of some $`U()`$-orbits and, as we will show later, also orbits of a particular action of the group $`GL()`$ of invertible complex linear operators on $``$. We show that $`𝒟()`$ is naturally a manifold stratified space with the stratification induced by the the rank of the state. Thus the space $`𝒟^k()`$ of rank-$`k`$ states, $`k=1,\mathrm{},n`$, is a smooth manifold of (real) dimension $`2nkk^21`$ and this stratification is maximal in the sense that every smooth curve in $`𝒟()`$, viewed as a subset of the dual $`u^{}()`$ to the Lie algebra of the unitary group $`U()`$, at every point must be tangent to the strata $`𝒟^k()`$ it crosses.
Section 5 is devoted to the geometry of $`u^{}()`$, to a global description of the Kählerian structure of $`U()`$-orbits by means of the canonical Poisson and Riemann-Jordan tensors. These Kählerian structure are well-known in algebraic geometry and can be easily generalized to analogous structures on general flag manifolds. The point which should be stressed here is that the geometry we develop is canonical, that it does not depend on the matrix form of an operator and the $`U()`$-orbits are treated as a collection rather than each orbit separately.
In the last section we investigate a Hilbert space decomposition $`=^1^2`$ which is usually understood as corresponding to a quantum composite system. We present an introduction to the problems of separability and entanglement together with an abstract scheme for measurement of entanglement.
Geometry of composite quantum systems was investigated in the literature from several points of view. First, it is of importance to distinguish classes of states which are equivalent under a restricted set of unitary transformations (dubbed local transformations in the physical literature), namely those which belong to the same orbit of $`U(^1)\times U(^2)`$. From the physical point of view all states on the same orbit contain an equal amount of quantum correlations between the subsystems, i.e., these can not be influenced by operations performed separately on each subsystem.
In order to characterize uniquely an orbit (i.e. a class of locally equivalent states) one can try to find a complete set of $`U(^1)\times U(^2)`$-invariant functions on $`𝒟()`$, such that the values of all functions at $`\rho 𝒟()`$ characterize uniquely the orbit through $`\rho `$ . The task can be effectively completed only for low-dimensional systems - in fact, only in the the case dim$`^1`$=dim$`^2=2`$ the explicit results were found . The same is true for multipartite composite systems i.e. when $`=^1^2\mathrm{}^K`$. Here also the explicit results are known for $`K`$ up to $`4`$ and dim$`^i=2`$, $`i=1,\mathrm{},4`$ .
Other (partial) characterization of local orbits is provided by their dimensions. These were investigated in and in all orbits of submaximal dimensionality in the case dim$`^1`$=dim$`^2=2`$ were explicitly identified and enumerated. The similar task of finding dimension of the local orbit through an arbitrary $`\rho `$ in the case of higher dimensional systems was never achieved. A much modest goal of determining dimensions and topology of local orbits stratifying the set of rank one (pure) states $`𝒟^1(^1^2)`$ was, however, completed for arbitrary finite-dimensional $`^1`$ and $`^2`$ .
The sets of pure states in two- and three-partite systems with dim$`^i=2`$ can be identified with, respectively, unit seven- and fifteen- dimensional spheres $`𝐒^7`$ and $`𝐒^{15}`$. In both cases there exist the Hopf fibrations $`𝐒^7𝐒^4`$ and $`𝐒^{15}𝐒^8`$ which were used to investigate the geometry of pure states in , whereas multipartite pure states were treated in using Segre variety.
Although in the present paper we limit ourselves to investigation of two-partite composite system, we would like to point out recent achievements in geometric characterization of entangled pure states of multipartite systems. When investigating entanglement in multicomponent system one aims at discriminating among different classes of entanglement, defined as different equivalence classes under appropriate group of transformations preserving entanglement properties. The goal can be achieved by identification of the so called entanglement monotones, i.e. measures of entanglement which are invariant under considered transformations. Construction of such invariants based on Plücker coordinates on Grassmannians, naturally appearing when considering pure states of multicomponent systems, were presented in and . The geometry of three-qubit pure states entanglement was recently investigated in , where geometric description of different classes of entanglement was given in terms of submanifolds of the so-called Klein quadric - a special quadric embedded in the five-dimensional complex projective space.
Of special interest is the set of separable states (defined in Sec. 6), as those which, from the physical point of view, do not carry any quantum correlations. From the construction they form a convex subset in $`𝒟(^1^2)`$. Only in the case of dim$`^1`$=dim$`^2=2`$ and dim$`^1=2`$ and dim$`^2=3`$ (or vice versa) there exist effective criteria which allow to discriminate separable and nonseparable (entangled) states. As a consequence only in these low-dimensional case one can relative easily investigate the geometry of the boundary of the set of separable states .
## 2 Kähler structure on the Hilbert projective space
Let $``$ be an $`n`$-dimensional Hilbert space with the Hermitian product $`x,y_{}`$ being, by convention, $`𝐂`$-linear with respect to $`y`$ and anti-linear with respect to $`x`$. The unitary group $`U()`$ acts on $``$ preserving the Hermitian product and it consists of those complex linear operators $`Agl()`$ on $``$ which satisfy $`AA^{}=I`$, where $`A^{}`$ is the Hermitian conjugate of $`A`$, i.e.
$$Ax,y_{}=x,A^{}y_{}.$$
The geometric approach to Quantum Mechanics is based on considering the realification $`_𝐑`$ of $``$ as a Kähler manifold $`(_𝐑,J,g,𝜔)`$ with canonical structures: a complex structure $`J:\text{T}_𝐑\text{T}_𝐑`$, a Riemannian metric $`g`$, and a symplectic form $`𝜔`$. The latter come from the real and the imaginary parts of the Hermitian product, respectively, $`g=\mathrm{}(,_{})`$, $`𝜔=\mathrm{}(,_{})`$. After the obvious identification of the tangent bundle $`\text{T}_𝐑`$ with $`_𝐑\times _𝐑`$, all these structures are constant structures induced from $``$:
$$J(x)=ix,g(x,y)+i𝜔(x,y)=x,y_{}.$$
We have obvious identities
$$J^2=I,𝜔(x,Jy)=g(x,y),g(Jx,Jy)=g(x,y),𝜔(Jx,Jy)=𝜔(x,y).$$
The tensors $`g`$ and $`𝜔`$ being non-degenerate have their inverses: the contravariant metric tensor $`G=g^1`$ and the Poisson tensor $`\mathrm{\Omega }=𝜔^1`$. They form together a Hermitian product
$$\alpha ,\beta _{^{}}=G(\alpha ,\beta )+i\mathrm{\Omega }(\alpha ,\beta )$$
on the dual real Hilbert space $`_𝐑^{}`$ equipped with the dual complex structure $`J^{}`$. Using the identification of $`_𝐑^{}`$ with $`_𝐑`$ via the metric tensor $`g`$, the latter can be interpreted as a contravariant complex tensor on $`_𝐑`$. This tensor induces two real brackets of smooth functions on $`_𝐑`$: $`\{f,h\}_g=G(\text{d}f,\text{d}h)`$ and $`\{f,h\}_𝜔=\mathrm{\Omega }(\text{d}f,\text{d}h)`$. The first one is the ‘Riemann-Jordan’ bracket associated with the contravariant version of the metric tensor $`g`$ and the second is just the symplectic Poisson bracket associated with $`𝜔`$. Of course both brackets can be extended to complex functions by complex linearity and give rise to the ‘total’ bracket
$$\{f,h\}_{}=\text{d}f,\text{d}h_{^{}}=\{f,h\}_g+i\{f,h\}_𝜔.$$
(1)
Fixing an orthonormal basis $`(e_k)`$ of $``$ allows us to identify the Hermitian product $`x,y_{}`$ on $``$ with the canonical Hermitian product on $`𝐂^n`$
$$a,b_{𝐂^n}=\overline{a_k}b_k$$
(2)
(we use the convention of summation on repeated indices), the group $`U()`$ of unitary transformations of $``$ with $`U(n)`$, its Lie algebra $`u()`$ with $`u(n)`$, etc. In this picture $`(a_{jk})^{}=(\overline{a_{kj}})`$ and $`(T^{}T)_{jk}=\alpha _j,\alpha _k`$, where $`\alpha _k=(t_{jk})𝐂^n`$ are columns of the matrix $`T=(t_{jk})`$. The choice of the basis induces (global) coordinates $`(q_k,p_k)`$, $`k=1,\mathrm{},n`$, on $`_𝐑`$ by
$$e_k,x_{}=(q_k+ip_k)(x)$$
in which $`_{q_k}`$ is represented by $`e_k`$ and $`_{p_k}`$ by $`ie_k`$. Hence the complex structure reads
$$J=_{p_k}\text{d}q_k_{q_k}\text{d}p_k,$$
the Riemannian tensor
$$g=(\text{d}q_k\text{d}q_k+\text{d}p_k\text{d}p_k)=\frac{1}{2}(\text{d}q_k\text{d}g_k+\text{d}p_k\text{d}p_k)$$
and the symplectic form
$$𝜔=\text{d}q_k\text{d}p_k,$$
where $`xy=xy+yx`$ is the symmetric, and $`xy=xyyx`$ is the wedge product. In complex coordinates $`z_k=q_k+ip_k`$ one can write the Hermitian product as the complex tensor
$$,_{}=\text{d}\overline{z}_k\text{d}z_k.$$
The contravariant tensor $`G+i\mathrm{\Omega }`$ has the form
$$G+i\mathrm{\Omega }=(_{q_k}_{q_k}+_{p_k}_{p_k})+i(_{q_k}_{p_k}_{p_k}_{q_k})$$
or, in complex coordinates,
$$G+i\mathrm{\Omega }=(_{q_k}i_{p_k})(_{q_k}+i_{p_k})=4_{z_k}_{\overline{z}_k}.$$
In other words,
$$\{f,h\}_g=\frac{f}{q_k}\frac{h}{q_k}+\frac{f}{p_k}\frac{h}{p_k},$$
$$\{f,h\}_𝜔=\frac{f}{q_k}\frac{h}{p_k}\frac{f}{p_k}\frac{h}{q_k},$$
and
$$\{f,h\}_{}=4\frac{f}{z_k}\frac{h}{\overline{z}_k}.$$
Every complex linear operator $`Agl()`$ on $``$ induces the quadratic function
$$f_A(x)=\frac{1}{2}x,Ax_{}.$$
The function $`f_A`$ is real if and only if $`A`$ is Hermitian, $`A=A^{}`$.
One important convention we want to introduce is that we will identify the space of Hermitian operators $`A=A^{}`$ with the dual $`u^{}()`$ of the (real) Lie algebra $`u()`$, according to the pairing between Hermitian $`Au^{}()`$ and anti-Hermitian $`Tu()`$ operators
$$A,T=\frac{i}{2}\text{Tr}(AT).$$
The multiplication by $`i`$ establishes further a vector space isomorphism $`u()TiTu^{}()`$ which identifies the adjoint and the coadjoint action of the group $`U()`$, $`\text{Ad}_U(T)=UTU^{}`$. Under this isomorphism $`u^{}()`$ becomes a Lie algebra with the Lie bracket $`[A,B]=\frac{1}{i}[A,B]_{}`$, where $`[A,B]_{}=ABBA`$ is the commutator bracket, equipped additionally with the scalar product
$$A,B_u^{}=\frac{1}{2}\text{Tr}(AB)$$
(3)
and an additional algebraic operation, the Jordan product $`[A,B]_+=AB+BA`$. The scalar product is invariant with respect to both: the Lie bracket and the Jordan product (or bracket)
$`[A,\xi ],B_{u^{}()}`$ $`=`$ $`A,[\xi ,B]_{u^{}()},`$ (4)
$`[A,\xi ]_+,B_{u^{}()}`$ $`=`$ $`A,[\xi ,B]_+_{u^{}()}.`$ (5)
and it identifies once more $`u^{}()`$ with its dual,
$$u^{}()A\widehat{A}=\frac{1}{i}Au(),$$
so vectors with covectors. Under this identification the metric (3) correspond to the invariant metric
$$\widehat{A},\widehat{B}_u=\frac{1}{2}\text{Tr}(AB)$$
(6)
on $`u()`$ which can be viewed also as a contravariant metric on $`u^{}()`$.
For a (real) smooth function $`f`$ on $`_𝐑`$ let us denote by $`grad_f`$ and $`Ham_f`$ the gradient and the Hamiltonian vector field associated with $`f`$ and the Riemannian and the symplectic tensor, respectively. In other words, $`g(,grad_f)=\text{d}f`$ and $`𝜔(,Ham_f)=\text{d}f`$ or $`grad_f=G(\text{d}f,)`$ and $`Ham_f=\mathrm{\Omega }(\text{d}f,)`$. Note that any $`Agl()`$ induces a linear vector field $`\stackrel{~}{A}`$ on $``$ by $`\stackrel{~}{A}(x)=Ax`$.
###### Lemma 1
For Hermitian $`A`$ we have
$$grad_{f_A}=\stackrel{~}{A}\text{and}Ham_{f_A}=\stackrel{~}{iA}.$$
Proof. If $`,`$ denotes the pairing between vectors and covectors then
$`\text{d}f_A(x),y`$ $`=`$ $`{\displaystyle \frac{1}{2}}(y,Ax_{}+x,Ay_{})=\mathrm{}(y,Ax_{})`$
$`=`$ $`g(a,Ax)=𝜔(y,iAx).`$
###### Corollary 1
For all $`A,Bgl()`$ we have
$$\{f_A,f_B\}_{}=f_{2AB}.$$
(7)
In particular,
$`\{f_A,f_B\}_g`$ $`=`$ $`f_{AB+BA},`$ (8)
$`\{f_A,f_B\}_𝜔`$ $`=`$ $`f_{i(ABBA)}.`$ (9)
Proof. For Hermitian $`A,B`$ we have
$`\{f_A,f_B\}_{}(x)`$ $`=`$ $`g(grad_A(x),grad_B(x))+i𝜔(Ham_f(x),Ham_g(x))`$
$`=`$ $`g(Ax,Bx)+i𝜔(iAx,iBx)=Ax,Bx_{}=x,ABx_{}=2f_{AB}(x).`$
But $`2AB=(AB+BA)+i(i(ABBA))`$, where $`AB+BA=[A,B]_+`$ and $`i(ABBA)=i[A,B]_{}`$ are Hermitian, thus $`f_{[A,B]_+}`$ and $`f_{i[A,B]_{}}`$ are real, so the thesis holds for Hermitian $`A,B`$. For general $`A,B`$ it follows by complex linearity.
The unitary action of $`U()`$ on $``$ is in particular Hamiltonian and induces a momentum map $`\mu :_𝐑u^{}()`$. The fundamental vector field associated with $`\frac{1}{i}Au()`$, where $`Au^{}()`$ is Hermitian, reads $`\stackrel{~}{iA}`$, since
$$\frac{\text{d}}{\text{d}t}_{t=0}\mathrm{exp}(\frac{t}{i}A)(x)=iA(x).$$
The Hamiltonian of $`\stackrel{~}{iA}`$ is $`f_A`$, so the momentum map is defined by
$$\mu (x),\frac{1}{i}A=f_A(x)=\frac{1}{2}x,Ax_{}.$$
But by our convention
$$\mu (x),\frac{1}{i}A=\frac{i}{2}\text{Tr}(\mu (x)\frac{1}{i}A)=\frac{1}{2}\text{Tr}(\mu (x)A),$$
so that $`\text{Tr}(\mu (x)A)=x,Ax_{}`$ and finally, in the Dirac notation,
$$\mu (x)=xx.$$
(10)
Note that for $`A`$ being Hermitian $`f_A`$ is the pullback $`f_A=\mu ^{}(\widehat{A})=\widehat{A}\mu `$, where $`\widehat{A}=A,_u^{}=\frac{1}{i}Au()`$. The linear functions $`\widehat{A}`$ generate $`\mathrm{T}^{}u^{}()`$, so that (8) and (9) mean that the momentum map $`\mu `$ relates contravariant tensors $`G`$ and $`\mathrm{\Omega }`$ on $``$, respectively, with the linear contravariant tensors $`R`$ and $`\mathrm{\Lambda }`$ on $`u^{}()`$ corresponding to the Jordan and Lie bracket, respectively. The Riemann-Jordan tensor $`R`$, defined in the obvious way,
$$R(\xi )(\widehat{A},\widehat{B})=\xi ,[A,B]_+_u^{}=\frac{1}{2}\text{Tr}(\xi (AB+BA)),$$
(11)
is symmetric and the tensor
$$\mathrm{\Lambda }(\xi )(\widehat{A},\widehat{B})=\xi ,[A,B]_u^{}=\frac{1}{2i}\text{Tr}(\xi (ABBA)),$$
(12)
is the canonical Kostant-Kirillov-Souriau Poisson tensor on $`u^{}()`$. They form together the complex tensor
$$(R+i\mathrm{\Lambda })(\xi )(\widehat{A},\widehat{B})=2\xi ,AB_u^{}=\text{Tr}(\xi AB)$$
(13)
and the momentum map relates this tensor with the dual Hermitian product:
$$\mu _{}(G+i\mathrm{\Omega })=R+i\mathrm{\Lambda }.$$
(14)
Example. For $`=𝐂^2`$ consider an orthonormal basis in $`u^{}(2)`$ consisting of
$$U=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),X=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),Y=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),Z=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right)$$
and the associated coordinates $`u,x,y,z`$, where $`u(A)=\frac{1}{2}\text{Tr(UA)}`$, etc. In these coordinates the Poisson tensor reads
$$\mathrm{\Lambda }=2(z_x_y+x_y_z+y_z_x)$$
and the Riemann-Jordan tensor reads
$$R=_u(2x_x+2y_y+2z_z)+u(_u_u+_x_x+_y_y+_z_z).$$
The rank of $`\mathrm{\Lambda }(u,x,y,z)`$ is 0 if $`x^2+y^2+z^2=0`$ and 2 if $`x^2+y^2+z^2>0`$. The rank of $`R(u,x,y,z)`$ is 0 at $`(u,x,y,z)=0`$, it is 2 for $`u=0`$ and $`x^2+y^2+z^2>0`$, it is 3 for $`x^2+y^2+z^2=u^2>0`$, and it is 4 for $`x^2+y^2+z^2u^2>0`$.
The image $`\mu (\{0\})`$ is the cone
$$𝒫^1()=\{xx:x0\}$$
of non-negatively defined Hermitian operators $`\xi =xx`$ of rank 1. The operator $`\xi `$ is proportional to the 1-dimensional projection $`\xi /\xi `$, so $`\xi ^2=\xi \xi `$, where $`\xi =x^2`$ is the operator norm of $`\xi `$. The manifold $`𝒫^1()`$ is foliated by $`U()`$-coadjoint orbits being complex projective spaces $`𝒟_r^1()=\{xx:x=r\}`$, $`r>0`$. In particular, the momentum map image of the (2n-1)-dimensional sphere $`S_{}=\{x:x^2=x,x_{}=1\}`$ is the complex projective space $`𝒟^1()=\{xx:x=1\}`$ of 1-dimensional projectors.
The coadjoint orbits $`𝒪`$ in $`u^{}()`$ possess canonical symplectic forms $`\eta ^𝒪`$ which build together the Poisson structure $`\mathrm{\Lambda }`$. These forms, as the inverses of $`\mathrm{\Lambda }_𝒪`$, are characterized by
$$\eta _\xi ^𝒪([A,\xi ],[B,\xi ])=[A,\xi ],B_{u^{}()}=\xi ,[A,B]_u^{}.$$
(15)
Indeed, the vectors $`[A,\xi ]=\frac{1}{i}[A,\xi ]_{}`$ form the tangent space of the $`U()`$-orbit through $`𝒪`$ and $`\eta ^𝒪`$ is the inverse of $`\mathrm{\Lambda }_𝒪`$. Due to invariance of the scalar product on $`u^{}()`$:
$$\mathrm{\Lambda }_\xi (\widehat{A},\widehat{B})=\xi ,[A,B]_{u^{}()}=[\xi ,A],B_{u^{}()}=[\xi ,A],\widehat{B}.$$
(16)
Hence $`\mathrm{\#}\mathrm{\Lambda }_\xi (\widehat{A})=[\xi ,A]`$ and
$$\mathrm{\#}\eta _\xi ^r([\xi ,A])=(\mathrm{\#}\mathrm{\Lambda }_\xi )^1([\xi ,A])=\widehat{A},$$
so
$`\eta _\xi ^r([A,\xi ],[B,\xi ])`$ $`=`$ $`\mathrm{\#}\eta _\xi ^r([A,\xi ]),[B,\xi ]=\widehat{A},[B,\xi ]`$
$`=`$ $`A,[B,\xi ]_{u^{}()}=\xi ,[A,B]_{u^{}()}.`$
The image $`\mathrm{\#}R(\text{T}^{}u^{}())`$ of the tensor $`R`$ is not an involutive (generalized)distribution, so its inverse $`\sigma =R^1`$ can be understood only as a ‘partial’ covariant tensor on $`u^{}()`$, i.e. as a ‘partial symmetric 2-form’ which at $`\xi u^{}()`$ is defined only on vectors from $`\mathrm{\#}R_\xi (\text{T}_\xi ^{}u^{}())`$. There is a completely analogous characterization of the tensor $`\sigma `$ to that of $`\eta `$. Both characterizations we can summarize as follows.
###### Proposition 1
(a) The symplectic form $`\eta ^𝒪`$ on the $`U()`$-orbit $`𝒪`$ is characterized by
$$\eta _\xi ^𝒪([A,\xi ],[B,\xi ])=[A,\xi ],B_{u^{}()}=\xi ,[A,B]_u^{},$$
(17)
where $`A,Bu^{}()`$ are arbitrary Hermitian operators.
(b) The ‘partial tensor’ $`\sigma `$ on $`u^{}()`$ is characterized by
$$\sigma _\xi ([A,\xi ]_+,[B,\xi ]_+)=[A,\xi ]_+,B_{u^{}()}=\xi ,[A,B]_+_u^{}.$$
(18)
where $`A,Bu^{}()`$ are arbitrary Hermitian operators.
Let us observe that the ‘partial tensor’ $`\sigma `$, when restricted to any $`𝒟_r^1()`$, induces a Riemannian structure $`\sigma ^r`$ which, together with the symplectic structure $`\eta ^r=\eta ^{𝒟_r^1()}`$, induces a Kähler structure.
###### Proposition 2
(a) The tensor $`\sigma ^r`$ being the restriction of the partial tensor $`\sigma `$ to the $`U()`$-orbit $`𝒟_r^1()`$ through $`\xi =\mu (x)`$, $`r^2=\xi `$, is proportional to the original scalar product on $`u^{}()`$:
$$\sigma _\xi ^r([A,\xi ],[B,\xi ])=\frac{1}{\xi }[A,\xi ],[B,\xi ]_{u^{}()}.$$
(19)
(b) The $`(1,1)`$-tensor $`𝒥`$ on $`𝒫^1()`$, $`𝒥_\xi (A)=\frac{1}{\xi }[A,\xi ]`$, satisfies $`𝒥^3=𝒥`$ and induces a complex structure $`{}_{}{}^{r}𝒥`$ on every $`𝒟_r^1()`$. Moreover,
$$\eta _\xi ^r([A,\xi ],^r𝒥_\xi ([B,\xi ]))=\sigma _\xi ^r([A,\xi ],[B,\xi ]),$$
(20)
and
$$\eta _\xi ^r(^r𝒥_\xi ([A,\xi ]),^r𝒥_\xi ([B,\xi ]))=\eta _\xi ^r([A,\xi ],[B,\xi ]),$$
(21)
i.e. $`(𝒟_r^1(),^r𝒥,\sigma ^r,\eta ^r)`$ is a Kähler manifold for each $`r>0`$.
Proof. Observe first that, due to the Leibniz rule,
$$[A,\xi ]=\frac{1}{\xi }[A,\xi ^2]=\frac{1}{\xi }[[A,\xi ],\xi ]_+.$$
Then, in view of (18),
$$\sigma _\xi ^r([A,\xi ],[B,\xi ])=\frac{1}{\xi ^2}\xi ,[[A,\xi ],[B,\xi ]]_+_{u^{}()}=\frac{1}{2\xi ^2}\text{Tr}(\xi [[A,\xi ],[B,\xi ]]_+).$$
But
$`\text{Tr}(\xi [[A,\xi ],[B,\xi ]]_+)`$ $`=`$ $`\text{Tr}(\xi [A,\xi ][B,\xi ]+\xi [B,\xi ][A,\xi ])`$
$`=`$ $`\text{Tr}([A,\xi ^2][B,\xi ][A,\xi ]\xi [B,\xi ]+\xi [B,\xi ][A,\xi ])`$
$`=`$ $`\text{Tr}([A,\xi ^2][B,\xi ])=\xi \text{Tr}([A,\xi ][B,\xi ])`$
$`=`$ $`2\xi [A,\xi ],[B,\xi ]_{u^{}()},`$
that proves (19).
To prove that $`𝒥`$ is a complex structure on every orbit, let us recall that $`\xi ^2=\xi \xi `$. Passing to $`\xi ^{}=\xi /\xi `$ if necessary, we can assume for all the further calculations that $`\xi =1`$ so that $`𝒥_\xi (A)=[A,\xi ]`$. Hence,
$$[[[A,\xi ],\xi ],\xi ]=\frac{1}{i}[(A\xi ^22\xi A\xi +\xi ^2A),\xi ]_{}=\frac{1}{i}[(A\xi ^3\xi ^3A]=[A,\xi ]$$
(22)
and (19) follows. Moreover, since vectors $`[A,\xi ]`$ form the tangent space $`\text{T}_\xi 𝒟_r^1()`$, (22) shows that $`𝒥`$ reduced to $`𝒟_r^1()`$ is an almost-complex structure $`{}_{}{}^{r}𝒥`$. We shall show that the Nijenhuis torsion of $`{}_{}{}^{r}𝒥`$ vanishes, so the structure is integrable. To do this, we must show that the distribution in the complexified tangent bundle $`\text{T}𝒟_r^1()𝐂`$ which corresponds to eigenvectors of complexified $`{}_{}{}^{r}𝒥`$ with the eigenvalue $`i`$ is involutive. But this distribution is generated by complex vector fields $`\overline{T}`$ for $`Tgl()`$, where $`\overline{T}(\xi )=\xi T(1\xi )`$. Indeed,
$$𝒥_\xi (\xi T(1\xi ))=[\xi T(1\xi ),\xi ]=\frac{1}{i}(\xi T(1\xi )\xi \xi ^2T(1\xi ))=i\xi T(1\xi )$$
and this is a generating set due to the decomposition
$$T=(\xi T\xi +(1\xi )T(1\xi ))+(1\xi )T\xi +\xi T(1\xi )$$
into eigenvectors of $`𝒥`$ with eigenvalues 0, $`i`$, and $`i`$, respectively. The bracket of vector fields $`[\overline{T}_1,\overline{T}_2]_{vf}`$ reads
$`[\overline{T}_1,\overline{T}_2]_{vf}(\xi )`$ $`=`$ $`\xi T_1(1\xi )T_2(1\xi )\xi T_2\xi T_1(1\xi )\xi T_2(1\xi )T_1(1\xi )`$
$`+`$ $`\xi T_1\xi T_1(1\xi )=\xi (T_1T_2T_2T_1)(1\xi )=\overline{([T_1,T_2]_{})}`$
that proves involutivity.
Finally, it is sufficient to combine (17) and (19) to get (20). Then
$`\eta _\xi ^r(^r𝒥_\xi ([A,\xi ]),^r𝒥_\xi ([B,\xi ]))`$ $`=`$ $`\sigma _\xi ^r(^r𝒥_\xi ([A,\xi ]),[B,\xi ])=^r𝒥_\xi ([A,\xi ]),[B,\xi ]_{u^{}()}`$
$`=`$ $`[[A,\xi ],\xi ],[B,\xi ]_{u^{}()}=[[[A,\xi ],\xi ],\xi ],B_{u^{}()}`$
$`=`$ $`[[A,\xi ],B_{u^{}()}=\eta _\xi ^r([A,\xi ],[B,\xi ]),`$
that proves (21).
###### Proposition 3
There is an identification of the orthogonal complement of the vector $`x`$ with the tangent space to the $`U()`$-orbit through $`\xi =\mu (x)`$ in $`u^{}()`$. For $`y,y^{}`$ orthogonal to $`x`$ with respect to the Hermitian product, the vectors $`(\mu _{})_x(y),(\mu _{})_x(y^{})`$ are tangent to the orbit through $`\xi `$ and
$`\sigma _\xi ^r((\mu _{})_x(y),(\mu _{})_x(y^{}))`$ $`=`$ $`g(y,y^{}),`$ (23)
$`\eta _\xi ^r(,(\mu _{})_x(y^{}))`$ $`=`$ $`𝜔(y,y^{}),`$ (24)
$`{}_{}{}^{r}𝒥_{\mu }^{}(x)((\mu _{})_x(y))`$ $`=`$ $`(\mu _{})_x(Jy).`$ (25)
Proof. Since
$$(\mu _{})_x(y)=P_y^x=yx+xy,$$
can be written as $`P_y^x=[A_y,\xi ]`$, where $`A_y`$ is a Hermitian operator such that $`Ax=iy`$ and $`Ay=i\frac{y^2}{x^2}x`$, the operators $`P_y^x,P_y^{}^x`$, viewed as vectors in $`u^{}()`$, are tangent to the orbit through $`\xi `$. Then, due to (27),
$`\sigma _\xi ^r(P_y^x,P_y^{}^x)`$ $`=`$ $`{\displaystyle \frac{1}{2x^2}}\text{Tr}\left(P_y^xP_y^{}^x\right)={\displaystyle \frac{1}{2x^2}}\text{Tr}\left(x^2yy^{}+y,y^{}_{}xx\right)`$
$`=`$ $`{\displaystyle \frac{1}{2x^2}}\left(x^2\left(y^{},y_{}+y,y^{}_{}\right)\right)=\mathrm{}(y,y^{}_{})=g(y,y^{}).`$
To prove (24), we use (17):
$`\eta _\xi ^r(P_y^x,P_y^{}^x)`$ $`=`$ $`\xi ,[A_y,A_y^{}]_{u^()}={\displaystyle \frac{1}{2}}\text{Tr}\left(\xi [A_y,A_y^{}]\right)`$
$`=`$ $`{\displaystyle \frac{1}{2}}x,[A_y,A_y^{}]x_{}={\displaystyle \frac{1}{2i}}\left(A_yx,A_y^{}x_{}A_y^{}x,A_yx_{}\right)`$
$`=`$ $`\mathrm{}(iy,iy^{}_{})=𝜔(y,y^{}).`$
Finally, (25) follows directly from $`𝜔(y^{},Jy)=g(y^{},y)`$ and (20).
The above theorem says that the Kähler manifold $`(𝒟_r^1(),^r𝒥,\sigma ^r,\eta ^r)`$ comes from a sort of a ‘Kähler reduction’ of the original linear Kähler manifold $`(_𝐑,J,g,𝜔)`$. In particular, the symplectic manifold $`𝒟_r^1()`$ is the symplectic reduction of $`(_𝐑,𝜔)`$ with respect to the isotropic submanifold $`S_r=\{x:x=r\}`$. The characteristic foliation of $`𝜔_{_{S_r}}`$ consists of orbits of the group $`S^1=\{z𝐂:|z|=1\}`$ acting on $``$ by multiplication. The fundamental vector field of this action is
$$\stackrel{~}{iI}=p_k_{q_k}q_k_{p_k}$$
which is simultaneously a Killing vector field for the Riemannian metric $`g`$. Therefore $`g`$ induces a Riemannian metric on $`𝒟_r^1()`$, etc.
## 3 Smooth manifold structure on $`𝒫^k()`$
Recall that the space of non-negatively defined operators from $`gl()`$, i.e. of those $`\rho gl()`$ which can be written in the form $`\rho =T^{}T`$ for a certain $`Tgl()`$, we denote by $`𝒫()`$. It is a cone as being invariant with respect to the homoteties by $`\lambda `$ with $`\lambda 0`$. The set of density states $`𝒟()`$ is distinguished in the cone $`𝒫()`$ by the equation $`\text{Tr}(\rho )=1`$, so we will regard $`𝒫()`$ and $`𝒟()`$ as embedded in $`u^{}()`$.
The space $`𝒟()`$ is a convex set in the affine hyperplane in $`u^{}()`$, determined by the equation $`\text{Tr}(\tau )=1`$. The tangent spaces to this affine hyperplane are therefore canonically identified with the space of Hermitian operators with trace $`0`$. It is known that the set of extreme points of $`𝒟()`$ coincides with the set $`𝒟^1()`$ of pure states, i.e. the set of one-dimensional orthogonal projectors $`xx`$ (see Corollary 3). Hence every element of $`𝒟()`$ is a convex combination of points from $`𝒟^1()`$. The space $`𝒟^1()`$ of pure states can be identified with the complex projective space $`P𝐂P^{n1}`$ via the projection $`\{0\}xxx𝒟^1()`$ which identifies the points of the orbits of the $`𝐂\{0\}`$-group action by complex homoteties. We have already seen that $`𝒟^1()`$ is canonically a Kähler manifold. This will be the starting point for the study of geometry of the set $`𝒟()`$ of all density states.
The (co)adjoint action of the group $`U()`$ in $`u^{}()`$ induces its action on the positive cone $`𝒫()`$ and on the space of density states. This action is transitive on pure states but it is no longer transitive on subsets $`𝒟^k()`$, $`k>1`$, where $`𝒟^k()=𝒟()𝒫^k()`$ and $`𝒫^k()`$ consists of non-negative operators of rank $`k`$. The rank is understood clearly as the rank of the corresponding operator (or matrix, if a basis in $``$ is chosen). The intersection of $`𝒟()`$ with any Weyl chamber in a Cartan subalgebra in $`u^{}()`$ is an $`(n1)`$-dimensional simplex, while the intersection of $`𝒟^k()`$ is the $`(k1)`$-skeleton of this simplex. However, the dimension of the orbit may vary even for points from a chosen $`𝒟^k()`$ if $`k>1`$. Thus, the set of density states is a union of smooth manifolds – orbits of $`U()`$ – but the differentiable structure of the stratum $`𝒟^k()`$ is a priori not clear (for $`k>1`$), since the decomposition into orbits is not a regular foliation, i.e. $`𝒟^k()`$ is the union of a family of various submanifolds of $`u^{}()`$ which differ even by dimensions. By the differential structure we mean here the differential structure inherited from $`u^{}()`$, so that the smooth curves in $`𝒟()`$ and hence the tangent spaces are uniquely defined.
Our aim in this section is to understand this differential structure. Of course, the interior of $`𝒟()`$, namely $`𝒟^n()`$, is an open subset, so a submanifold, in the affine subspace of trace=1 Hermitian operators and the real question is only the boundary, consisting of those density states $`\rho `$ for which $`\text{det}(\rho )=0`$. The best situation would be if the boundary were a submanifold, but this is not true in dimensions $`n>2`$ as we will show later. The stratification into $`U()`$-orbits is too small, since, as it will appear later, the subsets $`𝒟^k()`$ are coarser submanifolds in $`u^{}()`$. We will show also that the stratification by rank is the maximal one in the sense that the vectors tangent to $`𝒟()`$ at $`\rho 𝒟^k()`$ must be tangent to $`𝒟^k()`$ itself, so the largest $`u^{}()`$-submanifold through $`\rho 𝒟^k()`$ contained in $`𝒟()`$ is $`𝒟^k()`$.
We start with fixing an orthonormal basis in $``$ which allows us to identify $`u^{}()`$ with the space $`u^{}(n)`$ of Hermitian $`n\times n`$-matrices which is canonically an $`n^2`$-dimensional real manifold with respect to the identification
$$u^{}(n)(a_{ij})((a_{ii})_1^n,(a_{ij})_{i<j})𝐑^n\times 𝐂^{n(n1)/2}.$$
By $`𝒫(n)`$ we denote the space of non-negatively defined matrices from $`u^{}(n)`$, by $`𝒫^k(n)`$ the subset of rank $`k`$ matrices from $`𝒫(n)`$, etc. Let us denote by $`𝒫_J^k(n)`$ the set of matrices $`A=(a_{ij})_{i,j=1}^n𝒫(n)`$ being of rank $`k`$ and such that the minor $`\text{det}[(a_{rs})_{r,sJ}]`$ associated with a set of indices $`J=\{i_1,\mathrm{},i_k\}\{1,\mathrm{},n\}`$ is non-vanishing.<sup>4</sup><sup>4</sup>4The set $`J`$ is not to be confused in the following with the complex structure denoted accidentally by the same letter. From the context, however, the notion of $`J`$ is always obvious. The next lemma shows that any matrix from $`𝒫_J^k(n)`$ can be reconstructed from its rows (or columns, since it is Hermitian) indexed by $`J`$.
###### Lemma 2
Let $`A=(a_{ij})_{i,j=1}^n𝒫_J^k(n)`$, so that the matrix $`(a_{rs})_{r,sJ}`$ has the inverse $`(a^{rs})_{r,sJ}`$. Then the matrix $`A`$ is uniquely determined by $`\{(a_{ij}):iJ,j=1,\mathrm{},n\}`$ according to the formula
$$a_{ij}=\underset{r,sJ}{}a_{ir}a^{rs}\overline{a_{js}}.$$
(26)
Proof.- The matrix $`A`$ being non-negatively defined is of the form $`T^{}T`$ for certain $`n\times n`$-matrix $`T`$, so that $`a_{ij}`$ is the Hermitian product $`\alpha _i,\alpha _j`$ of columns of $`T`$ with respect to the standard Hermitian product (2). The matrix $`T`$ is not uniquely determined. However, the fact that $`A`$ is of rank $`k`$ with the non-vanishing minor associated with $`J`$ means that the columns $`\alpha _j,jJ`$ are linearly independent and span the rest of the columns of $`T`$. But the Hermitian product on the subspace in $`𝐂^n`$ spanned by $`\{\alpha _j:jJ\}`$ is given by the formula
$$x,y_{𝐂^n}=\underset{r,sJ}{}x,\alpha _r_{𝐂^n}\alpha ^{rs}\alpha _s,y_{𝐂^n},$$
(27)
where $`(\alpha ^{rs})_{r,sJ}`$ is the inverse of the matrix $`(\alpha _r,\alpha _s_{𝐂^n})_{r,sJ}`$. The proof of (27) is immediate, since the r.h.s. of (27) is $`𝐂`$-linear with respect to $`y`$, anti-linear with respect to $`x`$ and equals $`\alpha _i,\alpha _j_{𝐂^n}`$ for $`x=\alpha _i`$, $`y=\alpha _j`$, $`i,jJ`$, by definition. Since $`a_{ij}=\alpha _i,\alpha _j_{𝐂^n}`$, we get the formula (26) directly from (27).
Remark. It is worth noticing that the formula (27) is similar to the one describing the Dirac bracket on constraint manifolds induced by second class constraints.
For $`J`$ as above define a linear map
$$\mathrm{\Phi }_J:u^{}(n)u^{}(k)\times 𝐂^{(nk)k}𝐑^k\times 𝐂^{(2nkk^2k)/2}𝐑^{2nkk^2}$$
by
$$\mathrm{\Phi }_J((a_{ij})_{i,j=1}^n)=((a_{ij})_{i,jJ},(a_{rs})_{rJ,sJ})).$$
(28)
In particular, if we work with the principal minor, i.e. $`J=\{1,\mathrm{},k\}`$, then $`\mathrm{\Phi }_J`$ associates with a Hermitian matrix its first $`k`$ columns with removed, say, upper-triangular part which is irrelevant due to hermicity or, equivalently, its first $`k`$ rows with removed lower-triangular part.
For $`Au^{}(n)`$ by $`\mathrm{\Phi }_{J,A}`$ we denote the map $`\mathrm{\Phi }_{J,A}(X)=\mathrm{\Phi }_J(X)\mathrm{\Phi }_J(A)`$:
$$\mathrm{\Phi }_{J,A}((x_{ij})_{i,j=1}^n)=((x_{ij}a_{ij})_{i,jJ},(x_{rs}a_{rs})_{rJ,sJ})).$$
(29)
With some abuse of notation, its restriction to $`𝒫^k(n)`$ we will denote by the same symbol. It is clear from the above Lemma that the map $`\mathrm{\Phi }_J`$ is continuous and injective on $`𝒫_J^k(n)`$. Thus, for $`A𝒫_J^k(n)`$, also the map $`\mathrm{\Phi }_{J,A}`$ is continuous and injective on $`𝒫_J^k(n)`$.
Conversely, every point
$$((y_{ij})_{i,jJ},(y_{rs})_{rJ,sJ}))$$
of $`u^{}(k)\times 𝐂^{(nk)k}𝐑^{2nkk^2}`$, sufficiently close to $`0`$, is the value $`\mathrm{\Phi }_{J,A}(X)`$ for a certain $`X𝒫_J^k(n)`$. Indeed, adding a small Hermitian matrix to $`(a_{ij})_{i,jJ}`$ will not change its invertibility. Hence we have to reconstruct $`X`$ out of $`\mathrm{\Phi }_J(X)`$, i.e out of the columns (and rows, since $`X`$ should be Hermitian) with indices belonging to $`J`$ and knowing that $`(x_{ij})_{i,jJ}`$ has an inverse, say, $`(x^{rs})_{r,sJ}`$. Here $`x_{ij}=a_{ij}+y_{ij}`$ for $`jJ`$. An obvious choice is the formula (26), i.e.
$$x_{ij}=\underset{r,sJ}{}x_{ir}x^{rs}\overline{x_{js}}.$$
The only thing to be checked is that $`X=(x_{ij})_{i,j=1}^n`$ defined in this way is non-negatively defined and of rank $`k`$. Assume, for simplicity of notation, that $`J=\{1,\mathrm{},k\}`$. First, we can find vectors $`\beta _1,\mathrm{},\beta _k𝐂^k`$ such that
$$x_{ij}=\beta _i,\beta _j_{𝐂^k}$$
(30)
for $`i,j=1,\mathrm{},k`$. This can be done up to a unitary transformation. For example, $`\beta _i`$ can be columns of the matrix $`\sqrt{(x_{ij})_{i,j=1}^k}`$. Then, we find (this time unique) vectors $`\beta _{k+1},\mathrm{},\beta _n𝐂^k`$ satisfying the conditions $`x_{ij}=\beta _i,\beta _j_{𝐂^k}`$, $`i=k+1,\mathrm{},n`$, $`j=1,\mathrm{},k`$. It is easy to see now that, due to the formula (27), we have (30) for all $`i,j=1,\mathrm{},n`$. This immediately implies that $`X`$ is non-negatively defined and of rank $`k`$. Moreover, since
$$x_{ij}=\underset{r,sJ}{}(a_{ir}+y_{ir})a_y^{rs}(\overline{y_{js}}+\overline{a_{js}}),$$
(31)
where $`(a_y^{rs})_{r,sJ}`$ is the inverse of the matrix $`(a_{rs}+y_{rs})_{r,sJ}`$, the matrix elements $`x_{ij}`$ rationally depend on $`y_{ml}`$, so that $`\mathrm{\Phi }_{J,A}^1`$ is smooth, thus also regular, as a function from a neighbourhhod of $`0`$ in $`𝐑^{2nkk^2}`$ into $`u^{}(n)`$, so $`𝒫^k(n)`$ is a submanifold in $`u^{}(n)`$. To see the image of the differential of $`\mathrm{\Phi }_{J,A}^1`$ at $`0`$, i.e the tangent space $`\mathrm{T}_A𝒫^k(n)`$, let us consider the linear (with respect to $`y`$) part $`(v_{ij})`$ of the r.h.s. of (31):
$$v_{ij}=\underset{r,sJ}{}(y_{ir}a^{rs}\overline{a_{js}}a_{ir}a^{rm}y_{ml}a^{ls}\overline{a_{js}}+a_{ir}a^{rs}\overline{y_{js}}).$$
(32)
To see this better, let us change the orthogonal basis of $`𝐂^n`$ for such that $`J=\{1,\mathrm{},k\}`$ and $`A`$ is diagonal, $`a_{ii}=\lambda _i`$, $`\lambda _i=0`$ for $`i>k`$. Then one can easily find that (32) takes the form
$$v_{ij}=\{\begin{array}{cc}0,\hfill & \text{if }i,j>k\hfill \\ y_{ij},\hfill & \text{if }jk.\hfill \end{array}$$
(33)
This means that in the image are arbitrary Hermitian matrices $`V=(v_{ij})_{i,j=1}^n`$ such that $`v_{ij}=0`$ for $`i,j>k`$, that can be written in a coordinate-free way as $`Vx,y_{𝐂^n}=0`$ for all $`x,y\text{Ker}(A)`$. Note that the manifold $`𝒫^k()`$ is connected. Indeed. it consists of connected orbits of the group $`U()`$ which meet a Weyl chamber as the $`(k1)`$-dimensional skeleton of a simplex. However, the connected components of this skeleton are identified by the action of the Weyl group, so they form topologically a $`(k1)`$-dimensional simplex which is obviously connected. Therefore we have proved the following.
###### Theorem 1
Let $`A𝒫_J^k(n)`$. Then the map $`\mathrm{\Phi }_{J,A}:𝒫^k(n)𝐑^{2nkk^2}`$ defined by (29) is a local homeomorphism from a neighbourhood of $`A`$ in $`𝒫^k(n)`$ onto a neighbourhood of $`0`$ in $`u^{}(k)\times 𝐂^{(nk)k}𝐑^{2nkk^2}`$. Moreover, the collection of the maps $`\mathrm{\Phi }_{J,A}^1:𝒲_{J,A}𝒫^k(n)u^{}(n)`$ defined on sufficiently small neighbourhoods $`𝒲_{J,A}`$ of $`0`$ by the formula (31) constitutes a smooth manifold structure on $`𝒫^k(n)`$ which makes it into a smooth and connected submanifold of $`u^{}(n)`$. The tangent space $`\mathrm{T}_A𝒫^k(n)`$, viewed as a subspace of $`u^{}(n)`$ consists of matrices $`Vu^{}(n)`$ satisfying $`Vx,y_{𝐂^n}=0`$ for all $`x,y\text{Ker}(A)`$.
Remark. In section 5 we obtain the manifold structure on $`𝒫^k(n)`$ much simpler as the structure of an $`GL(n,𝐂)`$-orbit. But we find that Lemma 2 and Theorem 1 are of some interest per se providing explicit coordinate systems.
The next theorem shows that smooth curves in $`u^{}(n)`$ which lay in $`𝒫(n)`$ cannot cross $`𝒫^k(n)`$ transversally, i.e. $`𝒫^k(n)`$ is in a sense an edge for $`𝒫^{k+1}(n)`$ if $`k<n1`$.
###### Theorem 2
Let $`\gamma :𝐑u^{}(n)`$ be a smooth curve in the space of Hermitian matrices which lies entirely in $`𝒫(n)`$. Then $`\gamma `$ is tangent to the stratum $`𝒫^k(n)`$ it belongs, i.e. $`\gamma (t)𝒫^k(n)`$ implies $`\dot{\gamma }(t)\mathrm{T}_{\gamma (t)}𝒫^k(n)`$.
Proof.- Of course, it is enough to prove the above for an arbitrary $`t𝐑`$, say, $`t=0`$. Assume therefore that $`A=\gamma (0)𝒫^k(n)`$. Take $`x\text{Ker}(A)`$. Since
$$\frac{\gamma (\mathrm{\Delta }t)\gamma (0)}{\mathrm{\Delta }t}x,x0$$
for $`\mathrm{\Delta }t0`$, we have $`\dot{\gamma }(0)x,x0`$. Taking in turn $`\mathrm{\Delta }t0`$ we see in a similar way that $`\dot{\gamma }(0)x,x0`$, so
$$\dot{\gamma }(0)x,x=0.$$
(34)
By polarization of (34) we get
$$\dot{\gamma }(0)x,y+\dot{\gamma }(0)y,x=0$$
(35)
for all $`x,y\text{Ker}(A)`$. But $`\dot{\gamma }(0)`$ is Hermitian, so
$$\dot{\gamma }(0)y,x=y,\dot{\gamma }(0)x$$
and (35) yields that the real part $`\mathrm{}(\dot{\gamma }(0)x,y)`$ is $`0`$ for all $`x,y\text{Ker}(A)`$. On the other hand, the kernel of $`A`$ is a complex subspace and
$$\mathrm{}(\dot{\gamma }(0)x,iy)=\mathrm{}(\dot{\gamma }(0)x,y)$$
so
$$\dot{\gamma }(0)x,y=0$$
(36)
for all $`x,y\text{Ker}(A)`$. But, according to Theorem 1, (36) means that $`\dot{\gamma }(0)\mathrm{T}_A𝒫^k(n)`$.
## 4 Smooth stratification of density states
The set $`𝒟()`$ of density states on $``$ is the intersection of the cone $`𝒫()`$ with the affine subspace $`\{Au^{}():\text{Tr}(A)=1\}`$ or, in other words, it is the level set of the function $`\text{Tr}:𝒫()𝐑`$ corresponding to the value $`1`$. Since $`\text{Tr}(t\rho )=t\text{Tr}(\rho )`$ and $`𝒫^k`$ is invariant with respect to homoteties with positive $`t`$, it is clear that Tr is a regular function on each $`𝒫^k()`$, so that $`𝒟^k()`$ is canonically a smooth manifold. Since topologically $`𝒫^k()𝒟^k()\times 𝐑`$, the manifolds $`𝒟^k()`$ are connected. All these observations together with Theorems 1 and 2 can be summarized in the following.
###### Theorem 3
The spaces $`𝒟^k()`$ of density states of rank $`k`$, $`k=1,\mathrm{},n`$, are smooth and connected submanifolds in $`u^{}()`$ of (real) dimension $`2nkk^21`$. The tangent space $`\mathrm{T}_\rho 𝒟^k()`$ is characterized as the space of those Hermitian operators $`T`$ of trace $`0`$ which satisfy $`Tx,y=0`$ for all $`x,y\text{Ker}(\rho )`$. Moreover, the stratification into submanifolds $`𝒟^k()`$ is maximal in the sense that every smooth curve in $`u^{}()`$, which lies entirely in $`𝒟()`$, at every point is tangent to the strata $`𝒟^k()`$ to which it actually belongs.
###### Corollary 2
The boundary $`𝒟()=_{k<n}𝒟^k()`$ of the set of density states is not a smooth submanifold of $`u^{}()`$ if $`n=\text{dim}>2`$.
Proof.- If $`n>2`$ then the boundary $`𝒟()`$ has at least two different strata and the vectors orthogonal to, say, the stratum $`𝒟^1()`$ of pure states are not tangent to $`𝒟()`$. But the dimension of $`𝒟^1()`$ is smaller than the topological dimension of $`𝒟()`$.
Remark. It is well known that for $`n=2`$ the convex set of density states is affinely equivalent to the three-dimensional ball and its boundary – to the two-dimensional sphere, so it is a smooth manifold.
The last problem concerning the geometry of density states we will consider is the question of affine parts of the manifolds $`𝒟^k()`$. It is motivated by the fact that the set $`𝒟^1()`$ of pure states is exactly the set of extremal elements of $`𝒟()`$, so it does not contain intervals, but the other strata $`𝒟^k()`$ with $`k>1`$ must do as shows the following theorem. Recall that a non-empty closed convex subset $`K_0`$ of a closed convex set $`K`$ is called a face (or extremal subset) of $`K`$ if any closed segment in $`K`$ with an interior point in $`K_0`$ lies entirely in $`K_0`$; a point $`x`$ is called an extreme point of $`K`$ if the set $`\{x\}`$ is a face of $`K`$.
###### Theorem 4
If $`\rho 𝒟^k()`$ then the affine space in $`u^{}()`$ which is tangent to $`𝒟^k()`$ at $`\rho `$ intersects $`𝒟()`$ along a $`(k^21)`$-dimensional convex body which is affinely equivalent to the set $`𝒟(k)`$ of density states in dimension $`k`$. This convex body is exactly the face of $`𝒟()`$ at $`\rho `$. In other words, the face of $`𝒟()`$ at $`\rho 𝒟^k()`$ is affinely equivalent to $`𝒟(k)`$.
Proof.- Let us take coordinates in $`u^{}()`$, i.e. let us chose an orthonormal basis in $``$, in which $`\rho `$ is represented by a diagonal matrix $`(\rho _{ij})`$, $`\rho _{ij}=\delta _j^i\lambda _i`$, where $`\lambda _i=0`$ for $`i>k`$. According to the form of $`\mathrm{T}_\rho 𝒟^k()`$, matrices $`(x_{ij})`$ which belong to $`\rho +\mathrm{T}_\rho 𝒟^k()`$ have entries $`x_{ij}`$ with $`i,j>k`$ equal to $`0`$. If they belong as well to $`𝒟()`$, also $`x_{ij}=0`$ if $`i>k`$ or $`j>k`$. Indeed, since $`x_{ij}=\alpha _i,\alpha _j`$ for certain vectors $`z_i𝐂^n`$, we have $`x_{ii}=\alpha _i^2=0`$, so $`\alpha _i=0`$, for $`i>k`$, and further $`x_{ij}=\alpha _i,\alpha _j=0`$ if $`i>k`$ or $`j>k`$. In other words, the only non-zero part of $`X`$ is the block $`(x_{ij})_{i,j=1}^k`$ which is therefore an element of $`𝒟(k)`$. Conversely, every matrix $`X`$ with such a block form belongs simultaneously to $`𝒟()`$ and, since $`(X\rho )_{ij}=0`$ for $`i,j>k`$, to $`\rho +\mathrm{T}_\rho 𝒟^k()`$. To see that $`(\rho +\mathrm{T}_\rho 𝒟^k())𝒟()`$ is exactly the face of $`𝒟()`$ at $`\rho `$, consider a segment in $`𝒟()`$ for which $`\rho `$ is an interior point. The open segment is clearly a smooth curve in $`𝒟()`$, so, in view of Theorem 3, it is tangent to $`𝒟^k()`$ at $`\rho `$, thus belongs entirely to $`\rho +\mathrm{T}_\rho 𝒟^k()`$.
###### Corollary 3
Extremal points of $`𝒟()`$ are exactly pure states.
## 5 Geometry of $`u^{}()`$
Let us mention that a major part of what has been said about the differential structure of the space $`𝒫^k()`$ of rank-$`k`$ positive operators can be repeated for the space of all rank-$`k`$ Hermitian operators. Denote by $`u_{k_+,k_{}}^{}()`$ the set of those Hermitian operators $`\xi `$ whose spectrum contains $`k_+`$ positive and $`k_{}`$ negative eigenvalues (counted with multiplicities), respectively. Thus the rank of $`\xi `$ is $`k=k_++k_{}`$ and $`𝒫^k(n)=u_{k,0}^{}(n)`$.
Fixing an orthogonal basis in $``$ will identify $`u_{k_+,k_{}}^{}()`$ with the space $`u_{k_+,k_{}}^{}(n)`$ of $`n\times n`$ Hermitian matrices of rank $`k`$ with the corresponding spectrum. Denote by $`D_k_{}^{k_+}`$ the diagonal matrix $`diag(1,\mathrm{},1,1,\mathrm{},1,0,\mathrm{},0)`$ with $`1`$ coming $`k_+`$-times and $`1`$ coming $`k_{}`$-times. Denote by $`,_{k_+,k_{}}`$ the ‘semiHermitian’ product in $`𝐂^n`$ represented by $`D_k_{}^{k_+}`$:
$$a,b_{k_+,k_{}}=\underset{j=1}{\overset{k_+}{}}\overline{a_j}b_j\underset{j=k_++1}{\overset{k_++k_{}}{}}\overline{a_j}b_j.$$
(37)
It is easy to see the following.
###### Proposition 4
Any $`\xi =(a_{ij})u_{k_+,k_{}}^{}(n)`$ can be written in the form $`\xi =T^{}D_k_{}^{k_+}T`$ for certain $`TGL(n,𝐂)`$. In other words the entries of the matrix $`\xi `$ are semiHermitian products $`a_{ij}=\alpha _i,\alpha _j_{k_+,k_{}}`$, where $`\alpha _i`$ denotes the $`i`$th column of $`T`$.
Proof. We can diagonalize $`\xi `$ by means of an unitary matrix $`U`$,
$$U\xi U^{}=diag(\lambda _1,\mathrm{},\lambda _n),$$
where $`\lambda _1\mathrm{}\lambda _n`$, so $`\lambda _1,\mathrm{},\lambda _{k_+}>0`$ and $`\lambda _{k_++1},\mathrm{},\lambda _{k_++k_{}}<0`$. Hence $`\xi =T^{}D_k_{}^{k_+}T`$ for $`T=CU`$ with
$$C=diag(\sqrt{|\lambda _1|},\mathrm{},\sqrt{|\lambda _{k_++k_{}}|},1\mathrm{},1).$$
Now, we can reformulate Lemma 2 for $`u_{k_+,k_{}}^{}(n)`$ instead of $`𝒫^k(n)`$. The proof is essentially the same with the difference that we use the semiHermitian product $`,_{k_+,k_{}}`$ in $`𝐂^n`$ instead of $`,_{𝐂^n}`$.
###### Lemma 3
Let $`\xi =(a_{ij})_{i,j=1}^nu_{k_+,k_{}}^{}(n)`$. Assume that the matrix $`(a_{rs})_{r,sJ}`$ has the inverse $`(a^{rs})_{r,sJ}`$ for certain $`k=(k_++k_{})`$-element subset $`J=\{j_1,\mathrm{},j_k\}\{1,\mathrm{},n\}`$. Then the matrix $`\xi `$ is uniquely determined by $`\{(a_{ij}):iJ,j=1,\mathrm{},n\}`$ according to the formula
$$a_{ij}=\underset{r,sJ}{}a_{ir}a^{rs}\overline{a_{js}}.$$
(38)
One can now prove that $`u_{k_+,k_{}}^{}()`$ are submanifolds of $`u^{}()`$ in completely parallel way to the case of $`𝒫^k()`$. However, Proposition 4 suggest an easier (although less constructive) way to do it. Namely, we can see $`u_{k_+,k_{}}^{}()`$ as an orbit of a natural $`GL()`$ action on $`u^{}()`$.
###### Theorem 5
The family
$$\{u_{k_+,k_{}}^{}():k_+,k_{}0,k=k_++k_{}n\}$$
(39)
of subsets of $`u^{}()`$ is exactly the family of orbits of the smooth action of the group $`GL()`$ given by
$$GL()\times u^{}()(T,\xi )T\xi T^{}u^{}().$$
(40)
In particular, every $`u_{k_+,k_{}}^{}()`$ is a connected submanifold of $`u^{}()`$ and the tangent space to $`u_{k_+,k_{}}^{}()`$ at $`\xi `$ is characterized by
$$B\text{T}_\xi u_{k_+,k_{}}^{}()x,y\text{Ker}(\xi )[Bx,y_{}=0].$$
(41)
Moreover, the following are equivalent:
(1) $`u_{k_+,k_{}}^{}()`$ intersects $`𝒫()`$;
(2) $`u_{k_+,k_{}}^{}()`$ is contained in $`𝒫()`$;
(3) $`k_{}=0`$;
(4) $`u_{k_+,k_{}}^{}()=𝒫^k()`$, $`k=k_++k_{}`$.
Proof. The proof that (40) is a group smooth action is straightforward. Proposition 4 shows that $`u_{k_+,k_{}}^{}()`$ is contained in the $`GL()`$-orbit of $`D_k_{}^{k_+}`$.
On the other hand, although the spectrum is not fixed on every $`GL()`$-orbit, the number of positive and the number of negative eigenvalues (counted with multiplicities) are fixed along the orbit. Indeed, if $`x,\xi x_{}>0`$ (resp. $`x,\xi x_{}<0`$) for $`x`$ in a $`k_+`$-dimensional (resp. $`k_{}`$-dimensional) linear subspace $`V_+`$ (resp. $`V_{}`$), then $`x,T\xi T^{}x_{}=T^{}x,\xi T^{}x_{}>0`$ (resp. $`x,T\xi T^{}x_{}=T^{}x,\xi T^{}x_{}<0`$) for $`x`$ in a $`k_+`$-dimensional (resp. $`k_{}`$-dimensional) linear subspace $`(T^{})^1(V_+)`$ (resp. $`(T^{})^1(V_{})`$).
The corresponding infinitesimal action of $`vgl()`$ is $`\xi v\xi +\xi v^{}`$ and the operators $`\xi _v=v\xi +\xi v^{}`$ clearly satisfy $`x,\xi _vy_{}=0`$ for all $`x,y\text{Ker}(\xi )`$. Conversely, if for certain $`Bu^{}()`$ we have $`Bx,y_{}=0`$ for all $`x,y\text{Ker}(\xi )`$, then $`B`$ can be written in the form $`v\xi +\xi v^{}`$. To see this, consider the splitting $`=V_1V_2`$, where $`V_2=\text{Ker}(\xi )`$ and $`V_1=V_2^{}`$. According to this splitting $`\xi `$ can be written in the operator matrix form
$$\xi =\left(\begin{array}{cc}\xi _1& 0\\ 0& 0\end{array}\right),$$
where $`\xi _1`$ is Hermitian and invertible. Similarly, $`B`$ has the form
$$B=\left(\begin{array}{cc}B_{11}& B_{12}\\ B_{21}& 0\end{array}\right),$$
where $`B_{11}^{}=B_{11}`$ and, in the obvious sense, $`B_{21}=B_{12}^{}`$. Now, it is easy to see that $`B=v\xi +\xi v^{}`$, where
$$v=\left(\begin{array}{cc}\frac{1}{2}B_{11}\xi _1^1& \xi _1B_{12}\\ B_{21}\xi _1^1& 0\end{array}\right)$$
that proves (41).
Finally, if $`u_{k_+,k_{}}^{}()`$ intersects $`𝒫()`$, then it contains an element with non-negative spectrum. But the signs of the elements of the spectrum are constant along a $`GL()`$-orbit which means that $`k_{}=0`$ and $`u_{k_+,k_{}}^{}()=𝒫^k()𝒫()`$.
Note that the fundamental vector fields $`\stackrel{~}{a}(\xi )=a\xi \xi a^{}`$ of the $`GL()`$-action satisfy the commutation rules $`[\stackrel{~}{a},\stackrel{~}{b}]_{vf}=\stackrel{~}{[a,b]_{}}`$
The next results shows that the foliation into submanifolds $`u_{k_+,k_{}}^{}()`$ can be obtained directly from tensors $`\mathrm{\Lambda }`$ and $`R`$. We know already that the (generalized) distribution $`D_\mathrm{\Lambda }`$ induced by $`\mathrm{\Lambda }`$ is generated by vector fields $`\mathrm{\Lambda }_A(\xi )=\mathrm{\#}\mathrm{\Lambda }_\xi (\widehat{A})=[A,\xi ]`$ and the (generalized distribution $`D_R`$ induced by $`R`$ is generated by vector fields $`R_A(\xi )=\mathrm{\#}R_\xi (\widehat{A})=[A,\xi ]_+`$. The following is straightforward.
###### Theorem 6
The family $`\{\mathrm{\Lambda }_A,R_A:Au^{}()\}`$ of linear vector fields on $`u^{}()`$ is the family of fundamental vector fields of the $`GL()`$-action:
$`\mathrm{\Lambda }_A(\xi )`$ $`=`$ $`{\displaystyle \frac{1}{i}}(A\xi \xi A)=(iA)\xi \xi (iA)^{}=\stackrel{~}{iA}(\xi ),`$ (42)
$`R_A(\xi )`$ $`=`$ $`A\xi +\xi A=A\xi +\xi A^{}=\stackrel{~}{A}(\xi ).`$ (43)
In particular,
$$[\mathrm{\Lambda }_A,\mathrm{\Lambda }_B]_{vf}=\mathrm{\Lambda }_{[A,B]},[R_A,R_B]_{vf}=\mathrm{\Lambda }_{[A,B]},[R_A,\mathrm{\Lambda }_B]_{vf}=R_{[A,B]},$$
(44)
so the (generalized) distribution induced by jointly by the tensors $`\mathrm{\Lambda }`$ and $`R`$ is completely integrable and $`u_{k_+,k_{}}^{}()`$ are the maximal integrate submanifolds.
###### Corollary 4
The generalized distributions $`D_{gl}=D_R+D_\mathrm{\Lambda }`$, $`D_\mathrm{\Lambda }`$ and $`D_0=D_RD_\mathrm{\Lambda }`$ on $`u^{}()`$ are involutive and can be integrated to generalized foliations $`_{gl}`$, $`_\mathrm{\Lambda }`$, and $`_0`$, respectively. The leaves of the foliation $`_{gl}`$ are the orbits of the $`GL()`$ action $`\xi T\xi T^{}`$, the leaves of $`_\mathrm{\Lambda }`$ are the orbits of the $`U()`$-action.
Denote by $`\stackrel{~}{𝒥}`$ and $`\stackrel{~}{}`$ the $`(1,1)`$-tensors on $`u^{}()`$, viewed as a vector bundle morphism induced by the contravariant tensors $`\mathrm{\Lambda }`$ and $`R`$, respectively,
$`\stackrel{~}{𝒥},\stackrel{~}{R}`$ $`:`$ $`\text{T}u^{}()\text{T}u^{}(),`$
$`\stackrel{~}{𝒥}_\xi (A)`$ $`=`$ $`[A,\xi ]=\mathrm{\Lambda }_\xi (A),`$
$`\stackrel{~}{}_\xi (A)`$ $`=`$ $`[A,\xi ]_+=R_\xi (A),`$
where $`Au^{}()\text{T}_\xi u^{}()`$. The image of $`\stackrel{~}{𝒥}`$ is $`D_\mathrm{\Lambda }`$ and the image of $`\stackrel{~}{}`$ is $`D_R`$.
###### Lemma 4
The tensors $`\stackrel{~}{𝒥}`$ and $`\stackrel{~}{}`$ commute and
$$\stackrel{~}{𝒥}_\xi \stackrel{~}{}_\xi (A)=\stackrel{~}{}_\xi \stackrel{~}{𝒥}_\xi (A)=[A,\xi ^2].$$
(45)
Proof. We have
$$\stackrel{~}{𝒥}_\xi _\xi (A)=[[A,\xi ]_+,\xi ].$$
But, as easily seen,
$$[[A,\xi ]_+,\xi ]=[A,\xi ^2]=[[A,\xi ],\xi ]_+=\stackrel{~}{}_\xi \stackrel{~}{𝒥}_\xi (A).$$
(46)
Recall that $`U()`$-orbits $`𝒪`$, i.e. the orbits with respect to the action of the subgroup $`U()GL()`$, carry canonical symplectic structures $`\eta ^𝒪`$. The symplectic structures $`\eta ^𝒪`$ is $`U()`$-invariant, i.e. $`(𝒪,\eta ^𝒪)`$ is a homogeneous symplectic manifold. We will show that this symplectic structure is a part of a canonical Kähler structure. We know already this structure for the orbits $`𝒫_r^1()`$.
Recall also that on $`u^{}()`$ we have the Riemannian metric induced by the scalar product $`A,B_u^{}=\frac{1}{2}\text{Tr}(AB)`$ on $`u^{}()`$.
###### Theorem 7
(a) The image of $`\stackrel{~}{𝒥}_\xi `$ is $`\text{T}_\xi 𝒪`$ and $`\text{Ker}(\stackrel{~}{𝒥}_\xi )`$ is the orthogonal complement of $`\text{T}_\xi 𝒪`$.
(b) $`\stackrel{~}{𝒥}_\xi ^2`$ is a self-adjoint (with respect to $`,_u^{}`$) and negatively defined operator on $`\text{T}_\xi 𝒪`$.
(c) The $`(1,1)`$-tensor $`𝒥`$ on $`u^{}()`$ defined by
$$𝒥_\xi (A)=\left((\stackrel{~}{𝒥}_\xi )_{\text{T}_\xi 𝒪}^2\right)^{\frac{1}{2}}\stackrel{~}{𝒥}_\xi (A)$$
(47)
induces an $`U()`$-invariant complex structure $`𝒥`$ on every orbit $`𝒪`$.
(d) The tensor
$$\gamma _\xi ^𝒪(A,B)=\eta _\xi ^𝒪(A,𝒥_\xi (B))$$
(48)
is an $`U()`$-invariant Riemannian metric on $`𝒪`$ and
$$\gamma _\xi ^𝒪(𝒥_\xi (A),B)=\eta _\xi ^𝒪(A,B).$$
(49)
In particular, $`(𝒪,𝒥,\eta ^𝒪,\gamma ^𝒪)`$ is a homogeneous Kähler manifold. Moreover, if $`\xi u^{}()`$ is a projector and $`\xi 𝒪`$, then $`𝒥_\xi =\stackrel{~}{𝒥}_\xi `$ and $`\gamma ^𝒪(A,B)=A,B_u^{}`$.
Remark. The tensor $`𝒥`$ is canonically and globally defined. It is however not smooth as a tensor field on $`u^{}()`$. It is smooth on the open-dense subset of regular elements and, of course, on every $`U()`$-orbit separately.
Proof.
(a) The vector fields $`\mathrm{\Lambda }_A(\xi )=[A,\xi ]=\stackrel{~}{𝒥}_\xi (A)`$ are fundamental vector fields of the $`U()`$-action, so $`\text{T}_\xi 𝒪`$ is the image of $`\stackrel{~}{𝒥}_\xi `$. Moreover, the invariance of the Riemannian metric $`A,B_u^{}`$,
$$\stackrel{~}{𝒥}_\xi (A),B_u^{}=[A,\xi ],B_u^{}=A,\stackrel{~}{𝒥}_\xi (B)_u^{},$$
(50)
implies that
$$B\text{Ker}(\stackrel{~}{𝒥}_\xi )B\stackrel{~}{𝒥}_\xi (u^{}()).$$
(b) The identity (50) means that $`\stackrel{~}{𝒥}_\xi ^\times =\stackrel{~}{𝒥}_\xi `$, where $`\stackrel{~}{𝒥}_\xi ^\times `$ is the adjoint operator to $`\stackrel{~}{𝒥}_\xi `$ with respect to the scalar product $`A,B_u^{}`$. Consequently,
$$(\stackrel{~}{𝒥}_\xi ^2)^\times =\stackrel{~}{𝒥}_\xi ^2.$$
(51)
Moreover, $`\stackrel{~}{𝒥}_\xi ^2`$ is negatively defined on $`\text{T}_\xi 𝒪`$, since
$$\stackrel{~}{𝒥}_\xi ^2(A),A_u^{}=[[A,\xi ],\xi ],A_u^{}=[A,\xi ],[A,\xi ]_u^{}<0,$$
for $`[A,\xi ]\text{T}_\xi 𝒪`$, $`[A,\xi ]0`$.
(c) The tensor $`\stackrel{~}{𝒥}`$ is clearly $`U()`$-invariant:
$$\stackrel{~}{𝒥}_{U\xi U^{}}(UAU^{})=[UAU^{},U\xi U^{}]=U[A,\xi ]U^{}=U(\stackrel{~}{𝒥}_\xi (A))U^{},$$
(52)
so $`U()`$-invariant is the tensor $`\left(\stackrel{~}{𝒥}^2\right)^{\frac{1}{2}}`$ and its composition $`𝒥`$. The tensor $`𝒥`$ defines an almost complex structure on every orbit $`𝒪`$, since
$$\left[\left(\stackrel{~}{𝒥}^2\right)^{\frac{1}{2}}\stackrel{~}{𝒥}\right]^2=\left(\stackrel{~}{𝒥}^2\right)^1\stackrel{~}{𝒥}^2=I.$$
To show that this almost complex structure is integrable, it is sufficient to show that the distribution $`𝒩`$ in the complexified tangent bundle $`\text{T}𝒪𝐂`$ which consists of $`i`$-eigenvectors of (complexified) $`𝒥`$ is involutive. Since $`𝒥`$, and therefore $`𝒩`$, is invariant, it is sufficient to check it at one point, say $`\xi 𝒪`$ with respect to the complexified Lie algebra $`gl()=u^{}()𝐂`$ equipped with the bracket $`[a,b]=\frac{1}{i}[abba]`$.
Let $`\kappa _1^2,\mathrm{},\kappa _m^2`$, where $`\kappa _1,\mathrm{},\kappa _m>0`$, be the eigenvalues of $`(\stackrel{~}{𝒥}_\xi ^2)_{\text{T}_\xi 𝒪}`$ counted with multiplicities. The complexified $`\stackrel{~}{𝒥}_\xi `$, which with some abuse of notation we will denote by the same symbol, has therefore eigenvalues $`\pm i\kappa _k`$ with eigenvectors $`a_k^\pm `$, $`k=1,\mathrm{},m`$ and $`𝒥_\xi (a_k^\pm )=\pm ia_k^\pm `$. Thus $`𝒩_\xi `$ is spanned by the vectors $`a_k^+,k=1,\mathrm{},m`$, i.e. eigenvectors of $`\stackrel{~}{𝒥}_\xi `$, $`\stackrel{~}{𝒥}_\xi (a_k^+)=i\kappa _ka_k^+`$ with positive $`\kappa _k`$. This space is clearly a Lie subalgebra in $`gl()`$, since
$`𝒥_\xi ([a_k^+,a_l^+])`$ $`=`$ $`[[a_k^+,a_l^+],\xi ]=[[a_k^+,\xi ],a_l^+]+[a_k^+,[a_l^+,\xi ]]`$
$`=`$ $`[i\kappa _ka_k^+,a_l^+]+[a_k^+,i\kappa _la_l^+]=i(\kappa _k+\kappa _l)[a_k^+,a_l^+],`$
the vector $`[a_k^+,a_l^+]`$, if non-zero, is again an eigenvector of $`\stackrel{~}{𝒥}_\xi `$ corresponding to a ‘positive’ eigenvalue $`i(\kappa _k+\kappa _l)`$.
(d) The tensor
$$\gamma _\xi ^𝒪(A,B)=\eta _\xi ^𝒪(A,𝒥_\xi (B))$$
is clearly $`U()`$-invariant. From (50) and (51) it follows that $`𝒥_\xi ^\times =𝒥_\xi `$. Since $`\stackrel{~}{𝒥}`$ and $`𝒥`$ clearly commute, $`𝒥_\xi ([A,\xi ])=[𝒥_\xi (A),\xi ]`$, in view of (17),
$`\eta _\xi ^𝒪([A,\xi ],𝒥_\xi ([B,\xi ]))`$ $`=`$ $`[A,\xi ],𝒥_\xi (B)_{u^{}()}=𝒥_\xi ([A,\xi ]),B_{u^{}()}`$
$`=`$ $`\eta _\xi ^𝒪(𝒥_\xi ([A,\xi ]),[B,\xi ]).`$
This immediately implies that $`\gamma ^𝒪`$ is symmetric and proves (49). But (17) implies also that
$`\gamma _\xi ^𝒪([A,\xi ],[A,\xi ])`$ $`=`$ $`\eta _\xi ^𝒪([A,\xi ],𝒥_\xi ([A,\xi ]))=[A,\xi ],𝒥_\xi (A)_{u^{}()}`$
$`=`$ $`A,\stackrel{~}{𝒥}_\xi 𝒥_\xi (A)_{u^{}()}.`$
But
$$\stackrel{~}{𝒥}_\xi 𝒥_\xi =\left(\stackrel{~}{𝒥}^2\right)^{\frac{1}{2}}$$
is a positive operator, so
$$\gamma _\xi ^𝒪([A,\xi ],[A,\xi ])>0$$
for $`[A,\xi ]0`$.
Finally, if $`\xi `$ is a projector, $`\xi ^2=\xi `$, then (cf. (22))
$$\stackrel{~}{𝒥}_\xi ^2([A,\xi ])=[A,\xi ],$$
so $`𝒥_\xi =\stackrel{~}{𝒥}_\xi `$ and (cf. (5))
$$\gamma _\xi ^𝒪([A,\xi ],[B,\xi ])=[A,\xi ],𝒥_\xi (B)_{u^{}()}=[A,\xi ],[B,\xi ]_{u^{}()}.$$
We have some similar results for the tensor $`\stackrel{~}{}`$ which however are not completely analogous, since the distribution $`D_R`$ is not globally integrable. The proofs are analogous, so we omit them.
###### Theorem 8
(a) The image $`D_R(\xi )`$ of $`\stackrel{~}{}_\xi `$ is the orthogonal complement of $`\text{Ker}(\stackrel{~}{}_\xi )`$.
(b) $`\stackrel{~}{}_\xi ^2`$ is a self-adjoint (with respect to $`,_u^{}`$) and positively defined operator on $`D_R(\xi )`$.
(c) The $`(1,1)`$-tensor $``$ on $`u^{}()`$ defined by
$$_\xi (A)=|(\stackrel{~}{}_\xi )_{D_R(\xi )}|^1\stackrel{~}{}_\xi (A)$$
(55)
satisfies $`^3=`$.
###### Corollary 5
The distribution $`D_0`$ is the image of $`𝒥_\xi _\xi =_\xi 𝒥_\xi `$. In other words, $`D_0(\xi )=\{[A,\xi ^2]:Au^{}()\}`$. Moreover, the foliation $`_0`$ is $`U()`$-invariant, $`𝒥`$-invariant and $``$-invariant, so that $`𝒥`$ and $``$ induce on leaves of $`_0`$ a complex and a product structure, respectively. The leaves of the foliation $`_0`$ are also canonically symplectic manifolds with symplectic structures being restrictions of symplectic structures on the leaves of $`_\mathrm{\Lambda }`$, so the leaves of $`_0`$ are Kähler submanifolds of the $`U()`$-orbits in $`u^{}()`$.
Proof. The image of $`𝒥_\xi _\xi =_\xi 𝒥_\xi `$ is clearly contained in $`D_0`$. Conversely, let $`BD_0(\xi )=D_\mathrm{\Lambda }D_𝐑`$. According to (46), $`D_0(\xi )`$ is invariant with respect to both: $`𝒥_\xi `$ and $`_\xi `$ and $`𝒥_\xi `$ and $`_\xi `$ are injective, thus surjective, on $`D_0(\xi )`$. The distribution $`D_0`$ is therefore generated by vector fields $`X_A(\xi )=[A,\xi ^2]`$. It is a matter of simple calculations to show that these vector fields commute with the fundamental vector fields $`\mathrm{\Lambda }_B`$ of the $`U()`$ as $`[X_A,\mathrm{\Lambda }_B]_{vf}=X_{[B,A]}`$ that shows $`U()`$ invariance of $`D_0`$. One can also easily seen that the restrictions of $`\xi ^𝒪`$ to the leaves of $`_0`$ contained in $`𝒪`$ are non-degenerate. It follows also directly from the explicit calculations we present below.
Let us explain the above theorem in local coordinates, i.e. for the case of matrices. Suppose that $`\xi =diag(\lambda _1,\mathrm{}\lambda _n)u^{}(n)`$ is a diagonal matrix. For simplicity, it is better to start already with the complexified structures, i.e. with $`gl(n)=u^{}(n)𝐂`$ equipped with the bracket $`[a,b]=\frac{1}{i}(abba)`$ and the Hermitian product $`a,b_{gl}=\frac{1}{2}\text{Tr}(a^{}b)`$, so that $`u^{}(n)`$ is a real Lie subalgebra in $`gl(n)`$ with the induced scalar product. Let $`E_l^k`$ be the matrix whose the only non-zero entry is $`1`$ at $`k`$th row and $`l`$th column. We have
$$E_l^k,E_s^r_{gl}=\frac{1}{2}(\delta _k^r\delta _s^l),$$
(56)
$$[E_l^k,E_s^r]=i(\delta _l^rE_s^k\delta _s^kE_l^r),$$
(57)
and
$$[E_l^k,E_s^r]_+=\delta _l^rE_s^k+\delta _s^kE_l^r,$$
(58)
so that
$$\stackrel{~}{𝒥}_\xi (E_l^k)=[E_l^k,\xi ]=i(\lambda _k\lambda _l)E_l^k.$$
(59)
and
$$\stackrel{~}{}_\xi (E_l^k)=[E_l^k,\xi ]_+=(\lambda _k+\lambda _l)E_l^k.$$
(60)
In particular,
$$\stackrel{~}{𝒥}_\xi \stackrel{~}{}_\xi (E_l^k)=[E_l^k,\xi ^2]=i(\lambda _k^2\lambda _l^2)E_l^k.$$
(61)
Consequently,
$$\stackrel{~}{𝒥}_\xi ^2(E_l^k)=(\lambda _k\lambda _l)^2E_l^k$$
(62)
and
$$\stackrel{~}{}_\xi ^2(E_l^k)=(\lambda _k+\lambda _l)^2E_l^k,$$
(63)
so that
$$𝒥_\xi (E_l^k)=isgn(\lambda _k\lambda _l)E_l^k.$$
(64)
and
$$_\xi (E_l^k)=sgn(\lambda _k+\lambda _l)E_l^k.$$
(65)
The (complexified) tangent space $`\text{T}_\xi 𝒪𝐂`$ is spanned by those $`E_l^k`$ for which $`\lambda _k\lambda _l0`$, the space $`D_R(\xi )𝐂`$ is spanned by those $`E_l^k`$ for which $`\lambda _k+\lambda _l0`$, the space $`D_0(\xi )𝐂`$ is spanned by those $`E_l^k`$ for which $`\lambda _k^2\lambda _l^20`$, and the distribution $`𝒩`$ mentioned in the proof of the theorem is spanned by $`E_l^k`$ for which $`\lambda _k\lambda _l>0`$. The complexified symplectic form reads
$$\eta _\xi ^𝒪(i(\lambda _k\lambda _l)E_l^k,i(\lambda _r\lambda _s)E_s^r)=i(\lambda _k\lambda _l)E_l^k,E_s^r_{gl}=i(\lambda _k\lambda _l)\frac{1}{2}(\delta _l^r\delta _s^k),$$
i.e.
$$\eta _\xi ^𝒪(E_l^k,E_s^r)=\frac{1}{2i(\lambda _r\lambda _s)}(\delta _l^r\delta _s^k),$$
(66)
and the complexified Riemannian form
$$\gamma _\xi ^𝒪(E_l^k,E_s^r)=\eta _\xi ^𝒪(E_l^k,𝒥_\xi (E_s^r))=\frac{1}{2|\lambda _r\lambda _s|}(\delta _k^r\delta _s^l).$$
(67)
As a basis in $`u^{}(n)`$ let us take
$$A_l^k=E_l^k+E_k^l,kl,B_l^k=iE_l^kiE_k^l,k<l.$$
(68)
It is easy to see that this is an orthonormal basis and that
$$𝒥_\xi (A_l^k)=sgn(\lambda _k\lambda _l)B_l^k,𝒥_\xi (B_l^k)=sgn(\lambda _l\lambda _k)A_l^k.$$
(69)
and
$$_\xi (A_l^k)=sgn(\lambda _k+\lambda _l)A_l^k,𝒥_\xi (B_l^k)=sgn(\lambda _l+\lambda _k)B_l^k.$$
(70)
Moreover
$$\eta _\xi ^𝒪(B_l^k,A_s^r)=\frac{\delta _r^k\delta _s^l}{(\lambda _k\lambda _l)},\eta _\xi ^𝒪(B_l^k,B_s^r)=\eta _\xi ^𝒪(A_l^k,A_s^r)=0,\lambda _k\lambda _l,\lambda _r\lambda _s0$$
(71)
and
$$\gamma _\xi ^𝒪(B_l^k,A_s^r)=0,\gamma _\xi ^𝒪(B_l^k,B_s^r)=\gamma _\xi ^𝒪(A_l^k,A_s^r)=\frac{\delta _r^k\delta _s^l}{|\lambda _k\lambda _l|},\lambda _k\lambda _l,\lambda _r\lambda _s0.$$
(72)
In other words
$$\eta _\xi ^𝒪=\underset{\lambda _k\lambda _l0}{}\frac{1}{(\lambda _k\lambda _l)}\text{d}b_l^k\text{d}a_l^k,$$
(73)
and
$$\gamma _\xi ^𝒪=\underset{\lambda _k\lambda _l0}{}\frac{1}{|\lambda _k\lambda _l|}(\text{d}b_l^k\text{d}b_l^k+\text{d}a_l^k\text{d}a_l^k),$$
(74)
where
$$b_l^k=B_l^k,_u^{},a_l^k=A_l^k,_u^{}$$
are coordinates on $`u^{}(n)`$ such that $`B_l^k=_{b_l^k},A_l^k=_{a_l^k}`$. The reduction of the symplectic form $`\eta ^𝒪`$ to the leaves of the foliation $`_0`$
$$(\eta _\xi ^𝒪)__0=\underset{\lambda _k^2\lambda _l^20}{}\frac{1}{(\lambda _k\lambda _l)}\text{d}b_l^k\text{d}a_l^k,$$
(75)
is clearly non-degenerate and constitutes, together with the reduced Riemannian structure
$$(\gamma _\xi ^𝒪)__0=\underset{\lambda _k^2\lambda _l^20}{}\frac{1}{|\lambda _k\lambda _l|}(\text{d}b_l^k\text{d}b_l^k+\text{d}a_l^k\text{d}a_l^k),$$
(76)
a Kähler structure.
Remark. Of course, when $`\xi `$ is a projector, then $`\lambda _k=1,0`$, so $`\lambda _k\lambda _l0|\lambda _k\lambda _l|=1`$ and $`\gamma _\xi ^𝒪`$ reduces to the canonical scalar product. Note also that the leaves of $`_\mathrm{\Lambda }`$ and $`_0`$ through $`\xi `$ coincide, except for the rare case when there are $`\lambda ,\lambda ^{}0`$ in the spectrum of $`\xi `$ such that $`\lambda +\lambda ^{}=0`$. In particular, the foliations $`_\mathrm{\Lambda }`$ and $`_0`$ coincide when reduced to the subset $`𝒫()`$ of non-negative operators or to the set $`𝒟()`$ of density states. On such leaves the product structure $``$ is trivially the identity.
## 6 Composite systems and separability
Suppose now that our Hilbert space has a fixed decomposition into the tensor product of two Hilbert spaces $`=^1^2`$. This additional input is crucial in studying composite quantum systems and it has a great impact on the geometrical structures we have considered. The rest of this paper will be devoted to related problems.
Observe first that the tensor product map
$$:^1\times ^2=^1^2$$
(77)
associates the product of rays with a ray, so it induces a canonical imbedding on the level of complex projective spaces
$`\text{Seg}:P^1\times P^2`$ $``$ $`P=P(^1^2),`$ (78)
$`(x^1x^1,x^2x^2)`$ $``$ $`x^1x^2x^1x^2.`$ (79)
This imbedding of product of complex projective spaces into the projective space of the tensor product is called in the literature the Segre imbedding . The elements of the image $`\text{Seg}(P^1\times P^2)`$ in $`P=P(^1^2)`$ are called separable pure states (with respect to the decomposition $`=^1^2`$).
The Segre imbedding is related to the (external) tensor product of the basic representations of the unitary groups $`U(^1)`$ and $`U(^2)`$, i.e. with the representation of the direct product group in $`=^1^2`$,
$`U(^1)\times U(^2)(\rho ^1,\rho ^2)`$ $``$ $`\rho ^1\rho ^2U()=U(^1^2),`$
$`(\rho ^1\rho ^2)(x^1x^2)`$ $`=`$ $`\rho ^1(x^1)\rho ^2(x^2).`$
Note that $`\rho ^1\rho ^2`$ is unitary, since the Hermitian product in $``$ is related to the Hermitian products in $`^1`$ and $`^2`$ by
$$x^1x^2,y^1y^2_{}=x^1,y^1_^1x^2,y^2_^2.$$
(80)
The above group imbedding which, with some abuse of notation, we will denote by
$$\text{Seg}:U(^1)\times U(^2)U(),$$
gives rise to the corresponding imbedding of Lie algebras
$$\text{Seg}:u(^1)\times u(^2)u(),$$
or, by our identification, of their duals
$$\text{Seg}:u^{}(^1)\times u^{}(^2)u^{}().$$
(81)
The original Segre imbedding is just the latter map reduced to pure states. In fact, a more general result holds true.
###### Proposition 5
The imbedding (81) maps $`𝒫^k(^1)\times 𝒫^l(^2)`$ into $`𝒫^{kl}(^1^2)`$ and $`𝒟^k(^1)\times 𝒟^l(^2)`$ into $`𝒟^{kl}(^1^2)`$.
Proof. Let us take $`A^1𝒫^k(^1)`$ and $`A^2𝒫^l(^2)`$. Using bases of eigenvectors $`(e_i^1)`$ of $`A^1`$ and $`(e_j^2)`$ of $`A^2`$ to construct a basis $`(e_i^1e_j^2)`$ of eigenvectors of $`A^1A^2`$, one easily sees that the elements of the spectrum of $`A^1A^2`$ (counted with multiplicities) are $`\lambda _i\lambda _j^{}`$, where $`A^1(e_i^1)=\lambda _ie_i^1`$ and $`A^2(e_j^2)=\lambda _j^{}e_j^2`$, so that $`A^1A^2=\text{Seg}(A^1,A^2)`$ is non-negatively defined and has rank $`kl`$. If $`A^1,A^2`$ have trace 1, i.e. $`_i\lambda _i=1`$ and $`_j\lambda _j^{}=1`$, then $`_{i,j}\lambda _i\lambda _j^{}=_i\lambda _i_j\lambda _j^{}=1`$.
Let us denote the image $`\text{Seg}(𝒟^k(^1)\times 𝒟^l(^2))`$ by $`𝒮^{k,l}(^1^2)`$, the set $`𝒮^{1,1}(^1^2)`$ of separable pure states simply by $`𝒮^1(^1^2)`$, and the convex hull
$$\text{conv}\left(\text{Seg}\left(𝒟(^1)\times 𝒟(^2)\right)\right)$$
of the subset $`\text{Seg}\left(𝒟(^1)\times 𝒟(^2)\right)`$ of all separable states in $`u^{}()`$ by $`𝒮(^1^2)`$. The states from
$$(^1^2)=𝒟(^1^2)𝒮(^1^2),$$
i.e. those which are not separable, are called entangled states.
###### Proposition 6
The convex set $`𝒮(^1^2)`$ of separable states is the convex hull of the set $`𝒮^1(^1^2)`$ of separable pure states and $`𝒮^1(^1^2)`$ is exactly the set of extremal points of $`𝒮(^1^2)`$. Moreover, $`𝒮^1(^1^2)`$, thus $`𝒮(^1^2)`$, is invariant with respect to the canonical $`U(^1)\times U(^2)`$-action on $`u^{}(^1^2)`$,
$$(T_1,T_2)A=(T_1T_2)A(T_1T_2)^{}.$$
Proof. Let us start with showing that the convex hull of $`𝒮^1(^1^2)`$ contains $`\text{Seg}\left(𝒟(^1)\times 𝒟(^2)\right)`$ thus equals $`𝒮(^1^2)`$. Indeed $`𝒟^1(^i)`$ is the set of extreme points of $`𝒟(^i)`$, $`i=1,2`$, so that any $`A^i𝒟(^i)`$ is a convex combination $`A^i=t_s^i\rho _s^i`$ of elements $`\rho _s^i𝒟^1(^i)`$, $`i=1,2`$. Hence, $`A^1A^2`$ is the convex combination
$$A^1A^2=\underset{s,s^{}}{}t_s^1t_s^{}^2\rho _s^1\rho _s^{}^2.$$
On the other hand, every state $`\rho ^1\rho ^2`$, $`\rho ^i𝒟(^i)`$, $`i=1,2`$, is in $`𝒟^1(^1^2)`$, i.e. it is an extremal point of $`𝒟(^1^2)`$. Therefore it cannot be written as a non-trivial convex combination of elements from $`𝒟(^1^2)`$, thus from a smaller set $`𝒮(^1^2)`$. The invariance is obvious, since
$$(T_1T_2)(\rho _1\rho _2)(T_1^{}T_2)^{}=(T_1\rho _1T_2^{})(T_1\rho _2T_2^{})$$
and $`(T_1\rho _iT_2^{})𝒟^1(^i)`$ for $`\rho _i𝒟^1(^i)`$.
Since we are working in a finite-dimensional space, the closeness of the corresponding hulls is automatic that can be derived from the following lemma.
###### Lemma 5
If $`V`$ is an $`n`$-dimensional real vector space and $`x`$ is a convex combination $`x=_{i=1}^mt_ix_i`$ of certain points of $`V`$, then $`x`$ is a convex combination of at most $`(n+1)`$ points among $`x_i`$’s.
Proof. It suffices to prove that $`x`$ is a convex combination of $`(m1)`$ of $`x_i`$’s, provided $`m>n+1`$. Of course, we can assume that all $`t_i>0`$. If $`m>n+1`$, then there are $`a_i𝐑`$, not all equal 0, such that $`_1^ma_i=0`$ and $`_1^ma_ix_i=0`$. There is $`i_0`$ such that $`|a_{i_0}/t_{i_0}|`$ is maximal among $`|a_{i_0}/t_{i_0}|`$, $`i=1,\mathrm{},m`$. We can assume without loss of generality that $`i_0=m`$. Hence
$$x=\underset{i=1}{\overset{m1}{}}\left(t_i\frac{a_it_m}{a_m}\right)x_i$$
and the above combination is convex, since $`(t_i\frac{a_it_m}{a_m})0`$ and
$$\underset{1}{\overset{m1}{}}(t_i\frac{a_it_m}{a_m})=\underset{1}{\overset{m}{}}(t_i\frac{a_it_m}{a_m})=\underset{1}{\overset{m}{}}t_i=1.$$
###### Proposition 7
The convex hull $`\text{conv}(E)`$ of a compact subset $`E`$ in a finite dimensional real vector space $`V`$ is compact.
Proof. Suppose that the dimension of the space is $`n`$ and denote by $`\mathrm{\Delta }_{n+1}`$ the compact $`(n+1)`$-dimensional simplex
$$\mathrm{\Delta }_{n+1}=\{t=(t_1,\mathrm{},t_{n+1}):t_i0,\underset{1}{\overset{n+1}{}}t_i=1\}.$$
According to the above lemma, $`\text{conv}(E)`$ is the image of the compact set $`\mathrm{\Delta }\times E\times \mathrm{}\times E`$ ($`E`$ appears in the product $`(n+1)`$-times) under the continuous map
$$\mathrm{\Delta }\times E\times \mathrm{}\times E(t,x_1,\mathrm{},x_{n+1})\underset{1}{\overset{n+1}{}}t_ix_iV.$$
###### Corollary 6
The set $`𝒮(^1^2)`$ is a compact subset of $`u^{}(^1^2)`$.
The entangled states play an important role in quantum computing and one of main problems is to decide effectively whether a given composite state is entangled or not. An abstract measurement of entanglement can be based on the following observation (see also Ref. )
Let $`E`$ be the set of all extreme points of a compact convex set $`K`$ in a finite-dimensional real vector space $`V`$ and let $`E_0`$ be a compact subset of $`E`$ with the convex hull $`K_0=\text{conv}(E_0)K`$. For every non-negative function $`f:E𝐑_+`$ define its extension $`f_K:K𝐑_+`$ by
$$f_K(x)=\underset{x={\scriptscriptstyle t_i\alpha _i}}{inf}t_if(\alpha _i),$$
(82)
where the infimum is taken with respect to all expressions of $`x`$ in the form of convex combinations of points from $`E`$. Recall that that, according to Krein-Milman theorem, $`K`$ is the convex hull of its extreme points.
###### Theorem 9
For every non-negative continuous function $`f:E𝐑_+`$ which vanishes exactly on $`E_0`$ the function $`f_K`$ is convex on $`K`$ and vanishes exactly on $`K_0`$
Proof. It is completely obvious that $`f_K`$ vanishes on the convex hull of $`E_0`$. The function $`f_K`$ is convex, since for every convex combination $`x=t_iy_i`$ of points of $`K`$ and every $`\epsilon >0`$ we can find extreme points $`\alpha _j`$ with convex combinations $`y_i=s_i^j\alpha _j`$ and $`f_K(y_i)>s_i^jf(\alpha _j)\epsilon `$. Hence
$$f_K(t_iy_i)=f_K(t_is_i^j\alpha _j)t_is_i^jf(\alpha _j)<t_i(f(y_i)+\epsilon )=t_if_K(y_i)+\epsilon .$$
Due to arbitrariness of $`\epsilon >0`$ we get
$$f_K(t_iy_i)t_if_K(y_i).$$
Note finally that $`f_K`$ vanishes exactly on $`K_0`$. Indeed $`K_0`$ is compact due to proposition 7 and if $`xK_0`$, then $`x`$ and $`K_0`$ can be separated by a hyperplane, i.e. there is a linear functional $`\phi :V𝐑`$ such that $`\phi (x)=a>0`$ and $`\phi `$ is negative on $`K_0`$. Denote by $`E_1`$ the (compact) set of those points from $`E`$ on which $`\phi `$ takes non-negative values and by $`F`$ the minimum of $`f`$ on $`E_1`$. Of course, $`F>0`$, since $`E_1E_0=\mathrm{}`$. Let $`M𝐑`$ be the maximum of $`\phi `$ on $`E`$. Of course, $`M>0`$. For any realization $`x=t_i\alpha _i`$ of $`x`$ as a convex combination of points of $`E`$ we have
$$a=\phi (x)=\underset{i}{}t_i\phi (\alpha _i)\underset{\alpha _iE_1}{}t_i\phi (\alpha _i)M\underset{\alpha _iE_1}{}t_i.$$
On the other hand,
$$\underset{i}{}t_if(\alpha _i)\underset{\alpha _iE_1}{}t_if(\alpha _i)F\underset{\alpha _iE_1}{}t_i\frac{aF}{M},$$
so $`f_K(x)\frac{aF}{M}>0`$.
###### Corollary 7
Let $`F:𝒟^1(^1^2)𝐑_+`$ be a continuous function which vanishes exactly on $`𝒮^1(^1^2)`$. Then
$$\mu =F_{𝒟(^1^2)}:𝒟(^1^2)𝐑_+$$
is a measure of entanglement, i.e. $`\mu `$ is convex and $`\mu (x)=0x𝒮(^1^2)`$. Moreover, if $`f`$ is taken $`U(^1)\times U(^2)`$-invariant, then $`\mu `$ is $`U(^1)\times U(^2)`$-invariant.
Proof. The first part is a direct consequence of Theorem 9. Also the invariance of $`\mu `$ is clear:
$$\mu (T\rho T^{})=\underset{W}{inf}(t_if(\alpha _i))=\underset{W^{}}{inf}(t_if(T\alpha _iT^{}))=\underset{W^{}}{inf}(t_if(\alpha _i))=\mu (\rho ),$$
where $`T`$ is in the corresponding group,
$$W=\{(t_i,\alpha _i):T\rho T^{}=t_i\alpha _i,\alpha _i𝒮^1(^1^2),t_i0,t_i=1\},$$
and
$$W^{}=\{(t_i,\alpha _i):\rho =t_i\alpha _i,\alpha _i𝒮^1(^1^2),t_i0,t_i=1\}.$$
A careful study of the geometry of $`u^{}((^1^2)`$ and criteria of entanglement we postpone to a separate paper.
## 7 Acknowledgements
We would like to thank V. S. Varadarajan for useful discussions on the contents of this paper. This work was supported by the Polish Ministry of Scientific Research and Information Technology under the (solicited) grant No PBZ-Min-008/P03/03 and partially supported by PRIN SINTESI. |
warning/0507/cond-mat0507299.html | ar5iv | text | # Dynamic hysteresis in Finemet thin films
## I Introduction
The physics of thin and ultrathin magnetic films has been extensively studied in the recent past, because of its great importance in several applications, ranging from multilayers to high frequency devices. For this reason, many recent papers have been devoted to measure the magnetization reversal dynamics in two dimensional structures, revealing the existence of universal features and scale-invariant properties of the hysteresis loops . Despite these efforts, a general description of these features is still an open problem, as most experimental results are still to be interpreted in the framework of the existing models . In particular, the dynamic hysteresis loop area $`A`$ is often assumed to scale as $`AH_0^\alpha \omega ^\beta T^\gamma `$, where $`H_0`$ is the amplitude of sinusoidal external field of frequency $`\omega `$, $`T`$ is the temperature, and $`\alpha `$, $`\beta `$, $`\gamma `$ are three scaling exponents. As a matter of fact, experimental evaluation of the exponents $`\alpha `$ and $`\beta `$ in the low dynamic regime spans a quite large range from 0 to 0.8, with a general higher value for thinner films (see also Tab. I of ). The proposed theoretical models roughly span the same range, so that a clear identification of the fundamental properties of magnetization dynamics seems far to be reached. In order to investigate this complicated problem offering a new perspective, we present a series of dynamic hysteresis measurements on Finemet thin films by using the magneto-optical Kerr effect (MOKE), as employed in all the studies presented in the literature, and the fluxmetric inductive method, using a pick-up coil wound around the sample. This enables us to investigate the hysteresis properties not only considering the magnetization changes of the surface within the laser spot area, but also those of the total volume of the sample; this is particularly important in order to check the dependence of the loop area on the film thickness, and to understand the true nature of the magnetization dynamics. Quite unexpectedly, MOKE hysteresis loops show a remarkable variability with the frequency, about one order of magnitude higher than the inductive ones; on the other hand, both methods give a similar dependence on the frequency. We try to interpret these results with a simple domain wall depinning model which can be solved analytically, giving reason of the general behavior of the data.
## II Materials and measurement methods
Films having nominal composition Fe<sub>73.5</sub>Cu<sub>1</sub>Nb<sub>3</sub>Si$`{}_{1}{}^{}3.5`$B<sub>9</sub> have been deposited on glass substrates by rf magnetron sputtering under a 5 mTorr Ar atmosphere at room temperature. The sample thickness, ranging from about 21 nm to 5 $`\mu `$m, is measured by angle X-ray diffractometry, which also confirms the amorphous state of the samples. The hysteresis loops of the in-plane magnetization are measured in the as-prepared materials as a function of the applied field (up to a few kA/m) at 100 Hz, and as a function of the frequency. The longitudinal MOKE measurements are performed with an optical bench equipped with an He-Ne laser light source, covering a sample surface of about 1 mm<sup>2</sup>, and a photodiode having a frequency cutoff above 150kHz. Samples are cut to a maximum size of about 2 x 2 mm and inserted in a Helmholtz coil setup giving a maximum field of about 15 kA/m. In this configuration, we could perform measurements up to 300 Hz.
The fluxometric measurements of hysteresis loops are performed on larger samples, usually cut within an homogeneous region of 3 by 1.5 cm. For these samples, we prepared a 10 cm long solenoid with 720 coils and a N/L value of 13200. The sinusoidal applied field is measured detecting the voltage over a calibrated 1 $`\mathrm{\Omega }`$ resistance. The induced flux is detected with two sets of 50 coils wound around the sample and covering an area of about 1 cm; the two sets are wired in order to cancel out the air flux and directly detect the film magnetization changes. As the film cross section is much smaller than the area of the coils, a perfect cancellation of the air flux is often difficult, requiring a continuous adjustment of area of one of the sets. This problem further increases at high frequencies, due to the different coupling between the wires, making it hard to perform the measurement. We thus get the full cancellation by numerically subtracting a sinusoidal wave with proper amplitude and phase from the induced signal. This procedure has a certain degree of arbitrariness, as small changes of the sinusoidal amplitude and phase give slightly different loop shapes. In all cases, the loop area and the coercive field are not substantially affected.
## III Experimental results
The first fundamental result of both kind of measurements is the existence of a well defined static hysteresis, as usually found in magnetic materials (see for instance ). Therefore, the loop area $`A`$ is better described by
$$AA_0+H_0^\alpha \omega ^\beta T^\gamma $$
(1)
where the static loop area $`A_0`$ is estimated using data at the lowest frequencies. Clearly, the choice of (1) completely changes the experimental estimation of critical exponents. We believe that this simple observation could explain the so called dynamic transition, a sharp change in the value of the exponent $`\beta `$ at intermediate frequencies. Due to the usual large value of $`A_0`$ with respect to variation of the data with the frequency, log-log plots can mimic a dynamic transition when data actually follow a simple power law as assumed in (1). Our results and the theoretical analysis should help to clarify this important point.
The MOKE hysteresis loop area for different sample thickness show large variations with the frequency, as shown in Fig. 1 for the thickest (5 $`\mu `$m) and thinnest (21 nm) samples. Visual inspection of this Figure would suggest a dynamic transition around a few Hz, at least, for the 5 $`\mu `$m sample. As a matter of fact, the data show a much simpler behavior: in fact, the loop area $`A`$ follows a simple law of the type $`AA_0+k\omega ^{0.5}`$, where $`A_0`$ is estimated using low frequency data.
Surprisingly, hysteresis loops obtained with the fluxometric setup do not show the same frequency variability. As reported in Figs. 2 for the 21 nm film, the loop area changes only about 10% with respect to the static value $`A_0`$. On the other hand, when plotted as a function of the square root of frequency, these data follow reasonably well a linear behavior. This suggests a common mechanism responsible for the magnetization change, which we try to interpret considering the domain wall depinning models, as discussed below.
Inductive measurements of the hysteresis loops at 100 Hz as a function of the applied field amplitude do not show any scaling behavior as given by (1) or similar. Simple squared loops as the ones shown in Fig. 2 are observed only for the thinner films, while more complicated shapes appear for thicker samples, showing a clear evidence of multi-domains magnetization processes. We postpone the discussion of these rich but complex features to a further longer publication.
## IV Model
To understand the frequency dependence of our experimental data, we employ the domain wall depinning model described in Ref. . Under the assumption that hysteresis is mainly due to domain wall motion, we consider a phenomenological law for the wall velocity given by
$$v(H)=\mu (|H|H_p)\theta (|H|H_p),$$
(2)
where $`\mu `$ the wall mobility, $`H_p`$ is the depinning field, $`\theta `$ is the step function, and the applied field is $`H=H_0sin(\omega t)`$. Following , we solve (2) and compute the coercive field and the loop area, which at low frequency are given by
$$H_c(\omega )H_p+(2L\omega \sqrt{H_0^2H_p^2}/\mu )^{1/2},$$
(3)
where $`L`$ is the sample size, in the case of a single domain wall, or the typical distance between domain walls in a more general case. The lower branch of the hysteretic loop is then given by
$$M(h)=\{\begin{array}{cc}M_s\hfill & H<H_p\hfill \\ M_s\left(\frac{\mu }{L\omega }\frac{(HH_p)^2}{\sqrt{H_0^2H_p^2}}1\right)\hfill & HH_p\hfill \end{array}$$
(4)
where the second equation is valid as long as $`M<M_s`$. Fig. 3 displays the loop shapes for different values of frequency $`\omega `$. The loop area is thus easily computed and, for low frequencies, is given by
$$AA_0+M_s\frac{8}{3}\sqrt{\frac{2\omega L}{\mu }\sqrt{H_0^2H_p^2}},$$
(5)
where $`A_0=4M_sH_p`$ is the area of the loop in the quasistatic limit.
## V Discussion and conclusion
The simple model shown above gives a simple $`\sqrt{\omega }`$ dependence of the loop area (5), which we believe can account for some experimental data found in the literature and interpreted assuming the presence of a dynamic transition. The value $`\beta =0.5`$ is a strict consequence of the linear form of (2): in a more general case, shown in detail in , one gets $`\beta 0.5`$. It is worth noting that the model above describes the magnetization dynamics of a single domain wall having a well defined depinning field $`H_p`$. The model does not include the effects of any random disorder, the multiplicity of the domains, or the field-dependent nucleation of domain walls. As a consequence, it cannot describe minor hysteresis loops or other features, such as more complicated shapes of the loops. Despite these limitations, we believe it can account for the general behavior of the dynamic hysteresis, and could be successfully applied to describe thin films showing simple magnetization dynamics.
Using the results of the model, we found that our experimental data obtained with complementary measurement techniques are the consequence of the same magnetization dynamics. That means that the MOKE measurements on the surface are related to those of the volume, given by the inductive method. While this is expected in thin films, it is not necessarily valid for the thickest samples, for instance in the case of a few microns. It is not clear anyway what can account for the large differences in the change of the loop area. It is worth noting that the MOKE measurement uses a laser spot area close to the sample dimension so, in principle, it should detect the magnetization changes of the entire sample, as long as we consider the thinnest films. Therefore, we should observe similar variations using the two methods. On the other hand, one can suggest that samples with the same thickness but different later dimensions could not have the same magnetization dynamics, because of the different role of the demagnetizing fields, or of the number of domain walls. This possibility has been ruled out by performing fluxometric measurements on the same sample used in MOKE, having thickness 5 $`\mu `$m. In this case the fluxometric measurements show the same behavior described in Fig. 2. Unfortunately, measurements at lower thickness are not possible due to the very low intensity of the induced signal, given the reduce cross section available. However, we feel confident that this result proves that the sample lateral dimensions are not relevant, also considering that the data reported in literature refer to dependences only on the sample thickness and not, more generally, to the other two dimensions.
In summary, we have shown in this paper that the dynamic hysteresis of Fe<sub>73.5</sub>Cu<sub>1</sub>Nb<sub>3</sub>Si$`{}_{1}{}^{}3.5`$B<sub>9</sub> films exhibit different behavior as a function of sample thickness and magnetizing field frequency, and that the domain wall depinning model accounts for the general behavior of the data |
warning/0507/cond-mat0507207.html | ar5iv | text | # Structure and time-dependence of quantum breathers
## I Introduction
Nonlinearity creates localized structures through a variety of mechanisms. In a discrete lattice, extreme displacement of one or a few atoms pushes them to high oscillatory frequencies outside the range of ordinary (linear) phonons, so that the energy sequestered in these large excitations does not spread throughout the lattice. Such an excitation, known as a *discrete breather* (DB) or *intrinsic localized mode* (ILM), is a kind of soliton, although different in many ways from the moving, continuum soliton of Russell rus . The discrete breathers of primary interest to us, because of their relation to decay anomalies in doped alkali halides Schulman et al. (2002), are not “intrinsic” in the sense that they are induced by excitation of the impurity, and translational invariance has been lost. For convenience, and to minimize acronyms, we refer to the excitations we study as “breathers.”
There is an extensive literature on these breathers and an introduction can be found in Campbell et al. (January 2004). Because of the nonlinearity, a good deal of the theoretical work in this field is numerical. For systems well-approximated classically this is not a problem, but treating the system quantum mechanically demands potentially serious approximations.
An issue of experimental relevance is the lifetime of the quantum breather against decay. At the classical level, both analytical and numerical experience support an infinite lifetime, at least for the one-frequency breather Flach et al. (1994). In a diatomic lattice there are two breather frequencies (in the phonon gap and above the optical phonons) and although there are indications that combinations of the modes can allow energy to leave via the phonons, in our simulations this does not take place in any perceptible way.
But classical stability says little about its quantum counterpart. Quantum tunneling bypasses classical constraints, and for a translationally invariant system, guarantees that bound states that are not of finite support become bands. Indeed such bands—as well as tunneling—have been studied for breathers in a homogeneous lattice Wang et al. (1996); Fleurov et al. (1998); Flach and Fleurov (1997). Other breather quantization approaches have also been taken, including semiclassical and field-theoretical Schulman (2003); Konotop and Takeno (2001). In the semiclassical approach one uses Einstein-Brillouin-Keller quantization on the KAM torus-confined classical trajectory. Naturally this yields a stable state, since, as for a 1-D state hidden behind a barrier, the classical trajectory does not sense the possibility of escape. Nevertheless, in principle it is entirely plausible that corrections to semiclassical behavior, which generally begin at order $`\mathrm{}^2`$, would include the possibility of decay (although in one-dimension the tunneling decay correction is the much smaller $`\mathrm{exp}([\text{positive constant}]/\mathrm{})`$).
The focus of the present article is the stability of the quantized breather. At times it seems that the existence of even the classical breather is magical. Were this not a known phenomenon, one’s first reaction on finding this numerically should be to debug the program! Why then should one expect this “magic” to operate in the quantum case as well? Perhaps to first order in $`\mathrm{}`$, where semiclassical considerations are a guide, but a priori there is no reason that quantum corrections should not lead to decay channels.
Indeed, in a series of articles Hizhnyakov (1996, 1999); Hizhnyakov and Nevedrov (1997); Hizhnyakov et al. (2002) it has been claimed that such decay takes place, and estimates for its rate in alkali halides give lifetimes of about 10 ns. Decay is found to occur irrespective of dimension and, as far as we can tell, the decay rate is of order $`\mathrm{}`$, surprisingly large. We recall the framework used in these articles, since to some extent we will take a similar approach. It is typical of the classical breather that a small number of atoms (often just one) vibrate strongly while the others hardly move at all. This strong vibration should not be treated by perturbation theory, and one replaces this dynamic, nonlinear object by something more tractable. Refs. Hizhnyakov (1996, 1999); Hizhnyakov and Nevedrov (1997); Hizhnyakov et al. (2002) replace it by a fixed, rapidly oscillating forcing term which couples—including nonlinear coupling—to the phonon field of the crystal. Two approximations are inherent in this: nonlinear interactions other than those with the central atom are neglected and the motion of the central atom is treated classically, subject neither to quantum laws nor to the back reaction of the phonon field.
Recent work Flach et al. (2005) takes a similar perspective but comes to different conclusions regarding decay rates. For the infinite lattice the authors find no decay at all, with departures from this situation for finite $`N`$ (lattice size) going like $`1/N`$. In any case, this is in strong disagreement with nanosecond-scale lifetimes in physical crystals.
In the present article we treat the entire system quantum mechanically. We use two distinct approaches, having at their core fundamentally different approximations. In both cases however the unavoidably large impact of the breather is treated by introducing a *local mode*, either for comparison or for use as a zeroth-order system around which to perturb. As in other studies we drop nonlinear contributions except those that involve the particles most active in the breather.
Our first method employs the Feynman path integral. An early triumph of that technique was the calculation of polaron properties Feynman and Hibbs (1965); Feynman (1972), and the principle lying behind its success was that almost all forces involved were linear. Using the path integral most degrees of freedom could be integrated, leaving the polaron with a self-coupling that was non-local in time and reflected the back-reaction of the lattice to its motion. All this was at the quantum level. By dropping all nonlinear terms except that of the most active breather atom we arrive at a situation where the same method can be used. This also works if the active atom has only linear interactions, and allows us to compare properties of the breather to properties of a local mode, one whose motion is known to be fully localized. This is important because it is not enough to show the existence of eigenstates that correspond to the breather—one must show these eigenstates to be localized. What emerges from our calculations is that correlations of the breather-atom’s motion with other atoms in the lattice drop in the same way as those of a local mode, as the strength of the appropriate coupling constant is increased. The path integral method has the additional advantage that its extension to higher dimension should be possible. (Note that our use of the path integral is unrelated to certain other uses of functional integration for solitons, for example as described in the work of Faddeev and Korepin Faddeev and Korepin (1978).)
The second method is direct diagonalization of the Hamiltonian. Essentially the same calculation was performed in Wang et al. (1996), although they maintained translational invariance—they did not drop any nonlinear terms—and as a result obtained band structure. The difficulty of numerical diagonalization is that a cutoff on the size of the phonon Hilbert space is required. In a phonon basis, a quartic interaction couples to ever-higher phonon excitations, and perforce some of these must be neglected. Moreover, the cutoff must be fairly low. The elegance of the number-operator representation allows the analytic calculator to forget that the phonon spaces form a tensor product. Nevertheless, if $`n`$ phonons are considered, with cutoff Hilbert spaces of dimension $`N_c`$, the Hamiltonian will live in a space of $`N_{c}^{}{}_{}{}^{n}`$ dimensions, potentially a computational disaster. Our innovation is to introduce a local mode into the zeroth order Hamiltonian, one whose frequency is roughly that of the breather in which we are interested. In this way, when going over to nonlinear coupling very few phonons are excited other than those of the local mode, and a low cutoff does little harm.
Our conclusion from these calculations is that the quantum breather, at least in one dimension, is stable. (One-dimension, because of the intense distortion of the symmetry-breaking Jahn-Teller effect, is the case of physical interest to us Schulman et al. (2002).) From the path integral, the matching of correlations shows that the quantum breather has the same localization properties as a local mode, which is known to have an excitation profile (for atoms of the lattice) that drops off exponentially away from its source (as will be explicitly shown). For the direct diagonalization, the contribution of phonons other than those of the local mode to the true eigenstate is extremely small, on the order of $`10^3`$ for the parameters we use. We also address the issue of whether even very small contributions can lead to decay, which might seem to be implied by Fermi’s Golden Rule. It turns out that coupling to a continuum does not always mean decay to the continuum, as was found in Gaveau and Schulman (1995). In the present article we give explicit examples and in particular examine the time evolution of a state initially in an eigenstate of the local mode—which is *not* an eigenstate of the full Hamiltonian. It does not decay.
In Sec. III we show that dropping nonlinear terms, except for those affecting the central atom, leaves the classical mechanics qualitatively unchanged. Following that, in Sec. IV, we present the path integral approach. Direct diagonalization of the Hamiltonian appears in Sec. V and the associated time evolution in Sec. VI. A final section is devoted to discussion and conclusions.
## II Definitions and notation
The system is a ring of $`N+1`$ unit-mass atoms with Hamiltonian
$$H=\underset{k=0}{\overset{N}{}}\left\{\frac{1}{2}p_k^2+\frac{1}{2}\omega _s^2x_k^2+\frac{1}{2}\omega _0^2(x_kx_{k+1})^2+\frac{1}{4}\lambda x_k^4\right\},x,$$
(1)
Periodicity is expressed through $`x_0x_{N+1}`$ and mod-($`N`$+1) addition for atom labels. This is the nonlinearity studied by Wang et al. (1996). Another system, closer to our own models, is
$$H=\underset{k=0}{\overset{N}{}}\left\{\frac{1}{2}p_k^2+\frac{1}{2}\omega _s^2x_k^2+\frac{1}{2}\omega _0^2(x_kx_{k+1})^2+\frac{1}{4}\lambda (x_kx_{k+1})^4\right\},$$
(2)
also periodic. To avoid being awash in subscripts and superscripts we distinguish these by context, rather than by labels attached to $`H`$. As remarked, when a breather is present, one typically has a single atom where most of the energy is concentrated, with the other atoms relatively still. We take the dynamic atom to be #0 and consider an alternative Hamiltonian in which all nonlinear terms that do not involve $`x_0`$ are dropped. The corresponding Hamiltonians are
$`H`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{N}{}}}\left\{\frac{1}{2}p_k^2+\frac{1}{2}\omega _s^2x_k^2+\frac{1}{2}\omega _0^2(x_kx_{k+1})^2\right\}+\frac{1}{4}\lambda x_0^4,`$ (3)
$`H`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{N}{}}}\left\{\frac{1}{2}p_k^2+\frac{1}{2}\omega _s^2x_k^2+\frac{1}{2}\omega _0^2(x_kx_{k+1})^2\right\}+\frac{1}{4}\lambda \left[(x_0x_1)^4+(x_0x_N)^4\right].`$ (4)
For each of these we introduce a *local mode* Hamiltonian, namely a linear system which, at a level to be explored below, behaves similarly to those in Eqs. (3) and (4). They are
$`H_0`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{N}{}}}\left\{\frac{1}{2}p_k^2+\frac{1}{2}\omega _s^2x_k^2+\frac{1}{2}\omega _0^2(x_kx_{k+1})^2\right\}+\frac{1}{2}\omega _1^2x_0^2,`$ (5)
$`H_0`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{N}{}}}\left\{\frac{1}{2}p_k^2+\frac{1}{2}\omega _s^2x_k^2+\frac{1}{2}\omega _0^2(x_kx_{k+1})^2\right\}+\frac{1}{2}\omega _1^2\left[(x_0x_1)^2+(x_0x_N)^2\right].`$ (6)
As the zero subscript suggests, these will be the zero-order perturbation Hamiltonians in our numerical diagonalization. The perturbations will then have the form
$`V_I`$ $`=`$ $`\frac{1}{4}\lambda x_0^4\frac{1}{2}\omega _1^2x_0^2,`$ (7)
$`V_I`$ $`=`$ $`\frac{1}{4}\lambda \left[(x_0x_1)^4+(x_0x_N)^4\right]\frac{1}{2}\omega _1^2\left[(x_0x_1)^2+(x_0x_N)^2\right],`$ (8)
for the $`H_0`$’s of Eqs. (5) and (6), respectively.
All these systems have reflection symmetry about atom-0. Define $`P`$ by
$$P:k(Nk+1),k=0,\mathrm{},N_{\text{max}}\text{ with }N_{\text{max}}=\left[\frac{N+1}{2}\right],$$
(9)
and the square brackets represent “integer part of.” Classically an initial condition having this symmetry will retain it. Quantum mechanically, since $`V_I`$ shares this symmetry, only phonons even under $`P`$ need be considered.
## III Effect of dropping nonlinearity for “non-breather” atoms
Since we drop nonlinear terms that do not include $`x_0`$, we first study the impact of that change on the classical mechanics. The dynamical systems to be compared are
* $`H`$ of Eq. (1), “All nonlinear.”
* $`H`$ of Eq. (3), “One nonlinear” ($`x_0`$ in the equation).
* $`H_0`$ of Eq. (5), “Local mode” (around $`x_0`$).
We also make similar comparisons for nonlinear coupling:
* $`H`$ of Eq. (2), “All nonlinear.”
* $`H`$ of Eq. (4), “One nonlinear” ($`x_0`$ couples nonlinearly to its neighbors).
For the self-coupling cases, we first exhibit the kinetic energy as a function of time for a ring of 30 atoms. As for our subsequent graphs we have shifted the central atom to the middle of the ring to avoid cutting the breather. The kinetic energy is a clear way to see the breather and is shown in Fig. 1. In Fig. 2, we present position as a function of time for all three cases. Finally to illustrate the fact that differences do remain, we show the Fourier transform of the position of the central breather particle. Fig. 3 shows intensity, and it is clear that when all atoms see the nonlinear force there are more pronounced higher harmonics in the motion of the central atom as well. All show activity in the phonon band, notwithstanding the fact that all are stable.
The reason for the relative indifference of breather behavior to nonlinear couplings not involving #0 is that in any case atoms other than #0 have small amplitude and hardly feel the nonlinearity. It is certainly true that for smaller $`\lambda `$, when the breather spreads over a large number of atoms (as occurs in several cases in Mihóková et al. (2004)), this approximation would be inadequate. However, in this article we consider only the highly localized case so as to focus on the issue of in-principle quantum stability. Fig. 4, based on the full Hamiltonian Eq. (1), shows how small the nonlinear contributions to the energy are for any but the central breather atom. What are displayed are the time averages of $`\omega _s^2x^2`$ and $`\lambda x^4`$. It is clear that even for the nearest neighbors of the highly oscillatory atom the nonlinear contribution to the energy is negligible.
For the case of the truncation of nonlinear coupling, the parallel is not so close. The Hamiltonian, Eq. (2), gives rise to interesting patterns of behavior not seen in the self-coupling case. There is first the phenomenon of moving breathers Takeno and Hori (1991); Hori and Takeno (1992); Burlakov et al. (1990); Chubykalo et al. (1993); Flach and Willis (1994); Marín et al. (2001), which obviously will not occur when only one site and its neighbors are subject to nonlinearity. Movement can be suppressed by taking appropriate initial conditions (equal and opposite position displacements) so that the stationary breather consists of a pair of atoms oscillating with respect to each other. To reproduce this in a truncated Hamiltonian we would have to keep the nonlinearity for at least two atoms, which, for our later numerical diagonalization, puts a strain on computer memory. We therefore stayed with the form (4), which provides a local, nonlinear structure. The logic of Hizhnyakov (1996, 1999); Hizhnyakov and Nevedrov (1997); Hizhnyakov et al. (2002) predicts decay of this system when quantized, so that finding the eigenvectors of this system to be localized provides a further counterexample to that logic. In Fig. 5 we show the contrasting behavior of position as a function of time for the true breather and the truncated one. Apropos moving breathers, in Fig. 6 we present an illustration of a typical example, wandering around the ring in a seemingly random way.
### III.1 Quantum perspective
Quantum mechanically, dropping nonlinearity away from the center of the breather has a profound effect. It represents a loss of translational invariance, and, because of quantum tunneling, a qualitative change in the spectrum of the system. It follows that, in principle at least, for a translationally invariant system all breathers are nonlocal. The associated band structure for breather levels was found in Wang et al. (1996) and Ref. Fleurov (2003) has an extensive discussion of the significance of translational invariance.
It is therefore important to state what we mean in asserting that the breather is localized, and why our assertion is physically significant. If one considers a deeply bound atomic level, say the innermost shell of Cu when the atom is part of a crystal, that level—in principle—is part of a band. Practically, however, the band is irrelevant and the measurable properties of this level are independent of the near-zero tunneling probabilities that induce the theoretical band structure. In our work, we will consider breathers under circumstances where they are far from classical instability, so that quantum processes that connect to distinct classical breather states (e.g., the breather translated by one atom) would be of extremely small amplitude not (a). So we are saying that *except for this small tunneling amplitude* the breather is stable. Physically this means that the breather can have an extremely long lifetime, contrary to the claims made in the literature that we have cited Hizhnyakov (1996, 1999); Hizhnyakov and Nevedrov (1997); Hizhnyakov et al. (2002). We emphasize that in that literature translational invariance is also dropped.
A second consideration is that the physical systems of particular interest to us Schulman et al. (2002) do *not* have translational invariance. The breather formed by Jahn-Teller distortion at an impurity certainly has a preferred origin, and it is stability against low order phonon decay that concerns us.
## IV The Path integral approach
Feynman’s use of the path integral for the polaron Feynman (1955); Feynman and Hibbs (1965) exploited what has come to be one of the most useful features of the path integral: the propagator for a quadratic degree of freedom coupled linearly to something else can be evaluated explicitly, leaving only a self-coupling of that “something else.” The penalty is the nonlocality in time of the self coupling, the reward the reduction of the problem to a single degree of freedom.
### IV.1 Setting up the path integral
As discussed in Sec. III, for appropriate parameters, all atoms but one have negligible nonlinear energy contributions, and we drop the “$`\lambda x^4`$” for all but one of them. This allows us to integrate all but the single atom, around which the breather is centered. As in Eq. (3), we take that atom to be #0. As a result we consider the Lagrangian
$$=\frac{1}{2}\underset{n=1}{\overset{N+1}{}}\dot{x}_n^2\frac{1}{2}\omega _0^2\underset{n=1}{\overset{N+1}{}}(x_nx_{n+1})^2\frac{1}{2}\omega _s^2\underset{n=1}{\overset{N+1}{}}x_n^2\frac{\lambda }{4}x_0^4+\mu x_0(x_m+x_{N+1m}).$$
(10)
We have introduced an additional term, $`\mu x_0(x_m+x_{N+1m})`$. It is easy to carry along and will later allow us to study correlations. The derivative (at $`\mu =0`$) of the (appropriate form of the) propagator provides a (0-$`m`$)-correlation. We will take particle-$`m`$ to be distant from 0, typically about $`1/3`$ of the way around the ring. The use of a pair \[$`m`$ and $`(N+1m)`$\] maintains mirror symmetry (“$`P`$”) with respect to particle-0. We will also distinguish between the *ring*, all $`N+1`$ atoms, and the *chain*, atoms 1 through $`N`$. See Fig. 7.
The Lagrangian $``$ can be written in the form $`=_0+_C+_I`$. Here $`_0`$ is the Lagrangian of particle-0
$$_0=\frac{1}{2}\dot{q}^2\left(\omega _0^2+\frac{1}{2}\omega _s^2\right)q^2\frac{\lambda }{4}q^4,$$
(11)
where we have adopted the notation $`q`$ for $`x_0`$. $`_C`$ is the Lagrangian of the chain
$$_C=\frac{1}{2}\underset{n=1}{\overset{N}{}}\dot{x}_n^2\frac{1}{2}\omega _0^2\underset{m,n=1}{\overset{N}{}}J_{mn}x_mx_n\frac{1}{2}\omega _s^2\underset{n=1}{\overset{N}{}}x_n^2,$$
(12)
with $`J_{mn}=2\delta _{mn}\delta _{m+1,n}\delta _{m,n+1}`$ ($`m,n=1\mathrm{}N`$), the Jacoby matrix. The last term $`_I`$ is the interaction between the particle-0 and the others
$$_I=\omega _0^2q\left(x_1+x_N\right)+\mu q(x_m+x_{N+1m}).$$
(13)
We next perform the path integral for the chain and arrive at an effective action and path integral for $`q`$ alone Feynman and Hibbs (1965); Feynman (1972); Schulman (1981); Weiss (1999). Normal coordinates for the chain are defined by $`𝐐=\mathrm{𝐒𝐱}`$, where $`S_{mn}=\sqrt{2/\left(N+1\right)}\mathrm{sin}\left[\pi mn/\left(N+1\right)\right]`$ is a unitary matrix, and $`(x_1,\mathrm{},x_N)`$ is considered a column vector. This leads to
$$_C+_I=\frac{1}{2}\underset{n=1}{\overset{N}{}}\left\{Q_n^2\mathrm{\Omega }_n^2Q_n^2+f_n(t,q)Q_n\right\}.$$
(14)
Here $`\{\mathrm{\Omega }_n\}`$ are the *chain* normal-mode frequencies
$$\mathrm{\Omega }_n^2=\omega _s^2+4\omega _0^2\mathrm{sin}^2\frac{\pi n}{2\left(N+1\right)},$$
(15)
and $`f_n(t,q)=q\tau _n\left[1\left(1\right)^n\right]`$, where
$$\tau _n\sqrt{\frac{2}{N+1}}\left(\omega _0^2\mathrm{sin}\frac{\pi n}{N+1}+\mu \mathrm{sin}\frac{\pi nm}{N+1}\right).$$
(16)
The path integrals over the normal chain coordinates $`Q_n`$ can easily be evaluated, so that one has a propagator that is a function of the endpoints of *all* degrees of freedom. It is convenient to take the matrix element of this object in the chain ground state and divide by corresponding matrix elements of the unforced chain oscillators. This yields an effective action for the chain
$$e^{\frac{i}{\mathrm{}}S_{\text{eff}}[q]}=\left[\underset{n=1}{\overset{N}{}}𝒟Q_n\psi _0^f|e^{\frac{i}{\mathrm{}}(S_C+S_I)}|\psi _0^i\right]/\underset{n=1}{\overset{N}{}}𝒟Q_n\psi _0^f|e^{\frac{i}{\mathrm{}}S_C}|\psi _0^i,$$
(17)
where $`S_\gamma =_\gamma `$, $`\gamma =C,I`$, and (e.g.) $`\psi _0^f`$ is shorthand for $`_n\psi _0^{(\mathrm{\Omega }_n)}(Q_n^f)`$, and $`\psi _0^{(\omega )}(Q)=\left(\omega /\pi \mathrm{}\right)^{1/4}\mathrm{exp}(\omega Q^2/2\mathrm{})`$ is the ground state of a harmonic oscillator. The chain oscillator modes are subject to a forcing term due to the so-far-unintegrated $`q(t)`$.
The result of these standard calculations is that the propagator for particle-0 takes the form
$$𝒢(q^f,T;q^i,0)=𝒟qe^{\frac{i}{\mathrm{}}\left(S_0+S_{\text{eff}}\right)},$$
(18)
where the action $`S_0`$ arises from the original $``$ sans chain terms, specifically
$$S_0=𝑑t\left\{\frac{1}{2}\dot{q}^2\left(\omega _0^2+\frac{1}{2}\omega _s^2\right)q^2\frac{\lambda }{4}q^4\right\}.$$
(19)
The effective action, $`S_{\text{eff}}`$, is the result of the integration just described over chain degrees of freedom. It is given by
$$S_{\text{eff}}=_0^T_0^T𝑑t𝑑sK(|ts|)q(t)q(s),$$
(20)
with
$$K(u)=\underset{n=1,3,\mathrm{}}{\overset{N}{}}\tau _n^2\frac{\mathrm{cos}\mathrm{\Omega }_n\left(\frac{T}{2}u\right)}{\mathrm{\Omega }_n\mathrm{sin}\left(\frac{\mathrm{\Omega }_nT}{2}\right)}.$$
(21)
We shall study the (matrix element of the) propagator in the stationary phase approximation; that is, we focus on the action (including $`S_{\text{eff}}`$) along the extremal “classical paths.” The spectral expansion of the propagator \[as a function of *all* variables, i.e., *before* the chain-variable integrations and the division of Eq. (17)\] is
$$G(x_0^{\prime \prime },x_1^{\prime \prime },\mathrm{},t;x_0^{},x_1^{},\mathrm{})=\underset{\alpha }{}\mathrm{\Phi }_\alpha (x_0^{\prime \prime },x_1^{\prime \prime },\mathrm{})\mathrm{exp}(iE_\alpha t)\mathrm{\Phi }_\alpha ^{}(x_0^{},x_1^{},\mathrm{}),$$
(22)
where $`\alpha `$ is an eigenstate label. This implies that the operations of Eq. (17) (the integrations, *and* the division, which leads to energy-phase corrections) lead to
$$𝒢(q^{\prime \prime },t;q^{},0)=\mathrm{exp}\left(+i\mathrm{\Omega }_nt/2\right)\underset{\alpha }{}\varphi _\alpha (q^{\prime \prime })\mathrm{exp}(iE_\alpha t)\varphi _\alpha ^{}(q^{}),$$
(23)
where
$$\varphi _\alpha (q)\mathrm{\Psi }_0(x_1,\mathrm{})^{}\mathrm{\Phi }_\alpha (q,x_1,\mathrm{})𝑑x_1\mathrm{},$$
(24)
and $`\mathrm{\Psi }_0`$ is the ground state of the chain. Note that $`E_\alpha `$ in Eq. (23) is the total energy of the ring, while $`E_C\mathrm{\Omega }_n/2`$ is the total ground state energy of the chain. Going to imaginary time, $`t=iT`$, Eq. (23) becomes
$$\stackrel{~}{𝒢}(q^{\prime \prime },T;q^{},0)=\underset{\alpha }{}\varphi _\alpha (q^{\prime \prime })\mathrm{exp}\left[(E_\alpha E_C)T\right]\varphi _\alpha (q^{})^{}$$
(25)
This is the starting point for many calculations. The simplest would be to get the ground state energy of the chain by letting $`T\mathrm{}`$. For numerical work this turns out to be delicate. What one would do is find a semiclassical expression for the total action in Eq. (18) and consider its large $`T`$ limit. Unfortunately $`S`$ tends to a constant, and the energy emerges from its derivatives. The same occurs for the simple harmonic oscillator. Thus for $`_{\text{SHO}}=m\dot{x}^2/2m\nu ^2x^2/2`$, the propagator from $`a`$ to $`b`$ in imaginary time $`T`$ is
$`\stackrel{~}{G}_{\text{SHO}}(b,T;a)`$ $`=`$ $`\sqrt{{\displaystyle \frac{1}{2\pi }}{\displaystyle \frac{^2S}{ab}}}\mathrm{exp}\left(S_{\text{classical}}\right)`$ (26)
$`=`$ $`\sqrt{{\displaystyle \frac{m\nu }{2\pi \mathrm{sinh}\nu T}}}\mathrm{exp}\left({\displaystyle \frac{m\nu }{2\mathrm{sinh}\nu T}}\left[(a^2+b^2)\mathrm{cosh}\nu T2ab\right]\right)`$
For large $`T`$ the dominant term in the exponential approaches a constant, so that the energy comes from the prefactor, $`\sqrt{^2S/ab}`$, which as the explicit form above shows is essentially $`1/\sqrt{\mathrm{sinh}\nu T}`$, yielding the correct ground state energy, $`\nu /2`$. This means that one needs extremely accurate values in a second derivative of the action, exponentially small in $`T`$. For the action itself this precision was attainable in all cases, and for a single harmonic oscillator the second derivative was also reliable, but for the full ring we did not find that the second derivative had the needed precision.
In any case, evaluating the energy is not our primary goal. We know that the Hamiltonian has eigenstates, but we do not know whether they resemble breathers. We also expect that some states, not necessarily eigenstates, but at the least metastable states, should resemble the classical breather state. This much is guaranteed by semiclassical considerations. Our goal then is to show that the actual eigenstates resemble the breathers.
To see how to do this, we will consider first the local mode, a quantum state that we know to be localized. Specifically, the classical mode function has a sharp spatial dropoff and the quantum state represents atomic oscillations with that pattern. For reference we show the classical local mode for $`\omega _1=2`$ and $`\omega _0=1`$ in Fig. 8 \[a solution of Eq. (5)\]. (The value of $`\omega _s`$ does not affect the function.)
Our goal is to ascertain properties of $`\varphi _\alpha `$ and $`E_\alpha `$ from the semiclassical approximation to $`G`$. This goal will be attained in Sec. IV.3. However, before actually doing this calculation we discuss methodology. Both for the linear and nonlinear cases, implementing the semiclassical approximation presents technical challenges, which we now discuss. We will also indicate the extent to which independent confirmations of our technique are available, taking advantage of known asymptotic properties.
### IV.2 Methodology
The first observation to be made is that the development from Eq. (10) to Eq. (25) goes through in exactly the same manner, with the same $`K`$, etc., for the quadratic local mode as for the nonlinear excitation, with the sole exception of the final form of $`S_0`$, which, for the local mode case, has a term $`\omega _1^2q^2/2`$ instead of $`\lambda q^4/4`$. Going over to imaginary time, our problem (both quantum and classical) revolves about the following quantities:
$$S_{\text{total}}=_0^T𝑑t\left\{\frac{1}{2}\dot{q}^2+\frac{1}{2}(\omega _s^2+2\omega _0^2)q^2+U(q)\right\}_0^T_0^T𝑑t𝑑s\stackrel{~}{K}(|ts|)q(t)q(s),$$
(27)
where $`U(q)=\omega _1^2q^2/2`$ for the local mode problem. The same form holds for the nonlinear problem as well, but with $`U(q)=\lambda q^4/4`$. $`S_{\text{total}}`$ is the total classical action, and in the path integral appears as $`\mathrm{exp}(S_{\text{total}})`$ (note the minus sign, a result of the $`tit`$ transformation). The self-interaction kernel becomes
$$\stackrel{~}{K}(u)=\underset{n=1,3,\mathrm{}}{\overset{N}{}}\tau _n^2\frac{\mathrm{cosh}\mathrm{\Omega }_n\left(\frac{T}{2}u\right)}{\mathrm{\Omega }_n\mathrm{sinh}\left(\frac{\mathrm{\Omega }_nT}{2}\right)}$$
(28)
The classical equation of motion follows from the usual variational methods and is
$$\ddot{q}(\omega _s^2+2\omega _0^2)q\frac{U}{q}+2_0^T\stackrel{~}{K}(|ts|)q(s)=0.$$
(29)
For the semiclassical approximation one must solve Eq. (29). We used different methods for the linear and nonlinear cases. We discuss them separately.
#### IV.2.1 Linear, non-local in time propagators
For $`U(q)=\omega _1^2q^2/2`$, Eq. (29) can be discretized, say with $`t_k=kϵ`$, $`k=0,1,\mathrm{},M+1`$, $`ϵ=T/M`$. $`q`$ becomes a vector, $`q_k=q(t_k)`$ ($`k=1,\mathrm{},M`$) and the entire equation has the form $`Bq=q_0`$, with $`B`$ a matrix consisting of three parts: the second derivative operator, $`2/ϵ^2`$ on the diagonal, $`1/ϵ^2`$ above and below; a diagonal matrix proportional to $`(\omega _s^2+2\omega _0^2)`$ and a matrix (which by a slight abuse of notation we call) $`K`$ whose substantial non-diagonal components reflect the nonlocality in time. From Eq. (29) it may not be evident how a nonzero right-hand-side ($`q_0`$) comes into the picture. It arises from the boundary conditions. When discretizing, the first and last rows of the second derivative operator call on $`q`$-components outside the range $`1,\mathrm{},M`$. These other components are the boundary values (call them $`a=q_{\text{initial}}`$ and $`b=q_{\text{final}}`$), so that $`q_0(1)=a/ϵ^2`$ and $`q_0(M)=b/ϵ^2`$ (other components of $`q_0`$ are zero). This method of dealing with the two-time boundary value problem is a simpler version of what is done in Schulman (1974). It follows that the solution of the nonlocal-in-time equations of motion is $`q=B^1q_0`$. It is also immediate that for the linear equation the action along this “classical path” is given by $`S=[q(t)\dot{q}(t)/2]|_0^T`$, although one can also evaluate $`S`$ by explicit integration.
The substantial nonlocality in time means that the matrix $`K`$ is not sparse, putting a limit on the smallness of $`ϵ`$, hence on the accuracy of the action. This difficulty was overcome by the following device. For each such $`ϵ`$, call $`S(ϵ)`$ the action that results from the discretization and matrix inversion just described. Considered as a function of $`ϵ`$, one can write $`S(ϵ)=S(0)+ϵS^{}(0)+ϵ^2S^{\prime \prime }(0)/2+\mathrm{}`$, where the prime indicates derivative with respect to $`ϵ`$. By evaluating $`S(ϵ)`$ for several values of $`ϵ`$ one can then extrapolate to zero. For example, the simplest such extrapolation, for two values, $`ϵ_1`$ and $`ϵ_2`$ (and associated $`S_1`$ and $`S_2`$), gives $`S(0)=(ϵ_2S_1ϵ_1S_2)/(ϵ_2ϵ_1)`$. We typically used three values, chosen as small as possible, but far enough apart to keep denominators from being too small. As we indicate below, this gave excellent results when tested against quantities that could be independently calculated. The following tests were performed:
* Simple harmonic oscillator. For the energy this requires taking the second derivative with respect to position. This was done by evaluating $`G`$, hence $`S`$, at several points and using
$$\frac{^2S}{ab}|_{a=b}\frac{S(a+\delta ,a+\delta )+S(a\delta ,a\delta )S(a+\delta ,a\delta )S(a\delta ,a+\delta )}{(2\delta ^2)}$$
(30)
For this formula the extrapolation to $`\delta =0`$ was also used (as for $`ϵ`$ above), which is to say, several values of $`\delta `$ were used. Furthermore, there was also an extrapolation to $`T=\mathrm{}`$ (by fitting the logarithm of $`^2S/ab`$ as a function of $`1/T`$). Results were accurate to better than one part in 10<sup>4</sup>.
* Asymptotics of $`\varphi _0(q)`$ \[defined in Eq. (24)\]. This can be calculated by non-path integral methods as follows. All eigenstates—both ring and chain—are Gaussians. For the chain the frequencies are $`\mathrm{\Omega }_n^{(\text{chain})}=\sqrt{\omega _s^2+4\omega _0^2\mathrm{sin}^2(n\pi /2(N+1))}`$, $`n=1,\mathrm{},N`$. For the ring, for the case $`\omega _1=0`$, they are $`\mathrm{\Omega }_n^{(\text{ring})}=\sqrt{\omega _s^2+4\omega _0^2\mathrm{sin}^2(n\pi /(N+1))}`$, $`n=0,\mathrm{},N`$. If $`\omega _10`$, they are easily evaluated numerically. Aside from normalization, $`\varphi _0`$ is therefore given by \[cf. Eq. (24)\]
$$\varphi _0(q)=\underset{n=1}{\overset{N}{}}dQ_n^{(\text{chain})}\mathrm{exp}\left(_k\mathrm{\Omega }_k^{(\text{chain})}Q_k^{(\text{chain})2}/2\right)\mathrm{exp}\left(_{\mathrm{}}\mathrm{\Omega }_{\mathrm{}}^{(\text{ring})}Q_{\mathrm{}}^{(\text{ring})2}/2\right).$$
(31)
Given the transformations from the original variables, $`(q,x_1,\mathrm{},x_N)`$ to the chain and ring variables (known analytically for the chain, numerically for the ring), it is straightforward to get the coefficient of $`q^2`$ in $`\varphi _0`$. Let the matrices $`U^{(\text{ring})}`$ and $`U^{(\text{chain})}`$ be defined by $`Q_k^{(\text{ring})}=_0^NU_k\mathrm{}^{(\text{ring})}x_{\mathrm{}}`$ and $`Q_k^{(\text{chain})}=_1^NU_k\mathrm{}^{(\text{chain})}x_{\mathrm{}}`$ (with $`\{Q\}`$ the appropriate normal coordinates). The $`U`$’s can be taken to be real. Let
$`u_k`$ $``$ $`U_{k0}^{(\text{ring})},\mathrm{\Omega }_u{\displaystyle \underset{0}{\overset{N}{}}}\mathrm{\Omega }_k^{(\text{ring})}u_k^2,`$ (32)
$`g_n`$ $``$ $`{\displaystyle \underset{k=0}{\overset{N}{}}}{\displaystyle \underset{\mathrm{}=1}{\overset{N}{}}}\mathrm{\Omega }_k^{(\text{ring})}u_kU_k\mathrm{}^{(\text{ring})}U_\mathrm{}n^{(\text{chain})},`$ (33)
$`W_{nm}`$ $``$ $`\delta _{nm}\mathrm{\Omega }_n^{(\text{chain})}+{\displaystyle \underset{k=0}{\overset{N}{}}}{\displaystyle \underset{\mathrm{}=1}{\overset{N}{}}}{\displaystyle \underset{\mathrm{}^{}=1}{\overset{N}{}}}U_n\mathrm{}^{(\text{chain})}U_\mathrm{}k^{(\text{ring})}Q_k^{(\text{ring})}U_k\mathrm{}^{}^{(\text{ring})}U_\mathrm{}^{}m^{(\text{chain})}.`$ (34)
(Note that $`U^{(\text{chain})}=U^{(\text{chain})}`$.) Then the coefficient of $`q^2`$ in the exponent of $`\varphi _0^2`$ (call it $`\zeta `$) is
$$\zeta =\mathrm{\Omega }_ug^{}W^1g.$$
(35)
To evaluate $`\zeta `$ by the path integral method, the action was evaluated with various endpoints, $`a`$. Since $`S(a,T;a)`$ (the action from $`a`$ to $`a`$ in time $`T`$) appears in the exponent of the propagator, if the semiclassical method and our extrapolations are valid, $`S`$ should be quadratic in $`a`$. This was tested by fitting $`\sqrt{S}`$ against $`a`$. We did not do this for a variety of $`T`$ values since a single large $`T`$ was sufficient. Large in this case means with respect to the smallest spectral gap, and can be discerned numerically by a stabilizing of the slope in the fit. The desired result is that first, the fit of $`\sqrt{S}`$ against $`a`$ should be a straight line, and second, that the square of the slope of the line should agree with $`\zeta `$. We will not reproduce the graph of $`\sqrt{S}`$ against $`a`$ since there is little to note beyond its being an excellent fit. The values of slope-squared and $`\zeta `$ agreed to better than 1% for $`\mathrm{min}(\omega _1,\omega _s)1`$ and by no more than 3% when these quantities became smaller (they impact the spectral gap, which can force $`T`$ values too large for good accuracy).
#### IV.2.2 Non-linear, non-local in time propagators
With a nonlinear term in the action, the classical nonlocal two-time boundary value problem cannot be solved by matrix inversion. The equation to be solved is
$$\ddot{q}(\omega _s^2+2\omega _0^2)q\lambda q^3+2_0^T\stackrel{~}{K}(|ts|)q(s)=0.$$
(36)
Our method is an extension of Feynman’s approach Feynman and Hibbs (1965) in which the introduction of an auxiliary variable eliminates the nonlocality. Then one can use standard numerical methods for solving the two-time (even-if-nonlinear) boundary value problem. Recalling the definition of $`\stackrel{~}{K}`$, Eq. (28), one can define auxiliary variables $`z_n`$ by
$$z_n(t)_0^T𝑑s\frac{\mathrm{cosh}\mathrm{\Omega }_n\left(\frac{T}{2}|ts|\right)}{\mathrm{\Omega }_n\mathrm{sinh}\left(\frac{\mathrm{\Omega }_nT}{2}\right)}q(s).$$
(37)
Eq. (36) becomes
$$\ddot{q}(\omega _s^2+2\omega _0^2)q\lambda q^3+2\underset{n}{}\tau _n^2z_n=0.$$
(38)
Taking two derivatives of $`z_n`$ leads to
$$\frac{d^2z_n}{dt^2}=\mathrm{\Omega }_n^2z_n2q(t),n=1,3,\mathrm{},2\left[\frac{N1}{2}\right]+1.$$
(39)
One non-local equation has been replaced by a larger number of local ones. There is but a single pair of boundary conditions: $`q(0)=a`$, $`q(T)=b`$. The conditions on $`z`$ are forced by self-consistency arising from Eq. (37), which incidentally also imply $`z_n(0)=z_n(T)`$, for all $`n`$. It’s amusing that with the quartic interaction replaced by a quadratic one, this linear system is equivalent to the set of classical linear equations that apply to the ring not (b), with the boundary conditions inherited from the ground state averaging \[which is of course where Eq. (36) (made linear) comes from in the first place\].
For numerical solution of Eq. (36) we applied a variation of this method. As a function of its argument, $`\stackrel{~}{K}(u)`$, although a sum of many hyperbolic cosines, actually (for large $`N`$ and $`T`$) bears a strong resemblance to a single such function. For the parameter range of interest, one can, by judicious choice of $`\omega _{\text{eff}}`$, bring the difference between $`\stackrel{~}{K}`$ and its approximation, $`\stackrel{~}{K}_1(u)\stackrel{~}{K}(0)\mathrm{cosh}(\omega _{\text{eff}}u)`$, to about $`10^4`$ (integral of square of difference). With a sum of two hyperbolic cosines (and 3 adjustable parameters) one can do much better, but as we shall see below a single hyperbolic cosine was sufficiently accurate for excellent minimization of the action, i.e., solution of the classical motion.
With a single hyperbolic cosine in the nonlocal portion of the equation, only a single $`z`$ need be defined \[cf. the derivation of Eq. (39)\] so that our nonlocal equation becomes a pair of second order ODE’s, with boundary values for one of them and a self-consistency condition for the boundaries of the other. Using the optimal $`\omega _{\text{eff}}`$ defined in the last paragraph, the equations are
$`\ddot{q}`$ $`=`$ $`(\omega _s^2+2\omega _0^2)q+\lambda q^32z,`$ (40)
$`\ddot{z}`$ $`=`$ $`\omega _{\text{eff}}^2z2q(t),\text{with }z(t){\displaystyle _0^T}𝑑s{\displaystyle \frac{\mathrm{cosh}\omega _{\text{eff}}\left(\frac{T}{2}|ts|\right)}{\omega _{\text{eff}}\mathrm{sinh}\left(\frac{\omega _{\text{eff}}T}{2}\right)}}q(s).`$ (41)
The boundary values of $`q(t)`$ are given. For $`z`$ one proceeds iteratively. For given $`z(0)`$ (which automatically equals its value at $`T`$), the boundary value problem can solved using MATLAB’s program “bvp4c” mat . Using the associated solution, $`q(t)`$, one can recompute $`z(0)`$. In effect one has a function mapping $`z(0)`$ into a new value, and again, standard numerical search techniques can be use to find a fixed point of this map. That fixed point then provides a solution of the original boundary value problem for $`q`$.
In principle with a better approximation for $`\stackrel{~}{K}`$, using say two hyperbolic cosines, the parameter space of boundary value mappings becomes 2-dimensional. This, however, was not the way we proceeded.
Since the true solution of the classical nonlinear, nonlocal problem minimizes the action, $`S`$, it is possible to improve the solution by modifying the $`q`$ derived from the process given above in such a way as to reduce $`S`$. The class of functions to add to $`q(t)`$ for this purpose can be narrowed by the following consideration. The potential that we study is an inverted oscillator, linear or nonlinear. Therefore if the time interval for going between two not-too-small boundary values is large, for most of that time interval the particle will be near zero, with zero velocity: the path is thus an instanton and, except near the endpoints, will be exponentially small (in $`T`$). Profitable, i.e., $`S`$-reducing, variations of $`q`$ will thus have the same shape. Our basic variation consists of adding and subtracting hyperbolic cosine functions with varying amplitudes and angular frequencies.
Thus, starting with the approximate-$`\stackrel{~}{K}`$ solution to the nonlinear boundary value problem, we allowed modifications of the sort just described. The changes in $`S`$ that resulted were extremely small. The first such correction was on the order of 10<sup>-5</sup> of the action, and subsequent reductions were of order 10<sup>-15</sup>. No other functional forms (e.g., multiplying instanton-shaped curves by oscillatory functions) gave any improvement at all. The results presented below all use this method of optimizing $`S`$.
It is also possible to perform some of the checks made on the linear problem. In particular the asymptotic form of the wave function for a quartic anharmonic oscillator is $`\psi \mathrm{exp}\left([\text{positive const}q^3]\right)`$. With various $`\lambda `$ we found the action for large boundary values of $`q`$. There was a very good fit to the cubic. In Fig. 9 we plot various powers of the rescaled action, showing that the dropoff for what we have called $`\varphi _0(q)`$ is close to cubic (the exact wave function is also not exactly cubic).
### IV.3 Localization
The local mode represents a localized excitation. Fig. 8 shows the classical mode oscillation amplitude, hence the shape of the phonon; in particular it indicates an exponential dropoff with distance from the center of the mode. We will use our path integral formalism to show how it too reflects the localization of the “local mode.” Then we will use the same technique to establish the localization of the *nonlinear* breather wave function.
Recall the fictitious coupling, $`\mu x_0(x_m+x_{N+1m})`$, inserted in the Lagrangian of Eq. (10). Site-$`m`$ is far from site-0, where the large oscillations of the breather are taking place. Comparing the action for *small* $`\mu `$ and the action for *zero* $`\mu `$ provides a correlation function in the following way. The imaginary-time version of Eq. (18) is
$$𝒢(q^f,T;q^i,0)=𝒟qe^{S_{\text{total}}/\mathrm{}},$$
(42)
where there are slight differences from the definitions given in Eqs. (19) through (21). Specifically, the action ($`S_{\text{total}}`$) is now given by Eq. (27), with $`\stackrel{~}{K}`$ as defined there \[note that Eq. (27) includes both the cases of quadratic and quartic potentials\]. We emphasize that the fictitious coupling appears in the definition of $`\tau _n`$ \[Eq. (16)\], and *only* there.
Thinking back to the propagator before the integration over chain mode ground states, we consider the derivative $`G(a,iT;a)/\mu |_{\mu =0}`$. The importance of this derivative arises from the relation
$$A=\frac{}{\mu }\left[\mathrm{exp}(S+\mu A)𝑑x\right]|_{\mu =0}/\mathrm{exp}(S)𝑑x,$$
(43)
where $`\mathrm{exp}(S)`$ is a weight for averaging. Since we are taking $`/\mu `$ after integrating, we are getting an average of the $`qx_m`$-correlation in the chain mode ground state not (c).
Note too that because of our use of the semiclassical approximation (which happens to be exact in the linear case), study of $`G(a,iT;a)/\mu |_{\mu =0}`$ is essentially the same as study of $`\sigma (a,T,g)S(a,T)/\mu |_{\mu =0}`$, where “$`S`$” is the imaginary time action, $`a`$ is the common initial and final endpoint, and $`T`$ the imaginary time. The quantity $`g`$, implicit in $`S(a,T)`$, parameterizes the relevant particle-0 enhancement; for the linear local mode it is $`\omega _1`$ and for the quartic case it is $`\lambda `$.
Consider first $`\sigma (a,T,\omega _1)S(a,T)/\mu |_{\mu =0}`$ for moderate $`T`$ (such that states other than the ground state survive in the spectral sum for $`𝒢`$), for small $`\omega _1`$ and for large $`a`$. Because $`a`$ is large, the important terms in the spectral sum will not be those of lowest energy, but those that permit large excursions of $`q`$ \[cf. Eqs. (23) and (24)\]. However, since $`\omega _1`$ is assumed small, there will be no local mode. So pulling $`q`$ far from its equilibrium position, pulls *all* atoms far from their positions, and the correlation of atom-0 ($`q`$) and atom-$`m`$ ($`x_m`$) should be large. On the other hand, suppose $`\omega _1`$ to be large (with $`a`$ still large and $`T`$ moderate). In that case there is a pronounced local mode and pulling $`q`$ away from equilibrium has little impact on $`x_m`$. For the local mode, $`q`$ can have large excursions while other atoms hardly move. Therefore one expects $`\sigma `$ to be small. In Fig. 10a, the lowest curve shows just this behavior. The boundary value of $`q`$, $`a`$, is 4 and for small $`\omega _1`$ (no local mode) there is a large correlation with the motion of atom-$`m`$ (0 is at the *top* of the figure). As $`\omega _1`$ increases, this correlation shrinks.
By contrast, if $`a`$ is small, the variation of $`\omega _1`$ has little effect. This can be understood as follows. The requirement on the endpoints of $`q`$ imposes little demand on any other coordinates whether or not there is a local mode.
We next turn to the nonlinear case. Here we do not have a priori knowledge of the wave function but can use the correlation function as a test of localization. The behavior of the correlation function, as a function of $`a`$ and $`\lambda `$ exactly parallels that of the linear local mode. This is shown in Fig. 10b. For small $`\lambda `$ and large $`a`$ there is no breather and, as for the small $`\omega _1`$ case, the demand for a large excursion of $`q`$ forces a large excursion of $`x_m`$. And now the central observation: for large $`\lambda `$, forcing $`q`$ to be large has almost no impact on $`x_m`$, exactly as for the quadratic local mode, from which we deduce the localization of the breather excitation. We remark that corresponding values of $`\omega _1`$ and $`\lambda `$ are related by $`\lambda \omega _1^2`$. Finally we mention that from Fig. 10b alone it is difficult to tell whether $`dS/d\mu `$ is tending to a constant or to zero. In Fig. 11 is a plot of $`(dS/d\mu )^2`$ versus $`1/\lambda `$, in which the extrapolated value is close to zero.
Another check of correlation function behavior is its time dependence. For a localized stable state and fixed $`a`$ (the boundary value for $`q`$) $`dS/d\mu `$ should decrease with time. This follows from consideration of the propagator given in either Eq. (22) or Eq. (23). For large $`T`$ and $`a`$ there is a competition between the terms in the expansion, with large $`T`$ favoring states of lower energy and large $`a`$ favoring states that drop off most slowly in their spatial coordinate, which in general will be of higher energy. As $`T`$ increases lower energy states are increasingly important in the mix and the correlation correspondingly reduced (note that for the correlation the overall magnitude of the propagator, and in particular the factor $`\mathrm{exp}(E_0T)`$, drops out). In Fig. 12 this decline can be seen for both the quadratic local mode and the quartic breather. On the other hand, if the quantum breather were decaying in time the correlation should *increase*, since the true ground state of the system (if the breather were not stable) would resemble the $`\lambda =0`$ case, for which the correlation is high.
## V Direct diagonalization
The full Hamiltonian governing our system is given by Eq. (3) \[or Eq. (4)\] and it is natural to seek the properties of the system by numerical diagonalization of $`H`$ in an appropriate basis. Indeed this is the approach of Wang et al. (1996). The difficulty lies in the fact that a cutoff in phonon number is necessary in order to keep the problem finite. Since the overall Hilbert space is a product of individual phonon Hilbert spaces the total dimension is the product of that of the individual subspaces, so that the size of the matrix representing $`H`$ is large. Moreover, even when one has the eigenstates of $`H`$ (which by definition are stable), one still needs to show them to be localized in order to identify them with breathers. (Evidence for localization is also given in Wang et al. (1996).)
To reduce the burden of dimensional proliferation we used two devices. First, we introduce a second system for comparison, namely the Hamiltonian Eq. (5) \[or Eq. (6)\], representing a local mode. The idea is that if $`\omega _1`$ (the fictitious self-coupling at site-0) is chosen appropriately the frequency and to some extent the shape of the breather (induced by a quartic with coefficient $`\lambda `$) can be reasonably approximated. In that way, the perturbation, has minimal effect, and only a small number of non-localized phonon states enter the eigenstate associated with the breather. See Fig. 13 for an illustration of the efficacy of this method. The second device takes advantage of the reflection symmetry embodied in $`P`$ \[Eq. (9)\]. Since the perturbation ($`\lambda x_0^4/4\omega _1^2x_0^2/2`$) commutes with $`P`$ we can focus on states in the same symmetry class as the local mode, namely those with $`P`$-eigenvalue 1. This reduces the number of phonons by almost a factor two, lowering the Hilbert space dimension to a bit more than the square root of what we would otherwise need to consider. Note that using eigenstates of $`P`$ means that we are not using “traveling waves” for the phonons, and indeed one should no longer expect these to be the preferred basis once translational invariance has been dropped.
Some details of the calculation play a role in interpreting the results and we present them here. The normal modes of our monatomic ring \[Eq. (5)\] satisfy the classical equation of motion,
$$\ddot{x}_k+\omega _s^2x_k+\omega _0^2\left(2x_kx_{k1}x_{k+1}\right)+\delta _{k0}\omega _1^2x_k=0,k=0,\mathrm{},N+1,$$
(44)
with mod-$`(N+1)`$ addition and $`\delta _{jk}`$ the Kronecker delta. For $`\omega _1=0`$, this is trivially solvable, while for non-zero $`\omega _1`$ there is a local mode. Corresponding equations hold for the nearest-neighbor nonlinear-coupling model \[Eq. (6)\], but since the principles are the same we do not present the equations in detail. With $`\omega _0=1`$ and $`\omega _s`$ of order unity, a typical spectrum is close to that of the $`\omega _1=0`$ case, except for a single mode with frequency high above all the others. An example is shown in Fig. 14.
At the quantum level the Hamiltonian is written using the classical modes. Call $`\mathrm{\Omega }_{\mathrm{}}`$ the frequency of the $`\mathrm{}^{\text{th}}`$ mode, $`u^{(\mathrm{})}`$ the coordinates of the mode \[i.e., $`u^{(\mathrm{})}\mathrm{exp}(i\mathrm{\Omega }_{\mathrm{}}t)`$ solves Eq. (44)\], and let $`a_{\mathrm{}}`$ be its annihilation operator. The frequencies, $`\mathrm{\Omega }_{\mathrm{}}`$, are not the same as those defined in Eq. (15); those are for the chain, these for the ring. (When they appear together we distinguish using superscripts \[cf. Eq. (31)\]). We label the local mode “0” (although it is of highest frequency). As usual, $`H_0`$ becomes
$$H_0=\mathrm{\Omega }_{\mathrm{}}a_{\mathrm{}}^{}a_{\mathrm{}},$$
(45)
where we now restrict ourselves to modes symmetric under $`P`$, and call $`N_s`$ the number of such modes. The perturbation, given earlier, is
$$V=\frac{1}{4}\lambda x_0^4\frac{1}{2}\omega _1^2x_0^2,$$
(46)
so that all we need in order to proceed is the fact that
$$x_0=\frac{1}{\sqrt{2\mathrm{\Omega }_{\mathrm{}}}}\left(a_{\mathrm{}}+a_{\mathrm{}}^{}\right)u_0^{(\mathrm{})}.$$
(47)
Note that it is the zeroth component of each mode function, $`u^{(\mathrm{})}`$, that is important. (Had the asymmetric states been included they would now drop out since their zeroth components vanish.) In our general discussion we indicated that choosing to perturb around the local mode lessens the impact of the quartic term. In Eq. (47) the specific mechanism of that effect can be seen: it is built into the function $`u_0^{(\mathrm{})}`$, which is plotted in Fig. 15. Note that without $`\omega _1`$ (i.e., no local mode) the values of $`u_0^{(\mathrm{})}`$ fluctuate around 0.2 ($`u`$ is square normalized), while with a moderate $`\omega _1`$ all but $`u_0^{(0)}`$ drop to much smaller values.
With these definitions and observations the operator $`H`$ can be generated. It is a sparse matrix and selected eigenvalues and eigenvectors can be obtained for quite large dimension.
Our objective is to find the eigenstates corresponding to breathers. The most convenient for our purposes is the first excited state, since it presumably does not have enough excitation for the cutoff to be sensed (cf. Fig. 13). As a state around which to perform the perturbation we use $`|0,\mathrm{},0,1`$, which is the eigenstate of the local mode Hamiltonian, $`H_0`$ of Eq. (5), having a single excitation of the local mode, with all other phonons in their ground state. The values of the parameters that we use are $`\lambda =8`$, $`\omega _0=1`$, and $`\omega _s=1`$, which after adjusting for differences of convention, are the values used in Wang et al. (1996). For $`\omega _1`$ we used 2.5, although the structure of the eigenfunction was not sensitive to this. In Table 1 we show the $`H_0`$-eigenstates with the largest components in the true eigenstate. Clearly the local mode dominates. The next largest component is in fact the thrice excited local mode, which does not represent a spreading of amplitude, but rather a readjustment of the shape of the excitation. Other modes barely make to the $`10^3`$ level. Of particular interest is the observation that the highest excitation level for other phonons is 3, indicating that a cutoff of 6 is safe. In fact, even to probability levels of $`10^8`$ there is no excitation higher than 3 for anything but the local mode.
As a check of cutoff sensitivity we repeated this calculation with a cutoff of 8 (but with the local mode still at 13). The results are in the last column of Table 1. There is little sensitivity to the change—not only in the probabilities but in the composition of the state.
We also studied larger rings, but with smaller cutoffs. For eight atoms there were 5 phonon modes, with a cutoff of 7 for the local-mode phonon and 5 for the others. Results with $`\omega _1=2.3`$ are shown in Table 2. Once again the principal change from the unperturbed state arises from the excited breather state. And again, down to $`10^6`$ probability no non-local-mode phonon has excitation greater than one, an indication that the cutoff is again not felt. It is interesting that in other runs with different values of $`\omega _1`$ most of the change in the state was in the local-mode-phonon’s contribution, which makes sense since it may be a better or worse approximation to the true nonlinear state, depending on $`\omega _1`$.
Moving on to 12 atoms, the number of independent phonons is 7 and a cutoff of 3 was imposed on all but the local mode, for which the cutoff was 5. The results, in Table 3, continue to be consistent with our previous conclusions.
Of particular interest for our own application is the case of nonlinearity in the nearest neighbor interaction. The results of this calculation support what we have already seen: the breather is very well approximated by phonons of the (quadratic) local mode. In Tables 4 and 5 we show a small variation on the material displayed for the self interaction. Both the singly excited breather (more precisely, the true eigenstate closest to the singly excited local mode) and the ground state are shown. A variety of $`\lambda `$ and $`\omega _1`$ values are also used, to give a richer idea of variation of the state with changes in parameters.
We summarize: in all cases the (quadratic) local mode dominates the true eigenstate of the Hamiltonian, indicating, because that (quadratic) local mode is itself localized, that the stable eigenstates, or at least those generated by perturbation around the localized modes, are localized. Matching all these properties, including frequencies, we conclude that the quantized breather is stable.
## VI Time evolution
### VI.1 General discussion of decay
Although we have shown the breather eigenfunction (of the full Hamiltonian) to be dominated by the local modes, nevertheless there were small—very small—contributions from ordinary phonon states (in the local mode basis). Does this imply that a system initially in a state close to the eigenfunction (for example, in a local mode phonon eigenstate), must ultimately decay? An argument in favor of this perspective would be that so long as the initial state has matrix elements that connect it to a continuum of global phonons, it must have a decay rate, simply by Fermi’s Golden Rule. And a decay rate implies exponential decay, even if the multiplier of time in that exponent is small.
As a general guide the Golden Rule has wide applicability, but it should not be made into a false idol (a Golden Calf). Several years ago, in response to anomalies in computer decay studies, one of the present authors in collaboration with B. Gaveau Gaveau and Schulman (1995) studied one kind of breakdown. The context is a state that has energy $`h`$ and is coupled to a continuum with energies in the interval $`[E_1,E_2]`$, with $`E_1<h<E_2\mathrm{}`$. The level with energy $`h`$ is within the interval and has nonzero coupling under the full Hamiltonian to the other levels. It was found that depending on the threshold behavior of the coupling, its strength, and the proximity of $`h`$ to an edge, the system need not decay exponentially. What can happen is that when the interaction is turned on the initial state loses a bit of amplitude to other modes, but it eventually settles into an asymptotic state of norm well away from zero. To be precise, let the initial state (of energy $`h`$) be $`|0`$ and the full Hamiltonian be $`H`$. The survival probability is then $`p(t)|0|\mathrm{exp}(iHt/\mathrm{})|0|^2`$. Asymptotically, $`p(\mathrm{})`$ (or a smoothed average) could be anywhere in $`(0,1)`$.
In Gaveau and Schulman (1995) the emphasis was on the mathematical mechanism, namely the formation of what could be called a plasmon mode Fano (1992). Suppose a single mode is coupled to $`N`$ levels, so that neither that mode nor the $`N`$ levels are eigenfunctions of the full Hamiltonian. In the usual decay scenario, when that single mode is expanded in eigenstates of the full Hamiltonian its amplitude in each is of order $`1/\sqrt{N}`$. Evolving in time, if you wait long enough there is “nothing” left in the original state; in the limit $`N\mathrm{}`$ the decay is complete. But it can also happen, as it does for plasmons and for the systems studied in Gaveau and Schulman (1995), that the original state has significant \[not O($`1/\sqrt{N}`$)\] overlap with a true eigenstate, and as a result only decays to a value dependent on the value of that overlap.
We illustrate both the gradual and catastrophic failure of the Golden Rule for the $`(2N+2)\times (2N+2)`$ Hamiltonian
$$H=\left(\begin{array}{cc}h& C^{}\\ C& \mathrm{\Omega }\end{array}\right)$$
(48)
where $`C`$ is a column vector of coupling coefficients, $`h`$, and $`\mathrm{\Omega }=\text{diag}(\omega __N,\mathrm{}\omega __N)`$, with $`\omega _n[E_1,E_2]`$ ($`\mathrm{}=1`$). (Every decay system can be brought to this form.) Since threshold behavior will prove to be our main preoccupation, behavior at the other end of the energy range is irrelevant. For convenience we take the energy range to be $`[E/2,E/2]`$ and allow $`h`$ to approach $`E/2`$ from below. The $`\omega _n`$ are taken to be uniformly spaced—what matters for decay (also when going beyond the Golden Rule) is the product $`\rho |\gamma |^2`$ with $`\rho `$ the density of states and $`\gamma =C/\sqrt{N}`$, which is the appropriate scaling with $`N`$. Therefore we can make $`\rho `$ a constant \[$`(2N+1)/E`$\], and build dimensional or other density of states features into $`C`$. We shall take $`C`$ of the form $`C(c/\sqrt{N})\varphi `$, with $`c`$ a constant and $`\varphi `$ of the form $`\varphi _n=[1(n/(N+1))^2]^\delta `$, with $`n=N,\mathrm{},N`$, and $`\delta `$ another constant. If $`h`$ is placed near level-$`n`$ (with its associated $`\omega =\omega _n`$), the Golden Rule prediction for the decay rate is $`\mathrm{\Gamma }_{\text{GR}}=2\pi \rho (\omega )|\gamma (\omega )|^2`$, which in this case becomes $`\mathrm{\Gamma }_{\text{GR}}=2\pi (2/E)|c|^2|\varphi _n|^2`$ not (d).
In Fig. 16 we show how the Golden rule prediction gradually declines in accuracy as $`hE/2`$. For all points plotted, there is exponential decay to great accuracy and to values of $`p`$ that reach 10<sup>-4</sup> or less. (As lifetimes increase, $`N`$ is also increased to allow the calculation to go to long times and avoid a Poincaré recurrence. For the same $`c`$, $`\delta `$ and $`h`$, changing $`N`$ had little effect. More on this later.) Finally, for the parameter values stated in the figure, at about $`h=0.92E/2`$ a plasmon mode develops. At this point exponential decay is completely lost. In Fig. 17 we show decay behavior for $`h=0.9E/2`$, side-by-side with the corresponding graph for $`h=0.94`$, when a plasmon exists. With the formation of this excitation, a significant fraction of the initial amplitude never decays. This has nothing to do with Poincaré recurrence. The same asymptotic probability as well as overall pattern of the curve obtains whether we did the calculation with dimension 102 or 1002 ($`N=50`$ or 500, and several values in between). This is the continuum behavior calculated in Gaveau and Schulman (1995).
Remark: Although the slight shift in $`h`$ above creates markedly different long-term behavior, for short time the quantum Zeno effect for the two parameter sets is similar. In Fig. 18 we show in greater detail the “Zeno era,” $`t\tau __Z`$, with $`\tau __Z`$ what I have called the Zeno time Schulman (2002, 1997), $`\tau __Z\mathrm{}/\sqrt{0|H^2|00|H|0^2}`$.
Remark: Formally the structure of the plasmon Hamiltonian and the structure of our $`H_0`$ and $`V_I`$ \[Eqs. (5/6) and Eqs. (7/8)\] are similar. The matrix for finding the classical local modes has the form Eq. (48) if one first diagonalizes the “chain” as we do in Sec. IV. The adjustable parameter sitting in the 1-1 position of the Hamiltonian of Eq. (48) is then $`\omega _1^2`$.
Next we check that the decay inhibition due to the plasmon is independent of the matrix size. This is shown in Fig. 19, where the only change with increasing $`N`$ is the deferral of the Poincaré recurrence.
### VI.2 Quantum time evolution of the breather
We turn to the time dependence under quantum evolution of the breather state. Specifically, we study $`|\psi _0|\mathrm{exp}(iHt/\mathrm{})|\psi _0|^2`$ for $`\psi _0`$ a local mode phonon (in particular, $`|0,\mathrm{},1`$) and $`H`$ the full Hamiltonian. This is shown in Fig. 20 not (e). As we saw in the examples of Sec. VI.1, where a plasmon has formed, initial decay is followed by stabilization bounded away from zero. In the present situation the value at which it stabilizes is much larger than the $``$0.6 (cf. Fig. 19) of those examples, which reflects the much larger amplitude of the local mode eigenstate in the true eigenstate of the Hamiltonian (cf. Table 1). Although the coupling pattern for the two cases is not the same (we will turn to this in a bit), the basic idea is the same: the initial state has order unity overlap with a true eigenfunction, and although it has coupling to the continuum, that coupling does not lead to exponential decay.
#### VI.2.1 Off-diagonal matrix elements of the singly-excited local mode
Since for the plasmon mode the emphasis is on the threshold structure of the coupling, it is of interest to explore those properties for the breather. This is also useful for assessing properties of the eigenstate when, because of the size of the system, numerical diagonalization is out of reach.
States of the local mode are both our initial states for decay and our unperturbed states for the numerical diagonalization. One of these is $`|\text{local}|0,0,0,\mathrm{},1`$ where, as in Sec. V, the “1” refers to the eigenvalue of the local mode number operator and the zeros to other phonon levels. The coupling, the analogue of “$`C`$” of Eq. (48), is
$$g(|n_1,n_2,\mathrm{})\text{local}|x_0^4|n_1,n_2,\mathrm{}.$$
(49)
In Fig. 21 is a logarithmic plot of $`g`$ as a function of energy. With states ordered by energy \[as in Eq. (48)\] this clearly does not provide a smooth function. However, the principal issue is not a precise resemblance to the plasmon paradigm. The plasmon was discussed not because of its specific coupling pattern, but rather as an example of how, although the initial state *does* couple to a continuum, the initial state still has finite (i.e., not going to zero with system size) overlap with a true eigenstate of the system. The pattern seen in Fig. 21, shown for 14 oscillators, is the same as one sees for 4, 6, 8 or any number that we have been able to study. As such, the substantial overlap of the true eigenstate (induced by the nonlinear $`\lambda `$) with the local mode (induced by $`\omega _1`$) should continue as system size grows, since it these matrix elements that determine the overlap.
Remark: Note that stability may not persist in all dimensions; certainly threshold features of the density of states and spectrum are affected by dimension, and the usual intuitions regarding Fermi’s Golden rule may again hold sway. In Schulman et al. (2002) we made the point that the symmetry breaking of the Jahn-Teller effect makes this a one-dimensional problem, significantly enhancing the possibility of classical breathers. The same may be true quantum mechanically. It may even be that this plays a role in the temperature-dependent decay of the breather (through an effective increase in dimension), as evidenced by the high-temperature disappearance of anomalous decay in doped alkali halides Mihóková et al. (2002a). For numerical diagonalization, increase in dimension is difficult computationally. However, for the path integral there should be little problem, and indeed 3-dimensional phonon coupling is used in many significant problems (see for example Weiss (1999)).
## VII Conclusions
Both numerical diagonalization and path integration automatically deal with infinite-lifetime eigenstates of the Hamiltonian. The principal issue is therefore to show localization for the system we study. That system is not the full periodic system, but rather a version that neglects quantum-tunneling delocalization \[an order $`\mathrm{exp}([\text{positive constant}]/\mathrm{})`$ effect\], but otherwise differs little from the full translationally invariant system.
Localization is demonstrated in two ways. For the path integral we show that increasing value of the quartic coupling constant (“$`\lambda `$”) decouples the rapidly oscillating atom from distant atoms, in the same way the this phenomenon occurs for a quartic localized mode. For the method of numerical diagonalization we show that the breather state is, up to tiny corrections, entirely constructed of states of an appropriate quadratic local mode. We remark that there is no principle saying that such tiny corrections represent either delocalization or decay. For example, the harmonic oscillator state, $`\psi (x)\mathrm{exp}(m\omega x^2/2)`$ is certainly a localized object, even though it is nonzero for all $`x`$. As to decay, we devote an entire section to disabusing anyone of the notion that decay is a mathematical imperative, although what lies behind the non-decay is merely the fact that under the full Hamiltonian the overlap of a state $`\psi _0`$ with itself does not go to zero if $`\psi _0`$ has finite (bounded away from zero) overlap with an eigenstate of the full Hamiltonian.
It follows that the quantized breather is stable, except against spontaneous tunneling, a process that is in general far slower than the decay rates that have been ascribed to the breather.
###### Acknowledgements.
We thank Bernard Gaveau for useful discussions. This work was supported by NSF grant PHY 00 99471 and Czech grants ME 587 and GA AVCR A1010210. |
warning/0507/cond-mat0507169.html | ar5iv | text | # Orbital and spin physics in LiNiO2 and NaNiO2
## 1 Introduction
The low-temperature magnetic behaviour of LiNiO<sub>2</sub> has remained puzzling ever since its peculiar properties were discovered . For no apparent reason, LiNiO<sub>2</sub> does not show magnetic order nor a cooperative Jahn-Teller effect down to the lowest temperatures, which is very different from the conventional behaviour observed in its sister compound NaNiO<sub>2</sub>.
Structurally the two systems represent an interesting special category within the class of correlated transition metal (TM) oxides. LiNiO<sub>2</sub> has a layered structure \[see figure 1(a)\]: it is rhombohedral, consisting of successive (111) planes occupied by Li<sup>+</sup>, O<sup>2-</sup>, Ni<sup>3+</sup>, and O<sup>2-</sup> ions. Thus the Ni<sup>3+</sup> ions are on a triangular lattice, with each direct Ni–Ni bond lying along the diagonal of a nearly square Ni–O–Ni–O plaquette, the Ni–O–Ni bonds being close to 90 degrees. This is distinct from the more commonly encountered situation where the bond between two transition metal ions through the ligand ion connecting them is close to linear (180 degrees), as e.g. in the perovskites. As pointed out by Mostovoy and Khomskii , this difference should have important consequences for the orbital and magnetic superexchange (SE) in LiNiO<sub>2</sub>, since the SE within a Ni plane with Ni$`{}_{}{}^{3+}(t_{2g}^6e_g^1)`$ ions would originate predominantly from virtual charge transfer excitations $`e^12p^6e^1e^22p^5e^1e^22p^4e^2`$ along the 90 degrees Ni–O–Ni bonds.
Over the years, a number of experiments (magnetic susceptibility, ESR, NMR, neutron scattering) have been performed on LiNiO<sub>2</sub> samples of varying stoichiometry . From these data one has concluded that the presence of excess Ni ions in the lithium layers introduces extra ferromagnetic (FM) coupling between the nickel layers. In addition, for the samples closest to perfect stoichiometry a positive Curie-Weiss temperature was found, indicating that the in-plane exchange is also FM. This is not in accordance with the description given initially by Hirakawa et al, namely that LiNiO<sub>2</sub> would be a triangular lattice antiferromagnet (TALAF). Actually, this assumption was the original motivation for performing magnetic measurements on LiNiO<sub>2</sub>, since the TALAF for spin $`S=1/2`$ is a frustrated system and the ground state might be some kind of quantum liquid instead of the classical 120 rotated spin arrangement.
At first sight the FM nature of the in-plane correlations is not surprising when one looks at the three-dimensional (3D) crystal structure (figure 1). As the neighbouring Ni<sup>3+</sup> ions are connected via two Ni-O-Ni bridges, one has the case of 90 SE, which the classical Goodenough-Kanamori-Anderson rules apparently predict to be FM. In a strong crystal field the ground state of Ni<sup>3+</sup> is the low-spin ($`t_{2g}^6e_g^1`$) configuration, so that direct exchange between $`t_{2g}`$ electrons does not occur. The exchange interactions between different nickel layers should also be FM , so that one naively expects a long-range ordered FM ground state. However, for LiNiO<sub>2</sub> no long-range magnetic order was found down to temperatures very close to $`0`$ K and the magnetic susceptibility gradually diverges, giving the impression that the FM correlations mysteriously disappear. A suggestion was made by Feiner, Oleś and Zaanen that the $`e_g`$ orbital degeneracy of the Ni<sup>3+</sup> ion and partly antiferromagnetic (AF) interactions might be responsible for this peculiar behaviour.
The issue was then addressed by Mostovoy and Khomskii (MK) in an important paper in which they proposed a realistic spin-and-orbital model for the Ni plane, which includes the Coulomb repulsion and the Hund’s rule exchange splitting on oxygen. They arrived at the conclusion that there is a huge degeneracy in the orbital sector, which is not resolved at the mean-field (MF) level. Yet orbital order is favoured over an orbital liquid state by the order-out-of-disorder mechanism, while they claimed that anyway the magnetic interaction is always FM . From the absence of an orbital ordered state in LiNiO<sub>2</sub> they concluded that the difference between LiNiO<sub>2</sub> and NaNiO<sub>2</sub> is probably extrinsic, due to disorder or electron-lattice interaction. Apparently this has now become the predominant view, and is as yet not inconsistent with experiments.
However, a conclusion concerning the nature of the magnetic interactions and the origin of the peculiar properties of LiNiO<sub>2</sub> had better be drawn only after the theoretical prediction for the intrinsic in-plane behaviour is fully established. We believe that this is not the case and therefore reanalyze the situation in this paper. Our finding is that upon inclusion of the Hund’s rule splitting also on the Ni ions and correction of what is apparently a mistake in the MK analysis, both FM and AF interactions can occur in the Ni plane, depending upon the orbital arrangement. Admittedly, this still leaves the difference between LiNiO<sub>2</sub> and NaNiO<sub>2</sub> to be explained, but it reopens the case for an intrinsic mechanism since different orbital phases in the Ni plane could give rise to different magnetic interactions.
The paper is organized as follows. In section 2 we present the notation used in the following sections for the spin and orbital operators used to derive the microscopic model. In section 3 we review some issues concerning SE, in particular with regard to the Ni–O–Ni 90 bond, and compare this case with the standard situation encountered for 180 bonds in TM perovskites. The spin-orbital SE model for the triangular Ni planes of LiNiO<sub>2</sub> and NaNiO<sub>2</sub> is presented in section 4. Next we discuss the strong frustration of the orbital and spin interactions in this model, and we investigate its consequences and present possible ground states, obtained using MF theory for a pair of Ni ions (section 5) and for the entire plane (section 6). The implications of the model for the physical properties of LiNiO<sub>2</sub> and NaNiO<sub>2</sub> are discussed in section 7. Finally, in section 8 the main conclusions and a summary are given. Technical details of the derivation of the model are presented in A.
## 2 Pseudospin formalism for degenerate $`e_g`$ orbitals
For the twofold degenerate $`e_g`$ orbital state at each site local operators corresponding to pseudospin $`T=1/2`$ are introduced, i.e. $`T_i^x`$, $`T_i^y`$ and $`T_i^z`$, acting as half the Pauli matrices $`\sigma ^x,\sigma ^y`$ and $`\sigma ^z`$ on the two-dimensional (2D) orbital Hilbert space at site $`i`$ with basis
$$\left(\genfrac{}{}{0pt}{}{1}{0}\right)_i|izd_{i,3z^2r^2}^{}|0,\left(\genfrac{}{}{0pt}{}{0}{1}\right)_i|i\overline{z}d_{i,x^2y^2}^{}|0.$$
(1)
A general superposition is given by
$$|i\theta _i=\mathrm{cos}(\theta _i/2)\left(\genfrac{}{}{0pt}{}{1}{0}\right)_i+\mathrm{sin}(\theta _i/2)\left(\genfrac{}{}{0pt}{}{0}{1}\right)_i,$$
(2)
which for example at $`\theta _i=\frac{\pi }{3}`$ corresponds to a $`d_{z^2x^2}`$ orbital. The expectation values of the pseudospin operators in the general orbital state (2) are
$$T_i^z=\frac{1}{2}\mathrm{cos}\theta _i,T_i^x=\frac{1}{2}\mathrm{sin}\theta _i,T_i^y=0.$$
(3)
In order to make the formalism more flexible and to include explicitly the cubic symmetry of the $`e_g`$ orbitals, it is convenient to define two more equivalent basis sets by
$`|i\alpha d_{i,3\alpha ^2r^2}^{}|0,`$ $`(\{\alpha ,\beta ,\gamma \}\text{a cyclic}`$
$`|i\overline{\alpha }d_{i,\beta ^2\gamma ^2}^{}|0,`$ $`\text{permutation of}\{x,y,z\})`$ (4)
and corresponding rotated pseudospin operators $`I_i^\alpha `$ and $`\overline{I}_i^\alpha `$ behaving like $`T_i^z`$ and $`T_i^x`$ with respect to those basis sets, i.e.
$`I_i^x=\frac{1}{2}T_i^z\frac{\sqrt{3}}{2}T_i^x,`$ $`\overline{I}_i^x=+\frac{\sqrt{3}}{2}T_i^z\frac{1}{2}T_i^x,`$
$`I_i^y=\frac{1}{2}T_i^z+\frac{\sqrt{3}}{2}T_i^x,`$ $`\overline{I}_i^y=\frac{\sqrt{3}}{2}T_i^z\frac{1}{2}T_i^x,`$
$`I_i^z=T_i^z,`$ $`\overline{I}_i^z=T_i^x,`$ (5)
which satisfy the identities
$$I_i^x+I_i^y+I_i^z=0,\overline{I}_i^x+\overline{I}_i^y+\overline{I}_i^z=0.$$
(6)
We can now introduce on-site orbital projection operators by
$$𝒫_i^\alpha =(\frac{1}{2}\mathrm{I}_i+I_i^\alpha ),𝒫_i^{\overline{\alpha }}=(\frac{1}{2}\mathrm{I}_iI_i^\alpha ),$$
(7)
where $`\mathrm{I}_i`$ is the unit operator in the 2D orbital Hilbert space at site $`i`$.
We further define two sets of (mutually dependent) orbital-pair operators, $`\{_{ij}^{\alpha \beta },𝒥_{ij}^{\alpha \beta }\}`$ and $`\{𝒱_{ij}^{\alpha \beta },𝒲_{ij}^{\alpha \beta }\}`$, all of which refer to a pair of nearest neighbour TM ions with their bond $`ij`$ lying in the $`\alpha \beta `$ plane (see figure 1(b)),
$`_{ij}^{\alpha \beta }`$ $`=`$ $`(\mathrm{I}_i+𝒫_i^{\overline{\alpha }})(\mathrm{I}_j+𝒫_j^{\overline{\beta }})+(\mathrm{I}_i+𝒫_i^{\overline{\beta }})(\mathrm{I}_j+𝒫_j^{\overline{\alpha }})`$ (8)
$`=`$ $`(\frac{3}{2}\mathrm{I}_iI_i^\alpha )(\frac{3}{2}\mathrm{I}_jI_j^\beta )+(\frac{3}{2}\mathrm{I}_iI_i^\beta )(\frac{3}{2}\mathrm{I}_jI_j^\alpha ),`$
$`𝒥_{ij}^{\alpha \beta }`$ $`=`$ $`(\mathrm{I}_i𝒫_i^{\overline{\alpha }})(\mathrm{I}_j𝒫_j^{\overline{\beta }})+(\mathrm{I}_i𝒫_i^{\overline{\beta }})(\mathrm{I}_j𝒫_j^{\overline{\alpha }})`$ (9)
$`=`$ $`(\frac{1}{2}\mathrm{I}_i+I_i^\alpha )(\frac{1}{2}\mathrm{I}_j+I_j^\beta )+(\frac{1}{2}\mathrm{I}_i+I_i^\beta )(\frac{1}{2}\mathrm{I}_j+I_j^\alpha ),`$
$`𝒱_{ij}^{\alpha \beta }`$ $`=`$ $`\mathrm{I}_i(I_j^\alpha +I_j^\beta )(I_i^\alpha +I_i^\beta )\mathrm{I}_j=\mathrm{I}_iI_j^\gamma +I_i^\gamma \mathrm{I}_j,`$ (10)
$`𝒲_{ij}^{\alpha \beta }`$ $`=`$ $`2(I_i^\alpha I_j^\beta +I_i^\beta I_j^\alpha )=2(I_i^\gamma I_j^\gamma I_i^\alpha I_j^\alpha I_i^\beta I_j^\beta ).`$ (11)
Their expectation values in a pair state $`|i\theta _i|j\theta _j`$ are given by
$`_{ij}^{\alpha \beta }`$ $`=`$ $`\frac{1}{8}[35+12\mathrm{cos}(\theta _++\chi _\gamma )\mathrm{cos}\theta _{}+4\mathrm{cos}^2(\theta _++\chi _\gamma )2\mathrm{cos}^2\theta _{}],`$ (12)
$`𝒥_{ij}^{\alpha \beta }`$ $`=`$ $`\frac{1}{8}[34\mathrm{cos}(\theta _++\chi _\gamma )\mathrm{cos}\theta _{}+4\mathrm{cos}^2(\theta _++\chi _\gamma )2\mathrm{cos}^2\theta _{}],`$ (13)
$`𝒱_{ij}^{\alpha \beta }`$ $`=`$ $`\mathrm{cos}(\theta _++\chi _\gamma )\mathrm{cos}\theta _{},`$ (14)
$`𝒲_{ij}^{\alpha \beta }`$ $`=`$ $`\frac{1}{4}[4\mathrm{cos}^2(\theta _++\chi _\gamma )2\mathrm{cos}^2\theta _{}1]=\frac{1}{4}[2\mathrm{cos}(2\theta _++2\chi _\gamma )\mathrm{cos}2\theta _{}],`$ (15)
where $`\theta _\pm =(\theta _i\pm \theta _j)/2`$, and $`\chi _x=\frac{2\pi }{3}`$, $`\chi _y=\frac{2\pi }{3}`$, $`\chi _z=0`$.
Finally we introduce bond projection operators, needed for specifying the SE interactions between a pair of TM ions. For the spin part we will use the familiar projection operators for spin triplet and spin singlet,
$$\mathrm{Q}_{ij}^T=\frac{3}{4}\mathrm{𝟏}_{ij}+\mathrm{S}_i\mathrm{S}_j,\mathrm{Q}_{ij}^S=\frac{1}{4}\mathrm{𝟏}_{ij}\mathrm{S}_i\mathrm{S}_j,$$
(16)
where $`\mathrm{𝟏}_{ij}`$ is the unit operator in the four-dimensional (4D) spin Hilbert space on the bond $`ij`$. For the orbital part we will make use of
$`𝒬_{\mathrm{O},ij}^{\alpha \beta }`$ $`=`$ $`𝒫_i^\alpha 𝒫_j^\beta +𝒫_i^\beta 𝒫_j^\alpha ,`$ (17)
$`𝒬_{\mathrm{M},ij}^{\alpha \beta }`$ $`=`$ $`𝒫_i^\alpha 𝒫_j^{\overline{\beta }}+𝒫_i^{\overline{\alpha }}𝒫_j^\beta +𝒫_i^\beta 𝒫_j^{\overline{\alpha }}+𝒫_i^{\overline{\beta }}𝒫_j^\alpha ,`$ (18)
$`𝒬_{\mathrm{N},ij}^{\alpha \beta }`$ $`=`$ $`𝒫_i^{\overline{\alpha }}𝒫_j^{\overline{\beta }}+𝒫_i^{\overline{\beta }}𝒫_j^{\overline{\alpha }},`$ (19)
where the labelling will be explained in section 4. They are conveniently expanded in terms of the operators defined above, according to (with indices omitted for clarity)
$$\begin{array}{cccccccccccccc}\hfill 𝒬_\mathrm{O}& =& & & & & \hfill 𝒥& =& \hfill \frac{1}{2}\mathrm{I}& & \hfill \frac{1}{2}𝒱& +& \hfill \frac{1}{2}𝒲& \\ \hfill 𝒬_\mathrm{M}& =& \hfill 4\mathrm{I}& & \hfill \frac{1}{2}& & \hfill \frac{3}{2}𝒥& =& \hfill \mathrm{I}& & & & \hfill 𝒲& \\ \hfill 𝒬_\mathrm{N}& =& \hfill 2\mathrm{I}& +& \hfill \frac{1}{2}& +& \hfill \frac{1}{2}𝒥& =& \hfill \frac{1}{2}\mathrm{I}& +& \hfill \frac{1}{2}𝒱& +& \hfill \frac{1}{2}𝒲& ,\hfill \end{array}$$
(20)
where $`\mathrm{I}\mathrm{I}_{ij}`$ now denotes the identity in the 4D orbital-pair Hilbert space.
## 3 Superexchange for degenerate $`e_g`$ orbitals
If degenerate $`e_g`$ orbitals are partly filled, a spin-orbital Hamiltonian can be constructed from a degenerate-band Hubbard model in much the same way as one derives the AF Heisenberg model from the single-band Hubbard model at half-filling. The Hamiltonian of the $`e_g`$-band Hubbard model consists of two parts: there is a hopping term $`𝒯`$ modelling transfer of electrons between nearest-neighbour transition metal sites and a Hubbard (or Coulomb) term $`𝒰`$ describing the on-site interactions. In the situation where the number of sites equals the number of electrons, the ground state for $`𝒯=0`$ has twofold spin and twofold orbital degeneracy on each site. When we allow for a nonzero $`𝒯`$ as a small perturbation, the 4<sup>N</sup>-fold degeneracy is lifted by virtual electron hopping involving excited states, and the effective Hamiltonian in second order perturbation theory is
$$_{\mathrm{eff}}=\underset{n/\mathrm{G}}{}𝒫_\mathrm{G}\left[𝒯|n\frac{1}{E_\mathrm{G}E_n}n|𝒯\right]𝒫_\mathrm{G},$$
(21)
where $`𝒫_\mathrm{G}`$ is a projection operator onto the ground state manifold of $`𝒰`$. The spin-orbital Hamiltonian is obtained by writing out (21) in terms of spin and orbital projection operators at each site $`i`$.
For a cubic lattice, such as in the perovskites KCuF<sub>3</sub> and K<sub>2</sub>CuF<sub>4</sub>, one can describe by this approach also the 180 SE, thus treating it formally as direct exchange. The hopping term is taken as
$$𝒯=\overline{t}\underset{\alpha }{}\underset{ij\alpha }{}\underset{\sigma }{}d_{i\alpha \sigma }^{}d_{j\alpha \sigma },$$
(22)
expressing that hopping is only allowed between nearest-neighbour directional orbitals $`|\alpha d_{3\alpha ^2r^2}^{}|0`$ oriented along the connecting $`\alpha `$ axis ($`\alpha `$ being $`x`$, $`y`$ or $`z`$). In the following we will therefore call the $`\alpha `$-orbitals ‘hopping’ orbitals – since electrons in these orbitals can hop and contribute to the kinetic energy, and similarly call ‘non-hopping’ orbitals the $`\overline{\alpha }`$-orbitals ($`|\overline{\alpha }d_{\beta ^2\gamma ^2}^{}|0`$), orthogonal to the $`\alpha `$-orbital and oriented perpendicular to the bond. In a one-dimensional chain this situation leads to a characteristic competition between itinerant and localized phases .
The on-site interactions on a TM ion can be represented by
$`𝒰_{\mathrm{TM}}`$ $`=`$ $`U{\displaystyle \underset{i\lambda }{}}n_{i\lambda }n_{i\lambda }+\left(U{\displaystyle \frac{5}{2}}J_\mathrm{H}\right){\displaystyle \underset{i,\lambda <\mu }{}}n_{i\lambda }n_{i\mu }2J_\mathrm{H}{\displaystyle \underset{i,\lambda <\mu }{}}\mathrm{s}_{i\lambda }\mathrm{s}_{i\mu }`$ (23)
$`+`$ $`J_\mathrm{H}{\displaystyle \underset{i,\lambda <\mu }{}}\left(d_{i\lambda }^{}d_{i\lambda }^{}d_{i\mu }d_{i\mu }+d_{i\mu }^{}d_{i\mu }^{}d_{i\lambda }d_{i\lambda }\right),`$
and are characterized by two parameters: the intraorbital Coulomb energy $`U`$ and the exchange energy $`J_\mathrm{H}`$. The interorbital terms in equation (23) describe electron interactions between pairs of orthogonal orbitals, i.e., in the present subspace of $`e_g`$ orbitals one has $`\lambda ,\mu \{\alpha ,\overline{\alpha }\}`$. The excited states, generated in the virtual $`dd`$ transitions and relevant for SE, have two electrons on the same ion, and $`𝒰`$ contributes a Coulomb repulsion energy $`U`$. In the $`\overline{t}U`$ limit one thus derives an effective low-energy spin-orbital Hamiltonian with coupling constant $`J_{\mathrm{SE}}\overline{t}^{\mathrm{\hspace{0.17em}2}}/U`$, in which spin and orbital degrees of freedom are interrelated. By taking also the Hund’s rule exchange $`J_\mathrm{H}`$ into account one removes the classical degeneracy of magnetically ordered phases . The stable phase at low temperatures has long-range orbital order of a particular type of mixed orbitals, leading to AF interactions along the $`c`$-axis and FM interactions in the $`(a,b)`$-plane. This ordering was verified experimentally and has been shown to be stable with respect to quantum fluctuations for large $`J_\mathrm{H}`$ .
In a more realistic treatment of SE one takes into account explicitly that the hopping takes place via the ligand oxygen ion. The hopping term is then taken as
$$𝒯=t\underset{\alpha }{}\underset{ij\alpha }{}\underset{\sigma }{}\left(d_{i\alpha \sigma }^{}p_{j\alpha \sigma }+p_{j\alpha \sigma }^{}d_{i\alpha \sigma }\right),$$
(24)
and describes charge transfer with amplitude $`tt_\sigma `$ between the TM-orbital $`d_{3\alpha ^2r^2}`$ and the oxygen $`\sigma `$-type $`p`$-orbital $`p_\alpha `$, where again the orbitals are oriented along the connecting $`\alpha `$-axis. The on-site interaction on oxygen is given by
$`𝒰_\mathrm{O}`$ $`=`$ $`U_\mathrm{O}{\displaystyle \underset{j\lambda }{}}n_{j\lambda }n_{j\lambda }+\left(U_\mathrm{O}{\displaystyle \frac{5}{2}}J_\mathrm{O}\right){\displaystyle \underset{j,\lambda <\mu }{}}n_{j\lambda }n_{j\mu }2J_\mathrm{O}{\displaystyle \underset{j,\lambda <\mu }{}}\mathrm{s}_{j\lambda }\mathrm{s}_{j\mu }`$ (25)
$`+`$ $`J_\mathrm{O}{\displaystyle \underset{j,\lambda <\mu }{}}\left(p_{j\lambda }^{}p_{j\lambda }^{}p_{j\mu }p_{j\mu }+p_{j\mu }^{}p_{j\mu }^{}p_{j\lambda }p_{j\lambda }\right),`$
with intraorbital Coulomb energy $`U_\mathrm{O}`$ and exchange energy $`J_\mathrm{O}`$. Here the operators $`n_{j\lambda \sigma }`$, $`n_{j\lambda }`$ and $`\mathrm{s}_{j\lambda }`$ refer to an oxygen ion at site $`j`$. For two-hole excitations in the three $`2p`$ orbitals, as arises when the SE is derived (see below), this local problem and the excitation spectrum are isomorphic to those for three $`t_{2g}`$ orbitals filled by two electrons as in the vanadates . The effective Hamiltonian is now obtained in fourth order perturbation theory,
$$_{\mathrm{eff}}=\underset{k,l,m/\mathrm{G}}{}𝒫_\mathrm{G}\left[𝒯|k\frac{1}{E_\mathrm{G}E_k}k|𝒯|l\frac{1}{E_\mathrm{G}E_l}l|𝒯|m\frac{1}{E_\mathrm{G}E_m}m|𝒯\right]𝒫_\mathrm{G}.$$
(26)
In the case of a 180 TM–O–TM bond $`ij`$ this is unproblematic. For the so-called $`U`$-term (Anderson or delocalization process) , where an electron is effectively transferred from one TM ion to the other TM ion, schematically represented as
$$e^1p^6e^1e^1p^5e^2e^0p^6e^2e^1p^5e^2e^1p^6e^1,$$
one can simply replace $`\overline{t}`$ by $`t_\sigma ^2/\mathrm{\Delta }`$, where $`\mathrm{\Delta }`$ is the excitation energy for transferring an electron from O to TM, and so the coupling constant in the effective Hamiltonian becomes
$$J_U\frac{t_\sigma ^4}{\mathrm{\Delta }^2}\frac{1}{U}.$$
(27)
Note that the orbital filled in the second step is necessarily the same, also as regards spin, as the one emptied in the first step. The so-called $`\mathrm{\Delta }`$-term (Goodenough process or correlation effect) involves instead electron transfer of two electrons from the connecting oxygen ion, one to each of the TM neighbours,
$$e^1p^6e^1e^1p^5e^2e^2p^4e^2e^1p^5e^2e^1p^6e^1.$$
Here the oxygen $`2p^4`$ configuration involved has two holes with opposite spin on the same $`\sigma `$-type $`p`$-orbital, giving always the same intermediate state at the oxygen ion with energy $`2\mathrm{\Delta }+U_\mathrm{O}`$, and the contributions to the effective Hamiltonian have coupling constant
$$J_\mathrm{\Delta }\frac{t_\sigma ^4}{\mathrm{\Delta }^2}\left(\frac{1}{2\mathrm{\Delta }+U_\mathrm{O}}\frac{1}{2\mathrm{\Delta }}\right)=\frac{t_\sigma ^4}{\mathrm{\Delta }^2}\frac{U_\mathrm{O}}{2\mathrm{\Delta }(2\mathrm{\Delta }+U_\mathrm{O})}.$$
(28)
The reason for the subtraction of the term $`1/2\mathrm{\Delta }`$ will be discussed below. For a 180 TM–O–TM bond the two processes (Anderson and Goodenough) make qualitatively similar contributions to the effective Hamiltonian (at least for a single $`e_g`$ electron on each TM ion ) because both involve $`\sigma `$-type hopping. As the terms contributed to $`_{\mathrm{eff}}`$ have identical form, they can be formally generated by the second-order formalism of direct exchange above, even though this models only the process giving the $`U`$-term in the SE. To obtain the coupling constants one can simply add equations (27) and (28). As these are generally of the same order of magnitude, inclusion of the $`\mathrm{\Delta }`$-term is important quantitatively, and is essential to describe the trend in the strength of SE within the $`3d`$ TM series .
In the case of a 90 TM–O–TM bond, as on a triangular lattice, the situation is very different. In the $`U`$-term (Anderson) process the oxygen orbital through which the electron is being transferred is now $`\sigma `$-type for one TM ion but $`\pi `$-type for the other one. Therefore this process can only occur if it involves a (deep-lying) $`t_{2g}`$ orbital,
$$(t_2^6e^1)p^6(t_2^6e^1)(t_2^6e^1)p^5(t_2^6e^2)(t_2^5e^1)p^6(t_2^6e^2),$$
and thus contributes to the effective Hamiltonian a term with coupling constant
$$J_U\frac{t_\sigma ^2t_\pi ^2}{\mathrm{\Delta }^2}\frac{1}{U+\mathrm{\Delta }_{\mathrm{CF}}},$$
(29)
since the energy of the middle intermediate state is increased by the crystal field splitting $`\mathrm{\Delta }_{\mathrm{CF}}`$, and the hopping parameter for the $`t_{2g}`$ orbital is $`t_\pi `$ instead of $`t_\sigma `$. In the charge transfer terms ($`\mathrm{\Delta }`$ process) the oxygen $`2p^4`$ states now involve two holes on different $`p`$-orbitals, each of $`\sigma `$-type but with respect to a different TM neighbour. This implies that the oxygen Coulomb interaction involved is now the interorbital interaction $`U_\mathrm{O}^{}U_p`$ instead of the intraorbital interaction $`U_\mathrm{O}=U_\mathrm{O}^{}+2J_p`$. Moreover it leads to a singlet-triplet splitting, with energies $`U_s=U_p+J_p`$ and $`U_t=U_pJ_p`$, where $`J_pJ_\mathrm{O}`$ is the (Hund’s rule) exchange on oxygen, and thus the contributions to the effective Hamiltonian have coupling constants
$$J_\mathrm{\Delta }^\pm \frac{t_\sigma ^4}{\mathrm{\Delta }^2}\left(\frac{1}{2\mathrm{\Delta }+U_p\pm J_p}\frac{1}{2\mathrm{\Delta }}\right)=\frac{t_\sigma ^4}{\mathrm{\Delta }^2}\frac{U_p\pm J_p}{2\mathrm{\Delta }(2\mathrm{\Delta }+U_p\pm J_p)}.$$
(30)
From equations (29) and (30) one observes that the dominant contribution to the SE comes from the $`\mathrm{\Delta }`$-term (Goodenough process), since $`t_\sigma ^2/t_\pi ^24`$ , while $`U+\mathrm{\Delta }_{\mathrm{CF}}2\mathrm{\Delta }`$ and $`U_p2\mathrm{\Delta }`$ . This was already pointed out in reference , as well as the fact that the magnetic interaction should then be FM, as equation (30) shows, which is one of the famous Goodenough-Kanamori-Anderson rules. Nevertheless, one should be careful not to jump to conclusions here, since the implications of the orbital SE interaction were not fully considered by Goodenough (the case explicitly considered was the SE between Ni<sup>2+</sup> ($`t_{2g}^6e_g^2`$ : $`{}_{}{}^{3}A_{2}^{}`$) ions , where the two $`e_g`$ orbitals are both occupied by one electron, and no orbital effects can arise). It is further clear that the dominant $`\mathrm{\Delta }`$-term (Goodenough contribution) cannot be represented well as an effective second order direct exchange. Yet this was attempted in a recent paper : this has the merit of deriving the most general form of the effective Hamiltonian purely based on symmetry considerations, but it cannot capture the dependence on the most relevant parameters $`U_p`$ and $`J_p`$ (actually only the weaker $`U`$-term SE, dependent upon the splitting of the intermediate Ni<sup>2+</sup> configurations, is being described).
A notable peculiarity of the $`\mathrm{\Delta }`$-term, already included in equations (28) and (30), is that subtraction is needed of the contribution that would have been obtained if the electrons transferred to the TM ions would have come from two different oxygen ions and not from the connecting oxygen ligand, as pointed out by Mostovoy and Khomskii . The necessity for this subtraction can be understood as follows. The reference state (‘vacuum’) should be considered to be renormalized by all possible fourth order uncorrelated hopping sequences \[see figure 2(a)\], each contributing a term $`t_\sigma ^4/(2\mathrm{\Delta }^3)`$. In most cases they cancel out because the three contributions corresponding to hops 2 and 3 in figure 2(a) made along the three cubic axes with the axis of hops 1 and 4 kept fixed, add up to a constant, since the three projection operators $`𝒫_i^\alpha `$ do so because of equation (6), and thus only add to the vacuum energy. However, the cancellation fails if the contribution from one axis is missing because the oxygen ion there is joined with the other TM ion, as in figure 2(b). The cancellation is restored by adding this term to the other two and subtracting it from the SE for the pair of TM ions under consideration.
Note that the above correction implies a sign change of the coupling constant and so the opposite situation is favoured than one might naively expect. In particular, the largest diagonal SE is generally obtained for the configuration that permits the largest number of hopping sequences returning to itself, and so this configuration is now energetically disfavoured instead of favoured. This applies specifically to purely interorbital SE terms, whereas in spin-spin SE interactions, which generally originate from the difference of SE for two spin multiplets, the corrections cancel out and one gets the expected result.
## 4 Spin-orbital model for LiNiO<sub>2</sub> and NaNiO<sub>2</sub>
Let us now reconsider the derivation of the spin-orbital orbital model for the triangular lattice structure of LiNiO<sub>2</sub> , taking only the charge transfer process ($`\mathrm{\Delta }`$-term) into account, as done also by MK . So we consider two nearest neighbour Ni<sup>3+</sup> ions connected by two Ni–O–Ni 90 bridges, as in figure 1 or 2(b), and analyze the fourth order hopping sequences (which in this order of perturbation theory can be done for each bridge separately). In order that these virtual excitations lift the ground state degeneracy it is essential that the intermediate states are affected by the on-site Coulomb interactions on oxygen and/or nickel, described by $`𝒰`$, see equations (23) and (25).
As indicated above, the relevant states of the oxygen $`p^4`$ configuration, i.e. those occurring upon hopping, are a triplet $`{}_{}{}^{3}T_{1}^{}`$ and a singlet $`{}_{}{}^{1}T_{2}^{}`$, denoted for brevity by $`t`$ and $`s`$. They are split by $`2J_p`$, while moreover the interorbital Coulomb repulsion $`U_p`$ for $`p^4`$, absent for two $`p^5`$ configurations, must be taken into account. The relevant Ni $`e_g^2`$ states are $`{}_{}{}^{1}A_{1}^{}`$, $`{}_{}{}^{1}E`$ and $`{}_{}{}^{3}A_{2}^{}`$, for which we will use the abbreviations $`S`$ (singlet), $`D`$ (orbital doublet) and $`T`$ (triplet), respectively. These terms have equidistant energy levels given by $`U+J_\mathrm{H}`$, $`UJ_\mathrm{H}`$ and $`U3J_\mathrm{H}`$, with the triplet being lowest by Hund’s rule, where the Ni<sup>2+</sup> Coulomb repulsion and Hund’s rule coupling can be expressed in terms of Racah parameters as $`U=A+4B+3C`$ and $`J_\mathrm{H}=4B+C`$ .
The first transition in the hopping sequence is $`e^1p^6e^1e^2p^5e^1`$, with excitation energy
$$\mathrm{\Delta }_X=E_X^{\mathrm{Ni}}(e^2)E^{\mathrm{Ni}}(e^1)+E^\mathrm{O}(p^5)E^\mathrm{O}(p^6),$$
(31)
i.e. the charge transfer energy depends upon the Ni<sup>2+</sup> state accessed,
$$\mathrm{\Delta }_S=\mathrm{\Delta }+2J_\mathrm{H},\mathrm{\Delta }_D=\mathrm{\Delta },\mathrm{\Delta }_T=\mathrm{\Delta }2J_\mathrm{H},$$
(32)
where it is understood that $`UJ_\mathrm{H}`$ has been absorbed into $`\mathrm{\Delta }`$. Which states are reached depends on the $`e_g`$ electron already present on the Ni<sup>3+</sup> ion to which the electron hops, as illustrated in figure 3, which shows the six possible hopping channels. If the $`d_{3\alpha ^2r^2}`$ orbital is occupied, then a hop from the $`p_\alpha `$ oxygen orbital is only possible for an electron with opposite spin, and the excited states involved are the spin singlets $`{}_{}{}^{1}A_{1}^{}`$ and $`{}_{}{}^{1}E`$. On the other hand, if the $`d_{\beta ^2\gamma ^2}`$ orbital is occupied (and therefore the $`d_{3\alpha ^2r^2}`$ orbital is empty), then the spin of the hopping electron can have either sign, and the $`{}_{}{}^{1}E`$ and $`{}_{}{}^{3}A_{2}^{}`$ states are reached. In the second step $`e^2p^5e^1e^2p^4e^2`$, the excitation energy is raised further by $`\mathrm{\Delta }_Y+U_p\pm J_p`$, depending upon the $`e_g^2`$ state $`Y`$ accessed at the other Ni ion and the $`p^4`$ state ($`s`$ or $`t`$) at the oxygen. Figure 4 shows an example of such a double charge transfer excitation. In the third and fourth step the excitation is undone, either in the same or in reverse order.
To derive the complete spin-orbital Hamiltonian one must list all possible initial configurations, and for each of them list all possible hopping sequences that return to the ground state manifold. The initial and final states are described by means of projection operators, both for the orbital occupation and for the spin state, defined in section 2. The orbital bond projection operators specify whether the ‘hopping’ orbitals on the two Ni<sup>3+</sup> ions are both occupied by an electron (a situation denoted by ‘O’), whether one hopping orbital is occupied while the electron on the other ion is in the non-hopping orbital (‘M’, for mixed), or whether both non-hopping orbitals are occupied by the two electrons (‘N’). Since the contributions from the two Ni–O–Ni bridges are independent and may be simply added, the operators (1719) are defined to do so for each $`\alpha \beta `$ bond direction. Obviously these operators depend on the bond direction, and so this gives explicit orbital anisotropy. For specifying the overall spin state (i.e. formed by the two spins $`1/2`$ of the two Ni<sup>3+</sup> ions), we need the familiar projection operators (16) for spin triplet and spin singlet.
Collecting the contributions from all configurations and all hopping sequences with the help of diagrams like in figures 3 and 4 one then obtains the Hamiltonian in first instance in the form
$$_{\mathrm{eff}}=\underset{ij}{}\left(𝒬_{\mathrm{O},ij}^{\alpha \beta }[K_\mathrm{O}^\mathrm{T}\mathrm{Q}_{ij}^\mathrm{T}+K_\mathrm{O}^\mathrm{S}\mathrm{Q}_{ij}^\mathrm{S}]+𝒬_{\mathrm{M},ij}^{\alpha \beta }[K_\mathrm{M}^\mathrm{T}\mathrm{Q}_{ij}^\mathrm{T}+K_\mathrm{M}^\mathrm{S}\mathrm{Q}_{ij}^\mathrm{S}]+𝒬_{\mathrm{N},ij}^{\alpha \beta }[K_\mathrm{N}^\mathrm{T}\mathrm{Q}_{ij}^\mathrm{T}+K_\mathrm{N}^\mathrm{S}\mathrm{Q}_{ij}^\mathrm{S}]\right),$$
(33)
where it is understood that the form of the projection operators depends on the bond direction. In order to separate the spin dependent part from the purely orbital part this may be rewritten as
$$_{\mathrm{eff}}=\underset{ij}{}\left([J_\mathrm{O}^0𝒬_{\mathrm{O},ij}^{\alpha \beta }+J_\mathrm{M}^0𝒬_{\mathrm{M},ij}^{\alpha \beta }+J_\mathrm{N}^0𝒬_{\mathrm{N},ij}^{\alpha \beta }]\mathrm{𝟏}_{ij}+[J_\mathrm{O}^S𝒬_{\mathrm{O},ij}^{\alpha \beta }+J_\mathrm{M}^S𝒬_{\mathrm{M},ij}^{\alpha \beta }+J_\mathrm{N}^S𝒬_{\mathrm{N},ij}^{\alpha \beta }]\mathrm{S}_i\mathrm{S}_j\right),$$
(34)
with orbital and spin-orbital interactions
$$J_\mathrm{L}^0=\frac{3}{4}K_\mathrm{L}^\mathrm{T}+\frac{1}{4}K_\mathrm{L}^\mathrm{S},J_\mathrm{L}^S=K_\mathrm{L}^\mathrm{T}K_\mathrm{L}^\mathrm{S}(\mathrm{L}=\mathrm{O},\mathrm{M},\mathrm{N}).$$
(35)
To see how this works in detail and to compare with the analysis of MK, let us initially ignore the Hund’s exchange splitting on Ni (the full procedure is explained in more detail in A). If the initial configuration is N-type, the electrons hopping from the different $`2p`$ orbitals at the common oxygen ion into the empty hopping orbitals on the Ni neighbouring ions, can do so with their spins oriented in four ways: both up or both down, thus leaving the oxygen in the triplet $`p^4`$ state, or with their spins up-down or down-up, both with equal probability $`1/2`$ for leaving the oxygen in the triplet $`p^4`$ or in the singlet $`p^4`$ state, all in all making three triplet and one singlet contribution. As the $`e_g^2`$ terms on the Ni ions are all equivalent when $`J_H=0`$, this is independent of the initial orientation of the spins in the Ni non-hopping orbitals. Upon including an overall factor of 4, because both the two excitation transfers and the two deexcitation transfers can also be made in reverse order, one obtains
$$K_\mathrm{N}^\mathrm{T}=K_\mathrm{N}^\mathrm{S}=12[XtX]+4[XsX].$$
(36)
Here we denote the fourth-order perturbation expressions by giving a shorthand notation for the middle intermediate state (using here $`X`$ instead of $`S`$, $`D`$ or $`T`$, because we do not distinguish between those states yet). Their values are \[compare equation (30) above\]
$$[XtX]=\frac{t^4}{\mathrm{\Delta }^2}\left(\frac{1}{2\mathrm{\Delta }+U_pJ_p}+\frac{1}{2\mathrm{\Delta }}\right)=\frac{t^4}{\mathrm{\Delta }^2}\frac{U_pJ_p}{2\mathrm{\Delta }(2\mathrm{\Delta }+U_pJ_p)},$$
(37)
$$[XsX]=\frac{t^4}{\mathrm{\Delta }^2}\left(\frac{1}{2\mathrm{\Delta }+U_p+J_p}+\frac{1}{2\mathrm{\Delta }}\right)=\frac{t^4}{\mathrm{\Delta }^2}\frac{U_p+J_p}{2\mathrm{\Delta }(2\mathrm{\Delta }+U_p+J_p)}.$$
(38)
If the initial configuration is M-type, the result is equally independent of the initial Ni–Ni spin state. Although the spin of one transferred electron must be opposite to the spin of the electron occupying the hopping orbital on Ni, the other transferred electron can have its spin still either parallel to that of the first one, leaving a triplet $`p^4`$ state on oxygen, or antiparallel, yielding a triplet or a singlet with probability $`1/2`$ each. Including again the factor 4, one obtains
$$K_\mathrm{M}^\mathrm{T}=K_\mathrm{M}^\mathrm{S}=6[XtX]+2[XsX].$$
(39)
Only if the initial configuration is O-type, the result is different, because each of the transferred electrons must have its spin opposite to that of the electron already occupying the hopping orbital on Ni. So, if the spins of the electrons on Ni are parallel, the initial state thus being a Ni–Ni spin triplet, then the transferred electrons necessarily also have parallel spins, leaving oxygen in the triplet $`p^4`$ state. If the electrons on the Ni ions have antiparallel spins, i.e. being either in a spin triplet or in a spin singlet state depending upon the phasing between up-down and down-up, then the spins of the electrons being transferred are also antiparallel, and have the same phasing, so again the Ni–Ni spin triplet yields an oxygen $`p^4`$ triplet, while the Ni–Ni spin singlet yields an oxygen $`p^4`$ singlet. Therefore, upon inclusion of the factor 4,
$$K_\mathrm{O}^\mathrm{T}=4[XtX],K_\mathrm{O}^\mathrm{S}=4[XsX].$$
(40)
It follows that
$`J_\mathrm{N}^0=4J_T,J_\mathrm{M}^0=2J_T,J_\mathrm{O}^0=J_T,`$ (41)
$`J_\mathrm{N}^S=0,J_\mathrm{M}^S=0,J_\mathrm{O}^S=J_{TS},`$ (42)
$`J_T=\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}3}[XtX]+[XsX]={\displaystyle \frac{2t^4}{\mathrm{\Delta }^3}}{\displaystyle \frac{\mathrm{\Delta }(2U_pJ_p)+U_p^2J_p^2}{(2\mathrm{\Delta }+U_p)^2J_p^2}}{\displaystyle \frac{2t^4}{\mathrm{\Delta }^3}}{\displaystyle \frac{U_p}{2\mathrm{\Delta }+U_p}},`$ (43)
$`J_{TS}=4([XtX][XsX])={\displaystyle \frac{2t^4}{\mathrm{\Delta }^2}}{\displaystyle \frac{4J_p}{(2\mathrm{\Delta }+U_p)^2J_p^2}}{\displaystyle \frac{2t^4}{\mathrm{\Delta }^2}}{\displaystyle \frac{4J_p}{(2\mathrm{\Delta }+U_p)^2}},`$ (44)
where the final expressions on the right in equations (43) and (44) are the results in the limit $`J_p\mathrm{\Delta },U_p`$. So the effective Hamiltonian is
$`_{\mathrm{eff}}^{(0)}`$ $`=`$ $`{\displaystyle \underset{ij}{}}\left(J_T[𝒬_{\mathrm{O},ij}^{\alpha \beta }+2𝒬_{\mathrm{M},ij}^{\alpha \beta }+4𝒬_{\mathrm{N},ij}^{\alpha \beta }]\mathrm{𝟏}_{ij}J_{TS}𝒬_{\mathrm{O},ij}^{\alpha \beta }\mathrm{S}_i\mathrm{S}_j\right)`$ (45)
$`=`$ $`{\displaystyle \underset{ij}{}}\left(J_T_{ij}^{\alpha \beta }\mathrm{𝟏}_{ij}J_{TS}𝒥_{ij}^{\alpha \beta }\mathrm{S}_i\mathrm{S}_j\right),`$
with $`_{ij}^{\alpha \beta }`$ and $`𝒥_{ij}^{\alpha \beta }`$ given by (8) and (9), and understood to depend upon the direction of the bond $`ij`$.
Comparison with the results of MK shows that our analysis confirms their finding that the purely orbital interaction is stronger by one order of magnitude than any spin dependent interaction, $`J_TJ_{TS}`$, in agreement with the conjecture made by Reynaud et aland based on their experimental data . One recognizes that this comes about because, at $`J_\mathrm{H}=0`$, all spin dependence originates from the singlet-triplet splitting of the oxygen $`2p^4`$ configuration, and $`J_{TS}`$ is therefore smaller by a factor $`J_p/U_p0.1`$ than the orbital interaction $`J_T`$. Our analysis also confirms the form of the orbital interaction as being $`_{ij}^{\alpha \beta }`$. However, we find a different form for the orbital dependence of the mixed spin-orbital interaction, since MK give $`_{ij}^{\alpha \beta }`$ instead of $`𝒥_{ij}^{\alpha \beta }`$. Apparently their result is incorrect, because the reasoning above clearly demonstrates that only the O-type configuration can give spin dependence, and this is what is expressed by the operator $`𝒥^{\alpha \beta }=𝒫_\mathrm{O}^{\alpha \beta }`$. The difference is important, because the bounds on the expectation values are very different, $`\frac{49}{8}^{\alpha \beta }\frac{25}{8}`$ whereas $`\frac{9}{8}𝒥^{\alpha \beta }0`$ \[see equations (12) and (13)\], and so MK concluded that the spin exchange is effectively always FM for any orbital state, while we conclude that this is at best marginally so, i.e. the spin exchange interaction on a bond $`ij`$ given by (45) can vanish for some specific combination of orbital states at sites $`i`$ and $`j`$.
Moreover, the latter conclusion is only valid up to the present level of approximation, where only the Coulomb interactions on oxygen have been taken into account. If we include also the Coulomb interaction on nickel, i.e. allow $`J_\mathrm{H}`$ to be finite, and consider the leading correction to (45), we find that the spin exchange may even become weakly AF. To see this we first observe, as demonstrated in A, that quite generally the coupling constants satisfy
$$|J_\mathrm{O}^S|J_\mathrm{M}^S|J_\mathrm{N}^S|\mathrm{with}J_\mathrm{M}^S>0,$$
(46)
so that the next spin-dependent term to be considered is
$$_{\mathrm{eff}}^{(1)}=J_\mathrm{M}^S\underset{ij}{}𝒬_{\mathrm{M},ij}^{\alpha \beta }\mathrm{S}_i\mathrm{S}_j=J_\mathrm{M}^S\underset{ij}{}\left[\mathrm{I}_{ij}𝒲_{ij}^{\alpha \beta }\right]\mathrm{S}_i\mathrm{S}_j,$$
(47)
and next that it follows from equation (15) that $`\frac{7}{4}\mathrm{I}𝒲^{\alpha \beta }\frac{1}{4}`$. We will investigate below under which conditions an AF net spin exchange on a bond $`ij`$ could actually occur.
For that purpose it is useful, when allowing both $`J_\mathrm{M}^S`$ and $`J_\mathrm{N}^S`$ to be finite, to rewrite the full spin-orbital SE Hamiltonian (34) explicitly as the sum of a purely orbital part and a spin dependent part,
$$_{\mathrm{eff}}=_{\mathrm{eff},\mathrm{o}}+_{\mathrm{eff},\mathrm{s}},$$
(48)
$`_{\mathrm{eff},\mathrm{o}}`$ $`={\displaystyle \underset{i,j}{}}\left(\overline{J}_T^{}_{ij}^{\alpha \beta }+\overline{J}_T𝒥_{ij}^{\alpha \beta }\right)\mathbf{\hspace{0.25em}1}_{ij}=J_\tau {\displaystyle \underset{i,j}{}}𝒲_{ij}^{\alpha \beta }\mathbf{\hspace{0.25em}1}_{ij},`$ (49)
$`_{\mathrm{eff},\mathrm{s}}`$ $`={\displaystyle \underset{i,j}{}}\left(\overline{J}_{TS}^{\prime \prime }\mathrm{I}_{ij}\overline{J}_{TS}^{}_{ij}^{\alpha \beta }\overline{J}_{TS}𝒥_{ij}^{\alpha \beta }\right)\mathrm{S}_i\mathrm{S}_j,`$ (50)
$`={\displaystyle \underset{i,j}{}}\left(J_\sigma \mathrm{I}_{ij}+J_\nu 𝒱_{ij}^{\alpha \beta }J_\mu 𝒲_{ij}^{\alpha \beta }\right)\mathrm{S}_i\mathrm{S}_j,`$
with the exchange constants all positive and satisfying (see A)
$$0<J_\sigma <J_\nu <J_\mu J_\tau .$$
(51)
Note that $`J_\sigma `$ is the pure spin SE, while $`J_\nu `$ is the strength of the SE between a spin-and-orbital operator $`I_i^\alpha \mathrm{S}_i`$ on one site and a pure spin operator $`\mathrm{I}_j\mathrm{S}_j`$ on the other site and $`J_\mu `$ is the SE between spin-and-orbital operators on both sites. So equation (51) shows that the pure spin SE is actually always the weakest interaction (for the above case of $`J_\mathrm{H}=0`$, one has $`\overline{J}_T^{}=J_T`$, $`\overline{J}_T=0`$, $`J_\tau =\frac{1}{2}J_T`$; $`\overline{J}_{TS}^{\prime \prime }=\overline{J}_{TS}^{}=0`$, $`\overline{J}_{TS}=J_{TS}`$, $`J_\sigma =J_\mu =J_\nu =\frac{1}{2}J_{TS}`$). In the first expression in equation (49) we have dropped a constant, while in obtaining the second expression we have used that $`_{i,j}𝒱_{ij}^{\alpha \beta }=0,`$ as follows from equations (10) and (6). Note that the latter simplification is not possible in equation (50), where the spin inner product occurs instead of the unit operator in spin space.
To get an idea of the order of magnitude of the SE coupling constants one would need to estimate the parameters involved. Unfortunately, reliable estimates are available almost exclusively for the divalent TM ions, while only little work has been done on the trivalent ions . Based upon what is available we estimate (all values in eV): $`t1.5`$, $`J_\mathrm{H}1.1`$, $`U_p5`$, $`J_p0.8`$ and $`\mathrm{\Delta }2`$, where the last value is the most uncertain. The small value of $`\mathrm{\Delta }`$ suggests that the covalency effects are quite important in LiNiO<sub>2</sub>, and thus the present model has to be considered only as describing the generic structure of the interactions between states with spin $`S=1/2`$ and $`e_g`$-type orbital degeneracy, which originate from Ni<sup>2+</sup> ions surrounded by a hole shared between the Ni $`3d`$ orbitals and the oxygen $`2p`$ orbitals. It is anyway obvious that $`\mathrm{\Delta }`$ is too small to consider the perturbation theory a controlled expansion. We shall therefore regard the overall energy scale $`Jt^4/\mathrm{\Delta }^3`$ as an unknown parameter and take it as our energy unit. We can then study the relative size of the coupling constants, which are controlled by the ratios $`U_p/\mathrm{\Delta }`$, $`J_\mathrm{H}/\mathrm{\Delta }`$ and $`J_p/U_p`$, which we treat as parameters. Where any of them needed to be fixed for varying other variables, we took $`U_p/\mathrm{\Delta }=2`$, $`J_\mathrm{H}/\mathrm{\Delta }=0.05`$ and $`J_p/U_p=0.1`$ (which is maybe a bit small for the first and especially the second and third one).
The variation of the orbital SE coupling constant $`J_\tau `$ and the three spin and spin-orbital SE constants $`J_\sigma `$, $`J_\nu `$ and $`J_\mu `$ with the strength of the oxygen Coulomb repulsion $`U_p`$ is shown in figure 5. It illustrates once again that oxygen Coulomb interaction is crucial: at $`U_p=0`$ all coupling constants vanish. One further notes that the inequality (51) is well satisfied: indeed $`J_\tau `$ is by far the largest, and the three spin coupling constants are in the expected order. Note also that there is more or less saturation for $`U_p/\mathrm{\Delta }`$ larger than $`2`$ so that the inaccuracy in this parameter is not so important.
Figure 6 shows the variation with $`J_\mathrm{H}`$, the strength of the Hund’s rule exchange on nickel. One observes that a finite $`J_\mathrm{H}`$ enhances the orbital SE considerably \[figure 6(a)\], basically because this lowers the energy of the triplet and so further disfavours (see the remark at the end of section 3) the N configuration, already disfavoured because it allows the largest number of hopping sequences. Moreover, figure 6(b) shows explicitly that nonzero $`J_\mathrm{H}`$ is essential for having unequal spin interaction constants. Since this is equivalent to $`J_\mathrm{M}^S`$ being nonzero (and positive) and thus for having potentially AF interaction, as argued above when the Hamiltonian (47) was discussed, inclusion of the Ni<sup>2+</sup> term splitting makes a qualitative difference for the obtained spin-orbital model. It is amusing that in the unphysical regime of negative $`J_\mathrm{H}`$ the order of the spin SE constants is inverted, so that $`J_\mathrm{M}^S`$ would be negative and AF interactions would not be possible. Note further that magnetic frustration is apparently maximal at $`J_\mathrm{H}=0`$, where the three types of spin dependent interaction are equally strong and so experience the strongest mutual competition.
## 5 Interaction between a pair of Ni<sup>3+</sup> ions
In view of the dominating size of the orbital interaction $`J_\tau `$ it is natural to consider first the orbital interactions alone. Before addressing the Hamiltonian (49) on the entire triangular plane, it is useful to solve first the problem of a single Ni<sup>3+</sup>–Ni<sup>3+</sup> pair, which for definiteness we assume to have its bond in the $`XY`$ plane. For an arbitrary orbital pair state $`|i\theta _i|j\theta _j`$ the orbital energy of the pair is then given by \[compare equation (15)\]
$$E_T^{\mathrm{pair}}(\theta _i,\theta _j)=\frac{1}{4}J_\tau [4\mathrm{cos}^2(\frac{\theta _i+\theta _j}{2})2\mathrm{cos}^2(\frac{\theta _i\theta _j}{2})1],$$
(52)
and the corresponding energy surface is shown in figure 7(a). One notes that similar orbitals ($`\theta _i\theta _j`$) are generally favoured over dissimilar orbitals ($`\theta _i\theta _j\pm \pi `$). More in particular, for a single bond $`_{\mathrm{eff},\mathrm{o}}`$ is seen to favour specific identical orbitals, viz. those with $`\theta _i=\theta _j=\pm \frac{\pi }{2}`$, i.e. a pair configuration with orbitals $`(|z\pm |\overline{z})/\sqrt{2}d_{3z^2r^2}\pm d_{x^2y^2}`$ at both sites, as illustrated in figure 7(b), with energy $`E_T^{\mathrm{pair}}=\frac{3}{4}J_\tau `$.
Frustration prevents this optimum orbital arrangement to be realized on all bonds simultaneously, as already pointed out by MK. As will be seen in the next section, on the full triangular lattice the orbital interaction favours phases in which a pair has either a ferro orbital (FO) arrangement at a general angle $`\theta \theta _i=\theta _j`$ or similarly a canted orbital (CO) arrangement with $`\theta \theta _i=\theta _j`$. It is therefore of interest to consider the effective spin-spin interactions realized under those conditions.
Quite generally, once the orbitals of the pair are fixed, supposedly determined by the pure orbital interaction, the resulting effective spin SE interaction can be written as a Heisenberg Hamiltonian,
$$_{\mathrm{spin}}^{\mathrm{eff}}(\theta _i,\theta _j)=\underset{ij}{}J_{ij}^{\mathrm{eff}}(\theta _i,\theta _j)\mathrm{S}_i\mathrm{S}_j,$$
(53)
with the effective spin SE coupling $`J_{ij}^{\mathrm{eff}}`$ given by
$$J_{ij}^{\mathrm{eff}}(\theta _i,\theta _j)=\frac{1}{4}\left[4J_\sigma J_\mu \left(1+2\mathrm{cos}^2\theta _{}4\mathrm{cos}^2(\theta _++\chi _\gamma )\right)4J_\nu \mathrm{cos}(\theta _++\chi _\gamma )\mathrm{cos}\theta _{}\right],$$
(54)
where $`\theta _\pm =(\theta _i\pm \theta _j)/2`$, as obtained by inserting equations (14) and (15) into (50). For the two relevant cases of orbital order one thus finds, by adopting the appropriate values for the angles $`\theta _i`$ and $`\theta _j`$,
$`J_{\mathrm{FO}}(\theta )J_{ij}^{\mathrm{eff}}(\theta ,\theta )=\frac{1}{4}\left[4J_\sigma 3J_\mu +4J_\mu \mathrm{cos}^2(\theta +\chi _\gamma )4J_\nu \mathrm{cos}(\theta +\chi _\gamma )\right],`$ (55)
$`J_{\mathrm{CO}}(\theta )J_{ij}^{\mathrm{eff}}(\theta ,\theta )=\frac{1}{4}\left[4J_\sigma (14\mathrm{cos}^2\chi _\gamma )J_\mu 2J_\mu \mathrm{cos}^2\theta 4\mathrm{cos}\chi _\gamma J_\nu \mathrm{cos}\theta \right].`$ (56)
The dependence on $`\theta `$ is shown in figure 8 for the values of $`J_\sigma `$, $`J_\mu `$ and $`J_\nu `$ obtained with our standard parameter set. Figure 8(a) shows that in the FO case the spin SE constant is mostly FM (negative), as expected from the Goodenough-Kanamori-Anderson rules, but that AF values, though smaller, are indeed possible, notably for orbitals with $`\theta \pm \pi /3`$. Note that the coupling constant is that for a bond along $`XY`$; equation (55) shows that the curve should be shifted by $`\pm 2\pi /3`$ for bonds along $`YZ`$ or $`ZX`$. The CO case is shown in figure 8(b), for a bond in the $`YZ`$ or $`ZX`$ direction while the angle is still given with respect to the basis {$`|iz`$, $`|i\overline{z}`$} as in equations (1) and (2). Note that the coupling can again be AF as well, precisely for those angles ($`\theta \pm \pi `$) where the FO arrangement gives the largest FM coupling in the $`XY`$ direction.
## 6 Possible ground states
We have already noticed that the orbital SE is strongly frustrated — since the specific orbitals minimizing the pair energy are different for each bond direction, the interactions between neighboring Ni ions cannot be satisfied in all three bond directions simultaneously. The fact that frustration occurs although the interaction is of ferro type is somewhat unusual and is generic for orbital physics as compared to spin physics. The underlying reason is that the orbital interaction Hamiltonian (49) does not have global continuous rotational $`U(1)`$ symmetry with regard to the orbital angles $`\theta _i`$, but is only invariant under rotation by an angle $`\pm 2\pi /3`$, corresponding to a cyclic permutation of the bond directions $`XY`$, $`YZ`$ and $`ZX`$. This characteristic feature and the ensuing frustration are not captured by $`SU(4)`$ symmetric models .
Next we consider what states with long-range order (LRO) are possible. As pointed out by MK, the simplest such state, viz. with uniform ferro orbital order, is actually a minimum energy state in MF theory. Its energy is independent of $`\theta `$ and equal to $`\frac{3}{4}J_\tau `$ per site, while moreover this FO phase may be interrupted without energy cost by lines in which all orbitals have the opposite angle $`\theta `$, i.e. by canting of all orbitals along a line in the triangular lattice. In order to see how this comes about, one may write down the energy for a set of FO lines (for definiteness along the $`XY`$ direction) with orbitals in line $`n`$ given by $`\theta _n`$. From (49) and (15) one obtains
$`E_{\mathrm{FO}\mathrm{l}}=E_{\mathrm{intraline}}+E_{\mathrm{interline}}`$
$`=\frac{1}{4}J_\tau {\displaystyle \underset{n}{}}[2\mathrm{cos}2\theta _n1]`$
$`+\frac{1}{2}J_\tau {\displaystyle \underset{n}{}}\left[\mathrm{cos}(\theta _n+\theta _{n+1}\frac{2\pi }{3})+\mathrm{cos}(\theta _n+\theta _{n+1}+\frac{2\pi }{3})\mathrm{cos}(\theta _n\theta _{n+1})\right]`$
$`=\frac{1}{4}J_\tau {\displaystyle \underset{n}{}}[4\mathrm{cos}^2\theta _n3]J_\tau {\displaystyle \underset{n}{}}\mathrm{cos}\theta _n\mathrm{cos}\theta _{n+1}`$ (57a)
$`=\frac{1}{4}J_\tau {\displaystyle \underset{n}{}}\left[2(\mathrm{cos}\theta _n\mathrm{cos}\theta _{n+1})^23\right].`$ (57b)
Equation (57a) shows that upon variation of the angles $`\theta _n`$ the gain in intraline line energy is compensated by a loss in interline energy, leading on balance to $`E_{\mathrm{FO}\mathrm{l}}`$ being not dependent on the value of $`\theta `$ in the ordered phase, while equation (57b) shows that for any two neighbouring lines the same energy is obtained for $`\theta _n=+\theta _{n+1}`$ and $`\theta _n=\theta _{n+1}`$. So in addition to the degeneracy with respect to $`\theta `$ generated by the frustration, the model also gives rise to randomness due to uncorrelated switches from $`+\theta `$ to $`\theta `$ between successive lines. Such a random CO phase is probably best interpreted as the occurrence of twin boundaries or antiphase boundaries between equivalent FO domains. Because of the zero energy cost for formation of a boundary there is actually no preference for the formation of large domains.
It was argued by MK that the degeneracy in $`\theta `$ is an artefact of MF theory, and that this $`U(1)`$ symmetry of the FO state, not present in the original Hamiltonian (49), would be lifted in the FO domains by the ‘order-out-of-disorder’ mechanism , and that in this way the threefold symmetry associated with the triangular lattice would be restored. However, as the energies involved in the restoring quantum fluctuations are generally small, only a fraction of the orbital coupling constant $`J_\tau `$, one should at this stage reconsider the orbital-spin coupling. Even though the relevant coupling constants $`J_\sigma `$, $`J_\mu `$ and $`J_\nu `$ are much smaller than $`J_\tau `$, the spin-and-orbital dependent SE energy might still be more significant than the fluctuation energies in lifting the $`\theta `$ degeneracy. Figure 8 suggests that it could do so already at the MF level and so determine the ground state before fluctuations need to be considered.
We therefore consider the following LRO spin patterns in conjunction with orbital order: (i) FM order; (ii) the 120 three-sublattice order which is the ground state of the classical TALAF; (iii) AF line order similar to that of the orbitals, i.e. lines of parallel spins along the same direction as the orbital lines, and the spins alternating between up and down on successive lines. The result for the MF energy is shown in figure 9(a). It turns out that indeed the $`\theta `$ degeneracy is lifted, and the lowest MF energy is seen to be obtained for FO order with orbital angle $`\theta =+\pi `$ (or $`\theta =\pi `$), corresponding to a $`d_{x^2y^2}`$ orbital at every site, in combination with spin AF line order, illustrated in figure 9(b). This shows that even though the spin-orbital interaction is considerably weaker than the pure orbital interaction, it can still have a qualitative effect on the nature of the ground state by breaking the $`U(1)`$ symmetry generated by the orbital interactions alone and restoring threefold symmetry. Therefore the large difference in coupling strength does not necessarily lead to a decoupling of the spin and orbital degrees of freedom.
Remarkably, the $`+\theta `$/$`\theta `$ twin-boundary degeneracy persists, not only at the optimum angle $`\pm \pi `$, where it is immaterial since it corresponds only to a sign change of the orbital wavefunctions, but also at arbitrary $`\theta `$. This could in fact have been recognized from the $`\theta `$ dependence of the effective spin exchange constants: the spin exchange in the transverse direction (on the bonds along $`YZ`$ and $`ZX`$) as given for the FO case by shifts over $`\pm \frac{2\pi }{3}`$ in figure 8(a) is the same as that given for the CO case in figure 8(b).
Most significantly, we note that both for the ordered FO phase and for the disordered CO phase the effective spin-spin exchange is antiferromagnetic in the transverse direction, also for orbital angles deviating somewhat from the optimum value $`\pm \pi `$. However, for a disordered phase with truly random orbital angles, either a paraorbital state at high temperature or an orbital glass, the situation would be very different. Figure 8 indicates that orbital randomness would produce a wide distribution of mostly FM exchange constants, although small AF values would still occur but not in a regular pattern as in the FO and CO phases. Obviously a phase transition to an orbital ordered phase would manifest itself also in the magnetic properties, e.g. in the dispersion of spin waves.
Whereas the above analysis shows that the spin-orbital interaction could support orbital LRO, the degeneracy in the orbital sector and the small MF energy of only $`\frac{1}{4}J_\tau `$ per bond, both for the FO and for the CO phase, raise the question whether orbital LRO by itself is stable against formation of a disorded state of valence bond (VB) type. We therefore studied orbital valence bond (OVB) states in the following way. First we solved the eigenvalue problem of the orbital SE Hamiltonian $`_{\mathrm{eff},\mathrm{o}}`$ exactly on the three-site triangular plaquette shown in figure 10(a). The eigenstates are four doublets, with ground state energy $`E_0^{\mathrm{plaq}}=\frac{3}{2}J_\tau `$ and excited states at $`J_\tau `$, $`+\frac{1}{2}J_\tau `$ and $`+2J_\tau `$, the ground state doublet being
$`|e_1=\sqrt{\frac{4}{7}}(|x_1|z_2|y_3+|z_1|y_2|x_3),`$
$`|e_2=\sqrt{\frac{4}{7}}(|\overline{x}_1|\overline{z}_2|\overline{y}_3+|\overline{z}_1|\overline{y}_2|\overline{x}_3).`$ (57bf)
Note that particular combinations of orbitals are favoured, similar to the orbital order for the V triangle in LiVO<sub>2</sub> . It is remarkable that on a triangular lattice such triangular VB correlations in orbital space are favoured, in contrast to dimer singlet correlations in spin systems .
We then constructed an OVB solid by covering the triangular lattice with three-site plaquettes, as in figure 10(b), and assigning to each triangular plaquette $`l`$ an orbital wavefunction built from the groundstate components (57bf),
$$|\mathrm{\Delta }_l=\mathrm{cos}\theta _l|e_1_l+\mathrm{sin}\theta _l|e_2_l,$$
(57bg)
with $`\theta _l`$ still to be determined. This amounts to assigning a fictitious spin to each plaquette and solving a frustrated spin problem on the triangular lattice. In MF theory the total energy is now given by addition of the groundstate energies from all plaquettes and the interplaquette energies coming from the two bonds between nearest neighbour plaquettes,
$$E_{\mathrm{OVB}}=\underset{l}{}E_{\mathrm{\Delta }_l}+\underset{l,m}{}\mathrm{\Delta }_l,\mathrm{\Delta }_m|_{\mathrm{eff},\mathrm{o}}|\mathrm{\Delta }_l,\mathrm{\Delta }_m=\underset{l}{}E_0^{\mathrm{plaq}}+\underset{l,m}{}E_{\mathrm{VB}}^{\mathrm{p}\mathrm{p}}(\theta _l,\theta _m).$$
(57bh)
For a uniform solution the energy per bond is given by
$$\overline{E}_{\mathrm{OVB}}=\frac{1}{6}\left(2\times \frac{1}{3}\times E_0^{\mathrm{plaq}}+4\times \frac{1}{2}\times E_{\mathrm{VB}}^{\mathrm{p}\mathrm{p}}\right),$$
(57bi)
by counting the number of intraplaquette and interplaquette bonds. Although the plaquette eigenenergy amounts to $`\frac{1}{3}E_0^{\mathrm{plaq}}=\frac{1}{2}J_\tau `$ per bond of the plaquette, which for this exact solution is of course lower than the energy per bond for the FO and CO phases, the plaquettes cover only one third of the bonds of the lattice, and so the contribution from the interplaquette energy $`E_{\mathrm{VB}}^{\mathrm{p}\mathrm{p}}`$ is decisive.
The optimum values of the plaquette angles are found to be $`\theta _l=+\frac{\pi }{4}`$, $`\theta _m=\frac{\pi }{4}`$, for which $`E_{\mathrm{VB}}^{\mathrm{p}\mathrm{p}}=\frac{27}{49}J_\tau `$, and if this value could be attained between any pair of plaquettes one would have $`\overline{E}_{\mathrm{OVB}}=(\frac{1}{6}+\frac{9}{49})J_\tau 0.3503J_\tau `$, lower than $`\overline{E}_{\mathrm{FO}}=\overline{E}_{\mathrm{CO}}=0.25J_\tau `$. However, this would requiring alternating angles $`\pm \frac{\pi }{4}`$, which is impossible because of frustration on the triangular lattice. The optimum value that can be obtained is $`\overline{E}_{\mathrm{OVB}}0.228J_\tau `$, realized for instance by line order of $`\pm \frac{\pi }{4}`$. This is remarkably close to the energy per bond of the site-ordered states, and just confirms that the orbital sector is strongly frustrated. It is therefore likely that some quenched disorder would suffice to turn the system into an orbital glass below a freezing temperature $`T_{\mathrm{of}}`$ determined by $`J_\tau `$, as proposed to explain the behaviour of the orbitals observed in LiNiO<sub>2</sub> .
## 7 Origin of the difference between LiNiO<sub>2</sub> and NaNiO<sub>2</sub>
In view of the above results the properties of LiNiO<sub>2</sub> are puzzling , as neither orbital nor spin order sets in down to the lowest temperatures. Thus the high degeneracy of the ground state manifold apparently persists, with Ni<sup>3+</sup> ions being in the low-spin ($`t_{2g}^6e_g^1`$) state, with twofold orbital and twofold spin degeneracy. The standard scenario for such ions would be a cooperative Jahn-Teller effect lifting the orbital degeneracy below a structural transition at $`T_s`$, followed by a magnetic transition at $`T_N`$, lifting also the degeneracy in spin space. Both transitions could be induced already by the electronic mechanism involved in the SE, but the interactions with the lattice are expected to play a significant role in the structural transition, enhancing the value of $`T_s`$, as found in the manganites . In spite of the strong orbital interactions reported above, and although the $`S=1/2`$ spins interact by the predominantly FM SE described in the previous sections, and the magnetic susceptibility follows the Curie-Weiss behaviour, $`\chi (T\theta _{\mathrm{CW}})^1`$ with $`\theta _{\mathrm{CW}}35`$ K down to about 80 K , long-range order in either orbital or spin space is absent . This would suggest that the $`e_g`$ orbital degree of freedom plays an important role, possibly stabilizing strong singlet orbital correlations on individual bonds, and supporting a spin liquid state .
Before addressing the central question of the origin of the spin liquid state in LiNiO<sub>2</sub>, let us summarize shortly the structural and magnetic properties of the similar NaNiO<sub>2</sub> compound, which crystallizes in the same structure as LiNiO<sub>2</sub> (see figure 1). The standard scenario with two subsequent transitions described above is indeed realized in NaNiO<sub>2</sub>. The latter compound undergoes a first-order cooperative Jahn-Teller transition lowering the local symmetry from trigonal to monoclinic at $`T_s=480`$ K , and the Ni-O distances change from $`d=1.98`$ Å to two long bonds with $`d=2.14`$ Å and four short bonds with $`d=1.91`$ Å at $`T<T_s`$. In the low-temperature phase a magnetic transition at $`T_N=20`$ K to an $`A`$-type AF insulator follows — it was first derived from magnetization measurements on a single crystal by Bongers and Enz already in 1966 . This magnetic phase has also been confirmed by a complete static and dynamic magnetic study at low temperatures, where strong anisotropy between the effective magnetic interactions was established . An FM intralayer coupling of $`J_{\mathrm{FM}}=26`$ K and a weak AF interlayer coupling of $`J_{\mathrm{AF}}=2`$ K were found, values consistent with the observed Curie-Weiss temperature $`\mathrm{\Theta }=35`$ K and Néel temperature $`T_N=20`$ K. Similar values of the exchange constants ($`J_{\mathrm{FM}}29`$ K and $`J_{\mathrm{AF}}1.9`$ K) were also deduced recently by analyzing the dispersion of spin waves found in neutron scattering experiments . We remark that the parameters which fit the Curie-Weiss law change around the structural transition at $`T_s`$ , which indicates that the magnetic couplings depend on the actual orbital state, in agreement with the model presented here.
In contrast, in LiNiO<sub>2</sub> the structural transition is absent, but EXAFS experiments indicate the presence of local Jahn-Teller distortions below an orbital freezing temperature $`T_{\mathrm{of}}400`$ K, with two different Ni-O distances again favouring occupied directional $`d_{3z^2r^2}`$-type orbitals : two long bonds with $`d=2.09`$ Å and four short bonds with $`d=1.91`$ Å. While these local distortions are similar to those observed in NaNiO<sub>2</sub>, the absence of a macroscopic distortion in LiNiO<sub>2</sub> has been a mystery since its synthesis . However, very recent evidence from neutron diffraction indicates that below $`T_{\mathrm{of}}`$ the orbitals actually develop short-to-medium-range order in a trimerized state, although the associated strain field prevents the ordering from becoming long-range, and instead nanoscale domains of local orbital trimer order are formed . Nevertheless, even local distortions favouring occupied $`d_{3z^2r^2}`$ orbitals are intriguing, as we have seen above that the SE favours instead alternating symmetric and antisymmetric combinations of the basis orbitals, $`d_{3z^2r^2}\pm d_{x^2y^2}`$, for individual Ni<sup>3+</sup>-Ni<sup>3+</sup> pairs (so e.g. in a disordered state), and FO order of $`d_{x^2y^2}`$ orbitals, as obtained in the MF approximation. Actually, such an ordered phase with occupied $`d_{x^2y^2}`$ orbitals has been observed in large magnetic field , which demonstrates that the state predicted by the SE model is energetically close to the actual ground state, and a quantum phase transition between different orbital states can be triggered by an applied magnetic field. The very fact that instead, in the absence of a magnetic field, local distortions of NiO<sub>6</sub> octahedra associated with occupied $`d_{x^2y^2}`$-type orbitals dominate at low temperature, we interpret as a strong indication that the orbital physics is in the end determined by the Jahn-Teller coupling to the lattice and not by the orbital SE.
The ground state of LiNiO<sub>2</sub> remains the subject of intense debate. Several mechanisms for the spin liquid phase have been proposed so far: ($`i`$) quantum fluctuations on AF bonds , ($`ii`$) geometrical frustration of AF interactions in the $`SU(4)`$ spin-orbital model on the triangular lattice , ($`iii`$) decoupling of spin and orbital degrees of freedom because of the large difference, confirmed above, between orbital and spin-orbital SE coupling constants , ($`iv`$) disorder due to Ni<sup>2+</sup> ions in the Li planes . Interestingly, also in the $`SU(4)`$ model both AF and FM interactions are allowed , so at first sight the mechanism ($`ii`$) might appear similar to the present model of section 4. However, the $`e_g`$ orbital operators do not respect $`SU(2)`$ symmetry, and thus the actual magnetic interactions which follow from the $`SU(4)`$ model are quite different. For instance, in the $`SU(4)`$ model one may assume perfect FO order, with $`\mathrm{T}_i\mathrm{T}_j=\frac{1}{4}`$ for the orbital operators on each bond $`ij`$, which would give AF spin interactions on all the bonds, while such a situation cannot be realized for $`e_g`$ orbitals. Also a state with the orbitals staggered for all bonds is not allowed on a triangular lattice. Therefore, we suggest that the role of geometrical frustration of magnetic interactions is overestimated in the $`SU(4)`$ model and cannot explain the properties of the triangular Ni planes by itself. In contrast, all other mechanisms could contribute and stabilize the spin liquid state in a 3D lattice of LiNiO<sub>2</sub>. We propose a new scenario below in which the above aspects (iii) and (iv) are supplemented by two important features of the orbital glass or trimer nanodomain state below the transition at $`T_{\mathrm{of}}400`$ K: ($`v`$) the weakness of the (effective) SE interactions between Ni<sup>3+</sup> ions in adjacent Ni planes, and ($`vi`$) the Jahn-Teller coupling to the lattice, which induces local distortions favouring directional $`d_{3z^2r^2}`$-like orbitals.
We emphasize that the experimental studies on NaNiO<sub>2</sub> and LiNiO<sub>2</sub> have identified two competing effects which operate in triangular 2D Ni planes: on the one hand local Jahn-Teller distortions favouring occupied $`d_{3z^2r^2}`$ orbitals, produced by the on-site Jahn-Teller coupling, and on the other hand the SE, which would instead stabilize orbital order with occupied $`d_{x^2y^2}`$ orbitals, coexisting with FM interactions within the planes. Although Ni-Ni distances and Ni-O-Ni bond angles in the O–Ni–O trilayers of the two compounds are similar, it may well be that the in-plane elastic couplings and orbital SE differ just sufficiently that the 2D orbital susceptibilities are qualitatively different. In particular, the competition between Jahn-Teller effect and SE may favour orbital disorder or complex multiple-sublattice orbital order when one interaction is only marginally stronger than the other one. This is in fact suggested by the behaviour of LiNiO<sub>2</sub> in magnetic field, mentioned above, and by the properties of the mixed compound Li<sub>0.3</sub>Na<sub>0.7</sub>NiO<sub>2</sub> , discussed below. Whether orbital and magnetic order stabilize in a 3D system, however, depends also crucially on the interplane interactions. Without any coupling between the planes, no long-range order could occur in the 2D planes anyway, in accordance with the Mermin-Wagner theorem.
One thus expects that 3D orbital order can be induced only by a sufficiently strong effective interplane coupling, as apparently encountered in NaNiO<sub>2</sub>, while the absence of orbital order observed in Li<sub>1-x</sub>Ni<sub>1+x</sub>O<sub>2</sub> may be explained by a weaker effective interplane coupling in combination with an already weaker in-plane ordering tendency. Even without an explicit extension of the present microscopic spin-orbital model we can already suggest two mechanisms responsible for this weak interlayer coupling. First of all, the interplane SE is intrinsically weaker in LiNiO<sub>2</sub> than in NaNiO<sub>2</sub> because of the different covalency of the Ni–O–Li–O–Ni bonds as compared to the Ni–O–Na–O–Ni bonds, owing to the weaker hybridization between the oxygen $`2p`$ orbitals and the spatially less extended Li $`2s`$ orbitals than between the oxygen $`2p`$ orbitals and the bigger Na $`3s`$ orbitals. Although this difference is quantitative, it may lead to a qualitatively different behaviour. A further point is then how the individual Ni<sup>3+</sup>-Ni<sup>3+</sup> orbital interactions between two neighbouring planes add up. For two orbital ordered planes, as in NaNiO<sub>2</sub>, the interactions add up coherently, but if the orbitals preferentially avoid ordering in each Ni plane, e.g. because of a tendency towards domain formation, then in spite of the intrinsically identical spin-orbital interaction between Ni<sup>3+</sup> ions in adjacent Ni planes, orbital (and magnetic) couplings add up incoherently. Second, it is well known that in the case of NaNiO<sub>2</sub> an almost ideal structure is obtained, while Li<sub>1-x</sub>Ni<sub>1+x</sub>O<sub>2</sub> is never stoichiometric and extra Ni ions occur in Li planes . This creates frustration due to extrinsic randomness in the individual interplane interactions, and thus further weakens the effective interplane coupling.
Finally, the absence of 3D magnetic order in LiNiO<sub>2</sub> is apparently similarly due to the weakness of the magnetic effective interplane coupling. While it is tempting to ascribe it instead to the presence of both AF and FM SE interactions in the triangular Ni planes, demonstrated above in the spin-orbital model, this is unlikely to be the leading effect stabilizing the spin liquid. In particular, in the orbital glass state the SE would be FM on by far the majority of the bonds in the Ni planes, as follows from equation (54) \[compare also figure 8\], with the precise fraction depending on the degree of disorder, while moreover the few AF interactions would be much weaker than the FM interactions. That the effective interplane magnetic coupling in LiNiO<sub>2</sub> is weak, is not only due to the interplane SE being intrinsically weak owing to the small covalency. In addition, in spite of the identical spin-orbital interaction between Ni<sup>3+</sup> ions in adjacent Ni planes, the absence of orbital long-range order in the Ni planes will make individual Ni<sup>3+</sup>-Ni<sup>3+</sup> spin-spin SE interactions add up incoherently.
An additional extrinsic mechanism opposing magnetic order is provided by magnetic interactions due to the Ni<sup>2+</sup> defects in the Li planes, which by themselves frustrate weak AF SE interactions by stronger FM interactions along Ni<sup>3+</sup>-Ni<sup>2+</sup>-Ni<sup>3+</sup> units , additionally enhanced due to the $`S=1`$ spin states involved. As long as such defects do not occur, coherence in the magnetic interplane coupling is easier to obtain. Our scenario is not inconsistent with the recent observation of the decoupling of orbital and spin degrees of freedom in the mixed compound Li<sub>0.3</sub>Na<sub>0.7</sub>NiO<sub>2</sub> which has a very different orbital state and yet a similar magnetic ground state as NaNiO<sub>2</sub> , with similar exchange constants .
## 8 Discussion and summary
The above analysis demonstrates that the qualitative properties which follow from orbital and spin correlations within triangular Ni planes with 90 Ni–O–Ni bonds can be understood within a realistic spin-orbital SE model. This model demonstrates, in agreement with the SE model of Mostovoy and Khomskii and with experiment , that the orbital SE $`J_\tau `$ is stronger by one order of magnitude than any other (pure spin or spin-orbital) interaction, because all magnetic dependence for the SE along Ni–O–Ni bonds originates from the singlet-triplet splitting of the oxygen $`2p^4`$ configuration, and is therefore smaller by at least a factor $`J_p/U_p0.1`$, where $`U_p`$ ($`J_p`$) is the interorbital Coulomb (Hund’s exchange) interaction on oxygen. These parameters, as well as the effective hopping $`t`$ and the charge transfer energy $`\mathrm{\Delta }`$, are of importance to establish quantitative consequences of the SE model, in particular on such magnetic properties as the effective exchange constants and the magnetic susceptibility.
The orbital SE has interesting consequences for the orbital state and also for the magnetic interactions within the triangular Ni planes. First of all, for an individual (NiO)<sub>2</sub> plaquette (figure 7), a pair of identical orbitals, either both $`d_{3z^2r^2}+d_{x^2y^2}`$ or both $`d_{3z^2r^2}d_{x^2y^2}`$, is favoured. Therefore, in MF approximation a symmetry broken state arises in the Ni planes, with FO order characterized by a single orbital angle $`\theta `$ \[as defined in equation (2)\] along a particular direction. This is in contrast with the behaviour obtained for the cuprates or the manganites , where orbitals with angles $`\theta `$ and $`\theta `$ alternate on interlacing sublattices. As a result, one expects domains with either $`|\theta `$ or $`|\theta `$ orbitals, separated by twin-boundaries. Even though much weaker than the orbital SE, the spin-orbital SE plays a crucial role in the selection of the ground state. It breaks the $`U(1)`$ symmetry generated in MF theory by the orbital SE and restores the original threefold symmetry with respect to the orbital angle $`\theta `$. By this mechanism FO order with planar $`d_{x^2y^2}`$ orbitals in the triangular Ni planes is stabilized in MF theory. A characteristic feature of this state is the presence of both FM and AF spin interactions.
After having a closer look at the magnetic properties of LiNiO<sub>2</sub>, we suggest that this compound represents an interesting case of competition between the type of orbitals favoured locally by the Jahn-Teller effect and those favoured by the SE in an $`e_g`$ system with 90 Ni-O-Ni bonds, and that this is likely to induce orbital disorder. Remarkably, this case is reminiscent of the situation encountered in the $`t_{2g}`$ orbital model , and contradicts the experience and common knowledge from the $`e_g`$ models in perovskites, where the SE along 180 bonds favours alternating orbitals and both mechanisms support each other . Whether this competition might lead to dynamical disorder certainly depends on the phonons, and we conclude that further study of elastic coupling and phonons is urgently needed.
Summarizing, we presented and discussed the consequences of the 2D frustrated quantum spin-orbital model derived for the triangular Ni planes of two compounds with a very different behaviour: LiNiO<sub>2</sub> and NaNiO<sub>2</sub>. This model emphasizes the importance of charge transfer (Goodenough) processes and presents a complete description of the superexchange interactions on the 90 Ni-O-Ni bonds in triangular Ni planes, so we are confident that it provides a solid starting point for future progress in the theory. The frustration in the orbital sector is not of Ising type, as suggested before , but due to different orbitals being favoured for different bond directions. Providing a final answer concerning the origin of the different magnetic properties of LiNiO<sub>2</sub> and NaNiO<sub>2</sub> requires an extension of the present microscopic model by a microscopic description of both the interlayer coupling and of the coupling between orbitals and the lattice. We identified two possible reasons why the interlayer (spin and orbital) coupling is so weak: ($`i`$) an intrinsic effect due to the small covalency of interplane Ni–O–Li–O–Ni bonds, and ($`ii`$) an extrinsic effect due to Ni disorder in Li<sub>1-x</sub>Ni<sub>1+x</sub>O<sub>2</sub> — both these effects are expected to play an important role and to be responsible for the absence of $`A`$-type antiferromagnetic order in LiNiO<sub>2</sub>.
We thank D.I. Khomskii, G.A. Sawatzky and in particular M. Dolores Núñez-Regueiro and Jan Zaanen for stimulating discussions. A.M. Oleś would like to acknowledge support by the Polish State Committee of Scientific Research (KBN) under Project No. 1 P03B 068 26.
## Appendix A Derivation of the spin-orbital Hamiltonian
Here we give further details on the derivation of the spin-orbital Hamiltonian $`_{\mathrm{eff}}`$, in particular regarding the case where $`J_\mathrm{H}`$ is finite. First we take a closer look at the four-step charge transfer sequences. As an example, let us consider (see figure 11) a Ni–Ni pair in the $`XY`$ direction in an M-type initial configuration with the electron on Ni-site $`i`$ ($`j`$) in the non-hopping $`d_{z^2x^2}`$ orbital $`|\overline{y}_i`$ (hopping $`d_{3x^2r^2}`$ orbital $`|x_j`$) having spin up (down), i.e. an initial state $`|\mathrm{M},`$, and the electrons being transferred from the oxygen having opposite spin. In an obvious notation the corresponding contribution to $`_{\mathrm{eff}}`$ is $`\mathrm{M},;𝒫_i^{\overline{y}}𝒫_j^x\mathrm{Q}_{ij}^{}`$. In the middle intermediate state the ion states $`T`$ and $`D`$, $`t`$ and $`s`$, and $`D`$ and $`S`$ occur with equal amplitude at the Ni<sup>2+</sup>, O<sup>0</sup> and Ni<sup>2+</sup> ions, respectively. Denoting the contribution from a doubly excited state by $`[XvY]`$, with $`X,Y\{T,D,S\}`$ and $`v\{t,s\}`$, the contribution corresponding to figure 11 is thus given by
$$\mathrm{M},;=\frac{1}{2}([TtD]+[TtS]+[DtD]+[DtS]+[TsD]+[TsS]+[DsD]+[DsS]),$$
where again a factor 4 has already been included. However, the four sequences that lead to a particular middle state upon inverting the order of the exciting and/or deexciting charge transfers, are now in general inequivalent, as illustrated by the tree diagram of figure 12, because the energies of the first and third intermediate state may depend upon the sequence. The contributions $`[XvY]`$ therefore need to be expanded as
$$[XvY]=\frac{1}{4}\left[(X|XvY|X)+(X|XvY|Y)+(Y|XvY|X)+(Y|XvY|Y)\right],$$
(57bj)
where we denote in $`(Z_1|XvY|Z_3)`$ by $`Z_1`$ ($`Z_3`$) the state reached after the first (third) charge transfer step.
The explicit expression for this quantity is \[compare equation (30)\]
$`(Z_1|XvY|Z_3)={\displaystyle \frac{t^4}{\mathrm{\Delta }_{Z_1}\mathrm{\Delta }_{Z_3}(\mathrm{\Delta }_X+\mathrm{\Delta }_Y+U_v)}}+{\displaystyle \frac{t^4}{\mathrm{\Delta }_{Z_1}\mathrm{\Delta }_{Z_3}(\mathrm{\Delta }_X+\mathrm{\Delta }_Y)}}`$
$`={\displaystyle \frac{t^4U_v}{\mathrm{\Delta }_{Z_1}\mathrm{\Delta }_{Z_3}(\mathrm{\Delta }_X+\mathrm{\Delta }_Y)(\mathrm{\Delta }_X+\mathrm{\Delta }_Y+U_v)}},`$ (57bk)
where the second term in the first line is the renormalization correction discussed in section 3. Inserting (57bk) into equation (57bj) one now obtains
$$[XvY]=\frac{t^4}{4}\frac{\mathrm{\Delta }_X+\mathrm{\Delta }_Y}{\mathrm{\Delta }_X^2\mathrm{\Delta }_Y^2}\frac{U_v}{\mathrm{\Delta }_X+\mathrm{\Delta }_Y+U_v},$$
(57bl)
so that for example
$`[TsS]={\displaystyle \frac{t^4\mathrm{\Delta }}{2(\mathrm{\Delta }^24J_\mathrm{H}^2)^2}}{\displaystyle \frac{U_p+J_p}{2\mathrm{\Delta }+U_p+J_p}},`$ (57bm)
$`[DtD]={\displaystyle \frac{t^4}{2\mathrm{\Delta }^3}}{\displaystyle \frac{U_pJ_p}{2\mathrm{\Delta }+U_pJ_p}}.`$ (57bn)
The next step consists in listing all possible initial configurations and for each one the allowed spin arrangements of the two hopping electrons, and then working out, as above for $`\mathrm{M},;`$, with the help of diagrams like those shown in figure 11, which middle intermediate states contribute. The result is shown in table 1, which gives the coefficients $`C(..|..)`$ in
$$\mathrm{L},\sigma \sigma ^{};\overline{\sigma }\overline{\sigma }^{}=\underset{X,Y,v}{}C(\mathrm{L},\sigma \sigma ^{};\overline{\sigma }\overline{\sigma }^{}|XvY)[XvY].$$
(57bo)
The SE coupling constants $`\{K_\mathrm{L}^{\mathrm{T},\mathrm{S}}\}`$ for spin triplet and spin singlet, appearing in equation (33), are now obtained as follows. In each case (L = O, M or N) consider first the $`m=1`$ component of the triplet, i.e. the row(s) in table 1 in which there is an up-spin on both Ni<sup>3+</sup> ions. Adding them gives the triplet contribution $`K_\mathrm{L}^\mathrm{T}`$, while similarly adding the row(s) for which the two Ni<sup>3+</sup> spins are in up-down configuration gives $`K_\mathrm{L}^{}`$,
$`K_\mathrm{L}^\mathrm{T}={\displaystyle \underset{\overline{\sigma },\overline{\sigma }^{}}{}}\mathrm{L},;\overline{\sigma }\overline{\sigma }^{}={\displaystyle \underset{\overline{\sigma },\overline{\sigma }^{}}{}}{\displaystyle \underset{X,Y,v}{}}C(\mathrm{L},;\overline{\sigma }\overline{\sigma }^{}|XvY)[XvY],`$ (57bp)
$`K_\mathrm{L}^{}={\displaystyle \underset{\overline{\sigma },\overline{\sigma }^{}}{}}\mathrm{L},;\overline{\sigma }\overline{\sigma }^{}={\displaystyle \underset{\overline{\sigma },\overline{\sigma }^{}}{}}{\displaystyle \underset{X,Y,v}{}}C(\mathrm{L},;\overline{\sigma }\overline{\sigma }^{}|XvY)[XvY].`$ (57bq)
As $`K_\mathrm{L}^{}`$ corresponds to the $`m=0`$ components and so is shared between singlet and triplet, one now obtains the singlet contribution from $`K_\mathrm{L}^\mathrm{S}=2K_\mathrm{L}^{}K_\mathrm{L}^\mathrm{T}`$.
The result is
$`K_\mathrm{O}^\mathrm{T}=[DtD]+2[DtS]+[StS],`$ (57br)
$`K_\mathrm{O}^\mathrm{S}=[DsD]+2[DsS]+[SsS],`$ (57bs)
$`K_\mathrm{M}^\mathrm{T}=2[TtD]+[TsD]+2[TtS]+[TsS]+[DtD]+[DtS],`$ (57bt)
$`K_\mathrm{M}^\mathrm{S}=3[TtD]+3[TtS]+[DsD]+[DsS],`$ (57bu)
$`K_\mathrm{N}^\mathrm{T}=7[TtT]+2[TsT]+4[TtD]+2[TsD]+[DtD],`$ (57bv)
$`K_\mathrm{N}^\mathrm{S}=6[TtT]+3[TsT]+6[TtD]+[DsD].`$ (57bw)
The above procedure (considering first the $`m=1`$ component of the triplet) avoids the necessity of keeping track of the phases. Alternatively one may work out explicitly the states resulting after two perturbation steps from both $`|\mathrm{L},`$ and $`|\mathrm{L},`$, and from those, by taking their sum and difference, the states resulting from the $`m=0`$ components of the triplet and the singlet, and finally project the latter on all possible middle states. It is straightforward to verify that this gives the same results and equations (57br)–(57bw) are reproduced. One should note that, while the parallel-spin initial state $`|\mathrm{L},`$ corresponds necessarily to the triplet and is therefore associated with FM coupling, it is not correct to associate the antiparallel-spin initial state $`|\mathrm{L},`$ with AF coupling, since it projects on both the triplet and the singlet.
The SE constants $`J_\mathrm{L}^{0,S}`$ occurring in equation (34) then follow from equations (35) and (57br)–(57bw), with the result
$`J_\mathrm{O}^0=(DD)_++2(DS)_++(SS)_+,`$ (57bx)
$`J_\mathrm{M}^0=3(TD)_++3(TS)_++(DD)_++(DS)_+,`$ (57by)
$`J_\mathrm{N}^0=9(TT)_++6(TD)_++(DD)_+,`$ (57bz)
$`J_\mathrm{O}^S=(DD)_{}+2(DS)_{}+(SS)_{},`$ (57ca)
$`J_\mathrm{M}^S=(TD)_{}(TS)_{}+(DD)_{}+(DS)_{},`$ (57cb)
$`J_\mathrm{N}^S=(TT)_{}2(TD)_{}+(DD)_{},`$ (57cc)
where we have introduced the following abbreviations for the orbital exchange elements occurring in (57bx)–(57bz) and the spin-orbital exchange elements occurring in (57ca)–(57cc),
$`(XY)_+`$ $`=`$ $`\frac{3}{4}[XtY]+\frac{1}{4}[XsY],`$ (57cd)
$`(XY)_{}`$ $`=`$ $`[XtY][XsY].`$ (57ce)
One readily verifies that equations (57bx)–(57cc) reduce to equations (41)–(44) upon setting $`J_\mathrm{H}=0`$, whereupon all $`(XY)_+`$ become identical as do all $`(XY)_{}`$.
From equations (57cd)–(57ce) or from the explicit expressions
$`(XY)_+`$ $`=`$ $`{\displaystyle \frac{t^4}{4}}{\displaystyle \frac{\mathrm{\Delta }_X+\mathrm{\Delta }_Y}{\mathrm{\Delta }_X^2\mathrm{\Delta }_Y^2}}{\displaystyle \frac{(U_p\frac{1}{2}J_p)(\mathrm{\Delta }_X+\mathrm{\Delta }_Y)+U_p^2J_p^2}{(\mathrm{\Delta }_X+\mathrm{\Delta }_Y+U_p)^2J_p^2}},`$ (57cf)
$`(XY)_{}`$ $`=`$ $`{\displaystyle \frac{t^4}{2}}{\displaystyle \frac{(\mathrm{\Delta }_X+\mathrm{\Delta }_Y)^2}{\mathrm{\Delta }_X^2\mathrm{\Delta }_Y^2}}{\displaystyle \frac{J_p}{(\mathrm{\Delta }_X+\mathrm{\Delta }_Y+U_p)^2J_p^2}},`$ (57cg)
it is obvious that the orbital exchange elements $`\{(XY)_+\}`$ are much larger than the spin-orbital exchange elements $`\{(XY)_{}\}`$, the ratio being
$$\frac{|(XY)_{}|}{(XY)_+}=\frac{2J_p(\mathrm{\Delta }_X+\mathrm{\Delta }_Y)}{(U_p\frac{1}{2}J_p)(\mathrm{\Delta }_X+\mathrm{\Delta }_Y)+U_p^2J_p^2}\frac{4\mathrm{\Delta }J_p}{(2\mathrm{\Delta }+U_p)U_p},$$
(57ch)
where the second expression applies for $`J_p,J_\mathrm{H}\mathrm{\Delta },U_p`$. It then follows from equations (57bx)–(57bz) and (57ca)–(57cc) that the same holds for the orbital SE constant $`J_\tau `$ \[or $`\overline{J}_T^{}`$, see equation (49)\] when compared with any of the spin or spin-orbital SE constants $`J_\sigma `$, $`J_\nu `$ or $`J_\mu `$ \[or $`\overline{J}_{TS}`$, $`\overline{J}_{TS}^{}`$, or $`\overline{J}_{TS}^{\prime \prime }`$, see equation (50)\] — the same ratio (57ch) applies to $`\overline{J}_{TS}/\overline{J}_T^{}`$, in accordance with the analysis made above for $`J_\mathrm{H}=0`$ \[see equations (43)–(44)\].
As for the relative size of the spin and spin-orbital constants with respect to one another, we observe that while $`J_\mathrm{O}^S`$ is the sum of four spin-orbital exchange elements \[equation (57ca)\], $`J_\mathrm{M}^S`$ is the sum of two differences of such elements, while $`J_\mathrm{N}^S`$ is the difference of two differences of such elements,
$$J_\mathrm{M}^S=D_S+D_D,J_\mathrm{N}^S=D_DD_T,$$
(57ci)
with
$$D_X=(DX)_{}(TX)_{}\left(\frac{1}{\mathrm{\Delta }_T}\frac{1}{\mathrm{\Delta }_D}\right)(\frac{1}{\mathrm{\Delta }_T}+\frac{1}{\mathrm{\Delta }_D}+\frac{2}{\mathrm{\Delta }_X})C_{TDX},$$
(57cj)
where $`C_{TDX}`$ depends only weakly upon $`J_\mathrm{H}`$ and is given in good approximation by
$$C_{TDX}\frac{t^4}{2}\frac{J_p}{(2\mathrm{\Delta }+U_p)^2J_p^2}.$$
(57ck)
Obviously all $`D_X`$ are positive since $`\mathrm{\Delta }_T<\mathrm{\Delta }_D`$, and further $`D_T>D_D>D_S`$ since $`\mathrm{\Delta }_T<\mathrm{\Delta }_D<\mathrm{\Delta }_S`$. It follows that $`J_\mathrm{M}^S>0`$ and $`J_\mathrm{N}^S<0`$, and moreover, as clearly $`D_S+2D_D`$ is considerably larger than $`D_T`$, one concludes that
$$|J_\mathrm{O}^S|J_\mathrm{M}^S|J_\mathrm{N}^S|.$$
(57cl)
Finally we arrive at the SE constants associated with the more physical interactions as occurring in $`_{\mathrm{eff},\mathrm{o}}`$ \[equation (49)\] and $`_{\mathrm{eff},\mathrm{s}}`$ \[equation (50)\]. The orbital part is described by
$$J_\tau =\frac{1}{2}(J_\mathrm{O}^02J_\mathrm{M}^0+J_\mathrm{N}^0)=\frac{1}{2}[9(TT)_+6(TS)_++(SS)_+],$$
(57cm)
while the pure spin SE constant $`J_\sigma `$ and the two spin-orbital SE constants $`J_\nu `$ and $`J_\mu `$ are given by
$`J_\sigma ={\displaystyle \frac{1}{2}}(J_\mathrm{O}^S+2J_\mathrm{M}^S+J_\mathrm{N}^S)`$
$`={\displaystyle \frac{1}{2}}[(TT)_{}4(TD)_{}2(TS)_{}+4(DD)_{}+4(DS)_{}+(SS)_{}],`$ (57cn)
$`J_\mu ={\displaystyle \frac{1}{2}}(J_\mathrm{O}^S2J_\mathrm{M}^S+J_\mathrm{N}^S)`$
$`={\displaystyle \frac{1}{2}}[(TT)_{}+2(TS)_{}+(SS)_{}]`$
$`=J_\sigma +2D_S+2D_D,`$ (57co)
$`J_\nu ={\displaystyle \frac{1}{2}}(J_\mathrm{O}^SJ_\mathrm{N}^S)`$
$`={\displaystyle \frac{1}{2}}[(TT)_{}+2(TD)_{}+2(DS)_{}+(SS)_{}]`$
$`=J_\sigma +D_S+2D_DD_T.`$ (57cp)
From these expressions and the inequality (57cl) above, the inequality (51) follows immediately.
## References |
warning/0507/cond-mat0507723.html | ar5iv | text | # Nodal Quasiparticle Lifetimes in Cuprate Superconductors
## I Introduction
A $`d`$-wave BCS framework has proved useful in describing the thermal Zha01 ; Chi99 and microwave Hos99 conductivities of the superconducting cuprates HPS94 ; WS00 ; DL00 . This suggests that a quasiparticle description of the nodal excitations is adequate, at least at low energies. Here, the frequencies are small compared with the gap magnitude and the temperature dependence of the nodal quasiparticle lifetime plays the dominant role. Recent angle-resolved photoemission spectroscopy (ARPES) experiments on Bi<sub>2</sub>Sr<sub>2</sub>Ca<sub>1</sub>Cu<sub>2</sub>O<sub>8</sub> (BSCCO) are providing high resolution data on the momentum and energy dependence of the nodal quasiparticle lifetimes as well as their temperature dependence in this low energy range Despc ; Val00 ; Kam00 ; Kor04 . Thus, the question arises whether transport and ARPES data can be understood within a single framework. Establishing this agreement is important to confirming the nature of the superconducting phase as a BCS-like $`d`$-wave state with nodal quasiparticle excitations.
The width $`\mathrm{\Delta }k(\omega ,T)`$ of the momentum distribution curve (MDC) measured by ARPES experiments is proportional to the inverse of the on-shell quasiparticle lifetime. Along a nodal cut in momentum, with $`k_x=k_y`$, the magnitude $`k=\sqrt{k_x^2+k_y^2}`$ is set by $`\omega =\epsilon _k`$, i.e. $`k`$ becomes effectively a function of $`\omega `$; we refer to this as “on-shell”. Here, $`\epsilon _k`$ is the quasiparticle energy. Thus, the lifetime of such a nodal quasiparticle depends upon $`\omega `$ and the temperature $`T`$. As $`k`$ decreases from the nodal Fermi momentum $`k_N`$, $`\omega `$ can be swept over an energy range which is significantly larger than $`\mathrm{\Delta }_0`$. Therefore, there is a need to extend the inelastic scattering lifetime calculations to cover a wider range of energies. In addition, impurity scattering, particularly due to out-of-plane forward scattering, is believed to be important so that it is also of interest to determine its temperature and energy dependence AV00 ; ZHS04 .
In the following, we will find it convenient to discuss the elastic and inelastic quasiparticle scattering in terms of a scattering rate
$$\mathrm{\Gamma }(\omega ,𝐤,T)=\mathrm{\Sigma }^{\prime \prime }(\omega ,𝐤,T),$$
(1)
where $`\mathrm{\Sigma }^{\prime \prime }`$ is the imaginary part of the quasiparticle self-energy. The inverse of the quasiparticle lifetime $`1/\tau (\omega ,𝐤,T)=2\mathrm{\Gamma }(\omega ,𝐤,T)`$ and the width of the MDC about a given $`𝐤`$ determined from $`\epsilon _k=\omega `$ is
$$\mathrm{\Delta }k(\omega ,T)=\frac{2\mathrm{\Gamma }(\omega ,𝐤|_{\omega =\epsilon _𝐤},T)}{v(𝐤)},$$
(2)
with $`v(𝐤)=\frac{\epsilon _𝐤}{𝐤}`$ the bare band velocity. Here, we report on results obtained within a $`d`$-wave BCS framework for the energy and temperature dependence of the nodal quasiparticle elastic and inelastic scattering rates and discuss the $`\omega `$ and $`T`$ dependence of $`\mathrm{\Gamma }(\omega ,𝐤|_{\omega =\epsilon _k},T)`$.
We begin in Section II with a discussion of the temperature and energy dependence of the elastic impurity scattering contribution. In BSCCO, there are both in-plane and out-of-plane scattering centers Eis04 . The in-plane impurities are believed to give rise to strong scattering and are often treated in the unitary limit Hud03 . On the other hand, the out-of-plane impurity scattering is weak and tends to be forward. It will be treated within a self-consistent Born approximation. We will see that the temperature dependence of the elastic scattering rate arises from the opening of the gap as $`T`$ decreases below $`T_c`$, while its energy dependence reflects the decreased phase space as $`\omega `$ becomes smaller than the gap magnitude. The low energy inelastic scattering discussed in Section III is dominated by short range Coulomb scattering. At low energies, the linear $`\omega `$ dependence of the density of states gives rise to an inelastic scattering rate which varies as $`T^3`$ for $`T`$ larger than $`\omega `$ or $`\omega ^3`$ for $`\omega `$ larger than $`T`$. At higher energies, one can probe the $`\omega `$-dependence of the effective interaction. Here, we confine ourselves to the Hubbard model and explore the spin-fluctuation contributions to the higher energy behavior of the inelastic scattering. In Section IV, we discuss the combined effects of the elastic and inelastic scattering processes, and Section V contains our conclusions.
## II Elastic Scattering
Elastic scattering in BSCCO can arise from both impurities and disorder in the CuO planes as well as from regions outside these planes. The scattering rate due to unitary scatterers (possibly Cu vacancies) with concentration roughly $`n_u0.2\%`$ observed as zero-bias resonances in scanning tunneling microscopy (STM) experiments Hud03 is well understood theoretically. It will be treated as usual in the self-consistent $`T`$-matrix approximation,
$$\underset{¯}{\mathrm{\Sigma }}_{el,u}=\frac{n_u}{_𝐤G(𝐤,\omega )}\tau _0,$$
(3)
where $`\tau _i`$ are the $`2\times 2`$ Pauli matrices in Nambu notation. At low frequencies and temperatures, the unitary scatterers give rise to an impurity resonance of height $`\mathrm{\Sigma }_{el,u}^{\prime \prime }(\omega =0)\gamma _u\sqrt{\mathrm{\Gamma }_u\mathrm{\Delta }_0}10^2t`$, and roughly the same width, with $`t`$ the nearest-neighbor hopping matrix element. In the normal state we find a scattering rate of $`\mathrm{\Gamma }_un_uE_F10^3t`$ due to these defects, leading to $`\gamma _u`$ 1 to 2 meV. If, as we believe, the normal state elastic scattering rate at the node is a significant fraction of the gap, the contribution of the unitary scatterers will be difficult to resolve in ARPES experiments on BSCCO.
The elastic scattering from out-of-plane impurities, such as interstitial O ions, can be modeled by a momentum-dependent potential,
$$V(𝐤,𝐤^{})=V_0f(𝐤,𝐤^{})$$
(4)
which has a forward scattering form factor $`f(𝐤,𝐤^{})`$ which cuts off the scattering when $`|𝐤𝐤^{}|`$ exceeds a momentum $`\kappa `$, as has been discussed in Ref. ZHS04, . We will measure $`\kappa `$ in units of $`k_F`$ so that $`\kappa ^1`$ characterizes the range of the impurity potential in units proportional to the Cu-Cu spacing $`a`$. Since the scattering from an out-of-plane impurity is relatively weak, it can be treated within the self-consistent Born approximation Sch05 . For a nodal quasiparticle with $`k_x=k_y`$, the gap vanishes and the scattering rate due to elastic scattering is determined from the imaginary part of the $`\tau _0`$ Nambu self-energy in the self-consistent Born approximation,
$$\mathrm{\Sigma }_{el,0}(𝐤,\omega )=n_i\underset{k^{}}{}|V(𝐤,𝐤^{})|^2\frac{\stackrel{~}{\omega }}{\stackrel{~}{\omega }^2\epsilon _k^{}^2\stackrel{~}{\mathrm{\Delta }}_k^{}^2},$$
(5)
where $`n_i`$ is the planar density of the out-of-plane impurities. In addition, because the potential $`V(𝐤,𝐤^{})`$ is anisotropic, it is necessary to renormalize the order parameter as well:
$$\mathrm{\Sigma }_{el,1}(𝐤,\omega )=n_i\underset{k^{}}{}|V(𝐤,𝐤^{})|^2\frac{\stackrel{~}{\mathrm{\Delta }}_𝐤^{}}{\stackrel{~}{\omega }^2\epsilon _k^{}^2\stackrel{~}{\mathrm{\Delta }}_k^{}^2},$$
(6)
where $`\stackrel{~}{\omega }=\omega \mathrm{\Sigma }_{el,0}`$ and $`\stackrel{~}{\mathrm{\Delta }}_𝐤=\mathrm{\Delta }_𝐤+\mathrm{\Sigma }_{el,1}`$. The renormalization of $`\epsilon _𝐤`$ by disorder is also nonzero in general but was shown to be negligible in Ref. ZHS04, , and will therefore be neglected here. For our numerical calculations we will use a simple parameterization of the band structure with
$$\epsilon _k=2t(\mathrm{cos}k_x+\mathrm{cos}k_y)4t^{}\mathrm{cos}k_x\mathrm{cos}k_y\mu $$
(7)
and
$$\mathrm{\Delta }_k=\frac{\mathrm{\Delta }_0}{2}(\mathrm{cos}k_x\mathrm{cos}k_y).$$
(8)
Here, $`t^{}=0.35`$ and $`\mu =1.1`$, with energy measured in units of nearest neighbor hopping $`t`$. For the temperature dependence of the gap we use a common interpolation form
$$\mathrm{\Delta }_0(T)=\mathrm{\Delta }_0\mathrm{tanh}\left(\alpha \sqrt{\frac{T_c}{T}1}\right)$$
(9)
with $`\alpha =3`$, $`2\mathrm{\Delta }_0/kT_c=6`$, and $`\mathrm{\Delta }_0=0.2`$. These parameters are chosen to give the observed magnitude of the $`T=0`$ gap and a somewhat more rapid rise at $`T_c`$, consistent with experiment.
The scattering rate $`\mathrm{\Gamma }_{el}(𝐤_N,\omega )`$ is now determined by the imaginary part of the solution of Eq. (5) at the nodal wavevector $`𝐤_N`$. Before turning to a numerical evaluation of Eq. (5)-(6), we consider a simple analytic approximation similar in spirit to that discussed in Kee01 in the context of $`T_c`$ suppression by disorder Tcsupress . Measuring the momentum transfer from $`𝐤_N`$ along the $`k_{}`$ and $`k_{}`$ coordinates shown in Fig. 1, the imaginary part of Eq. (5) can be written as
$$\mathrm{\Gamma }_{el}=n_i|V_0|^2\mathrm{Im}\frac{dk_{}^{}d𝐤_{}^{}}{(2\pi )^2}|f(𝐤,𝐤^{})|^2\frac{\stackrel{~}{\omega }}{\stackrel{~}{\omega }^2v_1^2k_{}^2\stackrel{~}{v}_2^2k_{}^2}$$
(10)
Here, $`v_1=v_F`$ at the nodal point and $`\stackrel{~}{v}_2`$ is the renormalized gap velocity. We have assumed that the $`\kappa `$ cut-off prevents scattering to other nodes.
Since $`v_1v_2`$, the $`\kappa `$ cut-off primarily affects the $`k_{}`$ integral and one can integrate freely over the important range $`k_{}`$ giving
$$\mathrm{\Gamma }_{el}=\frac{1}{2}n_iN_0|V_0|^2\mathrm{Re}_\kappa ^\kappa 𝑑k_{}^{}\frac{\stackrel{~}{\omega }}{\sqrt{\stackrel{~}{\omega }^2\stackrel{~}{v}_2^2k_{}^2}}$$
(11)
with $`N_0=(2\pi v_F)^1`$ the band density of states in units in which $`k_F=1`$. Setting $`\stackrel{~}{\omega }=\omega +i\mathrm{\Gamma }_{el}`$, i.e. neglecting the real part of the self-energy, this becomes
$$\mathrm{\Gamma }_{el}=\frac{\mathrm{\Gamma }_0}{\stackrel{~}{v}_2\kappa }\mathrm{Re}\left[(\omega +i\mathrm{\Gamma }_{el})\mathrm{sin}^1\left(\frac{\stackrel{~}{v}_2\kappa }{\omega +i\mathrm{\Gamma }_{el}}\right)\right],$$
(12)
where $`\mathrm{\Gamma }_0=n_iN_0|V_0|^2\kappa `$ is the normal state elastic scattering rate and $`\stackrel{~}{v}_2\kappa `$ is the maximum renormalized order parameter probed by the scattering. From this result we recover immediately the $`\omega 0`$, $`k_{}0`$ limit that $`\mathrm{\Gamma }_{el}\mathrm{\Gamma }_0`$ from Ref. ZHS04, .
The gap velocity at the node in the pure system, $`v_22\mathrm{\Delta }_0/k_N`$, is significantly renormalized in the presence of disorder to $`\stackrel{~}{v}_2`$. From Eq. (6) we find, for $`k_{}\kappa `$,
$`(\stackrel{~}{v}_2v_2)k_{}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }_0}{2\kappa }}\mathrm{Im}{\displaystyle _\kappa ^\kappa }𝑑k_{}^{}{\displaystyle \frac{\stackrel{~}{v}_2(k_{}^{}+k_{})}{\sqrt{\stackrel{~}{\omega }^2\stackrel{~}{v}_2^2(k_{}^{}+k_{})^2}}}`$ (13)
leading to
$`\stackrel{~}{v}_2v_2`$ $``$ $`\mathrm{\Gamma }_0{\displaystyle \frac{\stackrel{~}{v}_2}{\sqrt{\mathrm{\Gamma }_{el}^2+\stackrel{~}{v}_2^2\kappa ^2}}}.`$ (14)
This result, which should be solved self-consistently together with Eq. (12), obtains only in the limit where the scattering is sufficiently forward such that the momentum integration takes place only near the node; it interpolates correctly to the isotropic limit $`\kappa \mathrm{}`$, however, where the gap velocity renormalization vanishes, $`\stackrel{~}{v}_2=v_2`$. In the forward scattering limit $`\kappa 0`$, it is easy to check that $`\stackrel{~}{v}_2v_2/2`$, and the dependence for small $`\kappa `$ is otherwise weak, as shown in Fig. 2.
We will consider two different ways of characterizing the scattering rates of quasiparticles with $`k_x=k_y`$. First, suppose $`𝐤`$ is fixed exactly at the nodal Fermi surface $`𝐤_N`$. In this case the quasiparticle energy $`\omega =0`$ and one can study $`\mathrm{\Gamma }_{el}(𝐤_N,\omega =0)`$ as a function of temperature. The contribution of the out-of-plane elastic impurity scattering to the temperature dependence of the MDC line width $`\mathrm{\Delta }k(T)`$ for a quasiparticle with $`𝐤=𝐤_N`$ is given by $`2\mathrm{\Gamma }_{el}(𝐤_N,\omega =0)/v_F`$. When $`\omega =0`$, Eq. (12) can be written as
$$1=\frac{\mathrm{\Gamma }_0}{\stackrel{~}{v}_2\kappa }\mathrm{ln}\left(\frac{1+\sqrt{1+(\mathrm{\Gamma }_{el}/\stackrel{~}{v}_2\kappa )^2}}{(\mathrm{\Gamma }_{el}/\stackrel{~}{v}_2\kappa )}\right).$$
(15)
At low temperatures, if $`\stackrel{~}{v}_2\kappa /\mathrm{\Gamma }_01`$, we find the familiar Born result
$$\mathrm{\Gamma }_{\mathrm{el}}(T)\stackrel{~}{v}_2\kappa e^{\left(\frac{\stackrel{~}{v}_2\kappa }{\mathrm{\Gamma }_0}\right)}$$
(16)
with the gap maximum $`\mathrm{\Delta }_0`$ in the isotropic case replaced by $`\stackrel{~}{v}_2\kappa `$. When $`\stackrel{~}{v}_2\kappa /\mathrm{\Gamma }_01`$,
$$\mathrm{\Gamma }_{\mathrm{el}}(T)\mathrm{\Gamma }_0\left(1\frac{1}{6}\left(\frac{\stackrel{~}{v}_2\kappa }{\mathrm{\Gamma }_0}\right)^2\right).$$
(17)
This limit, of course, always applies when $`T`$ goes to $`T_c`$, but it can also apply for all values of $`T/T_c`$ if the system is sufficiently dirty. Note that since $`\mathrm{\Gamma }_0`$ varies as $`|V_0|^2\kappa `$, if $`V_0`$ is kept fixed, the normal state scattering rate decreases as $`\kappa `$ becomes smaller. This simply reflects the reduction of phase space as the scattering is restricted to be in a more narrow forward cone. At the same time, the drop in the scattering rate as the gap $`\mathrm{\Delta }_0(T)`$ opens for $`T<T_c`$, is smaller for reduced values of $`\kappa `$. This reflects the fact that the maximum effective gap reached by a scattered nodal quasiparticle is reduced by $`\kappa `$ so that the opening of the gap is less effective in suppressing the scattering rate. Results obtained from solving Eqs. (12) and (14) for $`\mathrm{\Gamma }_0/\mathrm{\Delta }_0=0.5`$ and 0.25 are shown in Fig. 3 a) and c) as a function of $`\kappa `$. The temperature dependence of the elastic scattering arises from the temperature dependence of the $`d`$-wave gap, and therefore varies as $`T^3`$ at low $`T`$. Note that the interpolation form, Eq. (9), actually gives a qualitatively incorrect result at low temperatures for a $`d`$-wave superconductor. Eq. (9) implies an activated thermal depletion of the condensate, whereas the correct solution to the $`d`$-wave gap equation has a leading term varying as $`T^3`$ dwavegapTdependence , giving $`\mathrm{\Delta }_0(T)\mathrm{\Delta }_0(1bT^3/\mathrm{\Delta }_0^3)`$, where $`b`$ is a constant of order unity. In the clean limit, it follows from Eq. (16) that
$`\mathrm{\Gamma }_{\mathrm{el}}(T)`$ $``$ $`\mathrm{\Gamma }_{\mathrm{el}}(T=0)\left(1+{\displaystyle \frac{b(\stackrel{~}{v}_2\kappa \mathrm{\Gamma }_0)}{\mathrm{\Delta }_0^3\mathrm{\Gamma }_0}}T^3\right)`$ (18)
$``$ $`\mathrm{\Gamma }_{\mathrm{el}}(T=0)\left(1+2{\displaystyle \frac{b\kappa T^3}{\mathrm{\Gamma }_0\mathrm{\Delta }_0^2}}\right)`$
while in the dirty limit, Eq. (17) implies that
$`\mathrm{\Gamma }_{\mathrm{el}}(T)`$ $``$ $`\mathrm{\Gamma }_0\left(1{\displaystyle \frac{(\stackrel{~}{v}_2\kappa )^2}{6\mathrm{\Gamma }_0^2}}\right)+{\displaystyle \frac{b(\stackrel{~}{v}_2\kappa )^2}{3\mathrm{\Delta }_0^3\mathrm{\Gamma }_0}}T^3`$ (19)
$``$ $`\mathrm{\Gamma }_0\left(1{\displaystyle \frac{\mathrm{\Delta }_0^2\kappa ^2}{6\mathrm{\Gamma }_0^2}}\right)+{\displaystyle \frac{b\kappa ^2}{3\mathrm{\Delta }_0\mathrm{\Gamma }_0^2}}T^3.`$
In Eqs. (18) and (19) we have assumed $`v_22\mathrm{\Delta }_0`$ and used the asymptotic results for $`\stackrel{~}{v}_2/v_2`$ derived above to obtain simple expressions for the clean and dirty limits, respectively.
It is also interesting to study the $`\omega `$ energy dependence of $`\mathrm{\Gamma }_{el}`$ at a fixed temperature. Within the current approximate framework represented by Eqs. (12) and (14), we will study $`\mathrm{\Gamma }_{\mathrm{el}}(\omega ,𝐤_N,T)`$. For $`T`$ small compared with $`T_c`$, one obtains the results shown in Fig. 3 b) and d). The initial increase in the scattering rate as the binding energy $`|\omega |`$ increases reflects the opening of the nodal phase space with increased quasiparticle energy. When $`|\omega |`$ is large compared to $`\stackrel{~}{v}_2\kappa `$, the scattering saturates at its normal state value. By expanding Eq. (12), we find that
$$\mathrm{\Gamma }_{\mathrm{el}}(\omega )\mathrm{\Gamma }_0\left(1+\frac{1}{6}\left(\frac{\stackrel{~}{v}_2\kappa }{\omega }\right)^2\right)$$
(20)
so this saturation actually occurs from values of $`\mathrm{\Gamma }_{el}`$ somewhat above $`\mathrm{\Gamma }_0`$, as seen in Fig. 3.
For forward scattering, where $`\stackrel{~}{v}_2\kappa /\mathrm{\Gamma }_0`$ is small, one can obtain a useful estimate of the quasilinear slope. In this case, $`\mathrm{\Gamma }_{\mathrm{el}}(\omega )`$ rises from its value of $`\mathrm{\Gamma }_{\mathrm{el}}(\omega =0)`$ at $`\omega =0`$ to $`\mathrm{\Gamma }_0`$ over a frequency range of order $`\stackrel{~}{v}_2\kappa `$ with a quasi-linear slope $`(\mathrm{\Gamma }_0\mathrm{\Gamma }_{\mathrm{el}}(\omega =0))/(\stackrel{~}{v}_2\kappa )`$. In the dirty limit, Eq. (17), for $`0<|\omega |<\stackrel{~}{v}_2\kappa `$ this gives
$$\mathrm{\Gamma }_{\mathrm{el}}(\omega )\mathrm{\Gamma }_{\mathrm{el}}(\omega =0)+\frac{1}{6}\left(\frac{\stackrel{~}{v}_2\kappa }{\mathrm{\Gamma }_0}\right)\omega .$$
(21)
Finally, if the weak scattering is sufficiently isotropic $`(\kappa 1)`$, we recover the standard result for pointlike weak scatterers
$$\mathrm{\Gamma }_{\mathrm{el}}(\omega )\mathrm{\Gamma }_{\mathrm{el}}(\omega =0)+\frac{\pi }{2}\frac{n_iN_0|V_0|^2\omega }{\mathrm{\Delta }_0}.$$
(22)
Which of the results, Eq. (21) or (22) is valid depends upon the character of the disorder.
To check these estimates, and evaluate the scattering rate for various scattering potential ranges and impurity concentrations, as well as to treat momenta away from the Fermi surface more accurately, we introduce a more realistic model for the elastic scattering. We consider a screened exponential fall-off in 2D such that,
$`|V(𝐤,𝐤^{})|^2`$ $`=`$ $`{\displaystyle \frac{|V_0|^2}{𝐪^2+\kappa ^2}},`$ (23)
where $`𝐪=𝐤𝐤^{}`$ is the momentum transfer. We then perform a self-consistent solution of equations (5) and (6). Results for $`\mathrm{\Gamma }_{\mathrm{el}}`$ analogous to Fig. 3, but taken “on-shell” at momentum $`𝐤`$ such that $`\epsilon _𝐤=\omega `$ are presented in Fig. 4. Results have been scaled such that the (self-consistent) normal state scattering rate is the same for all curves; as discussed above, this implies that the more forward scattering cases correspond to larger impurity concentrations. It is seen that the qualitative behavior close to the Fermi level is captured rather well by the approximate model discussed above where the scattering is restricted to a cone of width $`\kappa `$. It is also interesting to note that the quasilinear behavior at the low $`\omega `$ values persists until extremely long scattering ranges, of order 5-10 lattice spacings, for reasonable total normal state scattering rates. This suggests on the one hand that elastic scattering is a likely explanation for the quasilinear behavior seen at low energies in the MDC width; on the other hand, it also means that it will be difficult to determine the range of the scatterers precisely from such a measurement alone.
## III Inelastic Scattering
In this section we discuss the scattering rates due to inelastic electron-electron scattering by a phenomenological short range Coulomb interaction based on the two dimensional Hubbard model. For the tight-binding bandstructure considered here, these scattering processes are dominated by exchange of antiferromagnetic spin-fluctuations. We are using a conventional Berk-Schrieffer-like Berk theory as has been used before to discuss the single-particle inelastic lifetime QSB94 and has proved to give a qualitative description of low-energy microwave and thermal conductivity lifetimes DSH01 ; QHS96 as well as NMR relaxation rates BS92 in the cuprates.
In second-order perturbation theory for the two dimensional Hubbard model, the imaginary part of the nodal quasiparticle self-energy due to inelastic scattering from the onsite Coulomb interaction $`U`$ can be written as
$`\mathrm{\Gamma }_{\mathrm{inel}}(\omega ,𝐤,T)=\mathrm{\Sigma }_{inel}^{\prime \prime }(\omega ,𝐤,T)=`$ (24)
$`{\displaystyle \frac{U^2}{N}}{\displaystyle \underset{q}{}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\mathrm{\Omega }[n(\mathrm{\Omega })+f(\mathrm{\Omega }\omega )]`$
$`\chi _0^{\prime \prime }(q,\mathrm{\Omega })N(kq,\omega \mathrm{\Omega })`$
with
$`N(k,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}\mathrm{Im}\left({\displaystyle \frac{\omega +i\delta +\epsilon _k}{(\omega +i\delta )^2E_k^2}}\right)`$ (25)
and
$`\chi _0(𝐪,\mathrm{\Omega })=`$ (26)
$`{\displaystyle \frac{1}{N}}{\displaystyle \underset{k}{}}\{{\displaystyle \frac{1}{2}}[1+{\displaystyle \frac{\epsilon _{k+q}\epsilon _k+\mathrm{\Delta }_{k+q}\mathrm{\Delta }_k}{E_{k+q}E_k}}]{\displaystyle \frac{f(E_{k+q})f(E_k)}{\mathrm{\Omega }(E_{k+q}E_k)+i0^+}}`$
$`+{\displaystyle \frac{1}{4}}\left[1{\displaystyle \frac{\epsilon _{k+q}\epsilon _k+\mathrm{\Delta }_{k+q}\mathrm{\Delta }_k}{E_{k+q}E_k}}\right]{\displaystyle \frac{1f(E_{k+q})f(E_k)}{\mathrm{\Omega }+(E_{k+q}+E_k)+i0^+}}`$
$`+{\displaystyle \frac{1}{4}}[1{\displaystyle \frac{\epsilon _{k+q}\epsilon _k+\mathrm{\Delta }_{k+q}\mathrm{\Delta }_k}{E_{k+q}E_k}}]{\displaystyle \frac{f(E_{k+q})+f(E_k)1}{\mathrm{\Omega }(E_{k+q}+E_k)+i0^+}}\}.`$
Here, $`E_k=\sqrt{\epsilon _k^2+\mathrm{\Delta }_k^2}`$. As before, we are considering a nodal quasiparticle with $`k_x=k_y`$ and the magnitude of $`k`$ is set by $`\omega =\epsilon _k`$. At low energies, where $`\omega `$ and $`T`$ are small compared to the zero temperature gap amplitude, Eq. (24) determines the form of the $`\omega `$ and $`T`$ dependence of the inelastic quasiparticle scattering. Higher order processes lead to a renormalization of $`U`$, but the phase space restrictions imposed by the $`d`$-wave nature of the superconducting state determine the low $`T`$ and $`\omega `$ dependence of $`\mathrm{\Gamma }_{\mathrm{inel}}(\omega ,𝐤|_{\omega =\epsilon _𝐤},T)`$ QSB94 ; analytic ; DSH01 .
For $`𝐤=𝐤_N`$ and $`\omega =0`$, the low temperature inelastic scattering varies as
$$\mathrm{\Gamma }_{\mathrm{inel}}(0,T)\frac{T^3}{\mathrm{\Delta }_0^2}$$
(27)
with a prefactor of order one QSB94 ; analytic . For the set of parameters discussed below the prefactor becomes 2.4. It has been argued that $`\mathrm{\Gamma }_{\mathrm{inel}}`$ should vary as $`T^{5/2}`$ at low temperatures Paaske . However, for the temperatures $`T>0.05T_c`$ we have studied, the $`T^3`$ behavior provides a much better fit to our numerical results in Fig. 6(a). Similarly, at $`T=0`$ the $`\omega `$ dependence is given by
$$\mathrm{\Gamma }_{\mathrm{inel}}(\omega ,0)\frac{\omega ^3}{\mathrm{\Delta }_0^2}.$$
(28)
The extra power $`T^3`$ and $`\omega ^3`$ versus the $`T^2`$ and $`\omega ^2`$ dependence of the usual (3D) Fermi liquid inelastic quasiparticle Coulomb scattering is due to the $`\omega /\mathrm{\Delta }_0`$ dependence of the $`d`$-wave density of states. Note that the usual $`T^3`$ and $`\omega ^3`$ phonon contributions to the low energy quasiparticle scattering in a normal metal give contributions which vary as $`T^4`$ and $`\omega ^4`$ in a $`d`$-wave superconductor.
At higher energies, the problem becomes more complicated. New, collective channels may open such as the so-called $`\pi `$-resonance or possibly the $`B_{1g}`$ phonon. Here, we replace $`\chi _0(q,\omega )`$ by the RPA form
$$\chi (q,\omega )=\frac{\chi _0(q,\omega )}{1U\chi _0(q,\omega )}$$
(29)
and consider the spin-fluctuation contributions to the scattering of the nodal quasiparticle. Following the usual spin-fluctuation notation, we replace the coupling $`U^2`$ in Eq. (24) by $`\frac{3}{2}U^2`$ which in any case is simply a phenomenological coupling constant in these calculations. In principle, the $`U`$ that enters the coupling is different from the effective $`U`$ in the denominator of $`\chi `$, Eq. (29) because of vertex corrections. Here, we will for simplicity ignore this distinction so we have just the basic RPA form of the spin-fluctuation interaction. A similar approach was used to discuss the microwave and thermal conductivity lifetimes DSH01 . This approach will in principle include resonant spin excitations like those which have been proposed to explain the $`\pi `$ resonance, i.e. the peak seen by neutron scattering at around 40 meV in various cuprates, but we do not investigate these effects in detail here. Rather, we will examine what this approximation gives for the higher energy dependence of $`\mathrm{\Gamma }_{\mathrm{inel}}(\omega ,𝐤|_{\omega =\epsilon _𝐤},T)`$ numerically.
In Fig. 5, we show the numerical evaluation of Eq. (24), as well as the RPA result with $`\chi _0`$ replaced by $`\chi `$ as in Eq. (29), using the band parameters discussed above in Eqs. (7) and (8) appropriate for an optimally doped cuprate. Here, the numerical integrations in Eq. (24) and (26) are done using the technique described in Ref. DT95 . The value of $`U=2.2t`$ chosen in the RPA case comes from early fits of spin-fluctuation theory to NMR data, and more recently to microwave and thermal conductivity data on YBCO, but is expected to be reasonable for BSCCO as well because transport data suggest that the inelastic scattering rates at $`T_c`$ are very similar. The value of $`U=6.7t`$ in the case of the second order perturbation theory result was chosen so that the $`\omega ^3`$ behavior of both results matches for $`|\omega |/\mathrm{\Delta }_0<1`$ as shown in the inset. As noted, when $`|\omega |\mathrm{\Delta }_0`$, the $`\omega `$-dependence of the inelastic scattering is determined by the phase space, and one gets an $`\omega ^3`$ dependence from a constant $`U`$ interaction as well as from the spin fluctuation (RPA) interaction. However, as shown in Fig. 5, at higher energies the scattering reflects the frequency dependence of the effective interaction. In the RPA case, the classical quasilinear energy dependence expected in spin fluctuation theory when the energy exceeds the spin-fluctuation energy is clearly visible, and the crossover to the low-energy $`\omega ^3`$ form takes place around $`3\mathrm{\Delta }_0`$, as also found in Ref. QHS96 .
In Figures 6 a) and b), we study the effect of temperature on the inelastic scattering. First, in Fig. 6 a) we consider the nodal scattering rate at $`\omega =0`$, which collapses rapidly at $`T_c`$ due to the opening of the gap. In Fig. 6 b), the energy dependence of the scattering rate is plotted for various temperatures. It is interesting to note that for these band parameters, there is a considerable amount of upward curvature at low $`\omega `$ in $`\mathrm{\Gamma }_{\mathrm{el}}(\omega ,𝐤|_{\omega =\epsilon _𝐤},T)`$ even at $`T_c`$, in contrast to the simpler nearly nested bands considered in QSB94 ; QHS96 .
## IV Total scattering rate
To include both types of scattering effects, we neglect interference processes between electron-electron collisions and impurity scattering entirely, and approximate the total scattering rate at the node by
$$\mathrm{\Gamma }_{\mathrm{tot}}=\mathrm{\Gamma }_{\mathrm{el}}+\mathrm{\Gamma }_{\mathrm{inel}},$$
(30)
i.e. higher order processes like the influence of the inelastic scattering on the elastic scattering and vice versa are neglected. However, this is partially taken into account since we have chosen the parameters for the elastic and inelastic scattering on a phenomenological basis.
In Fig. 7, we plot the on-shell total scattering rate $`\mathrm{\Gamma }_{\mathrm{tot}}`$ of a nodal quasiparticle as a function of both temperature and energy. It is seen that the generic features, for reasonable assumptions about the magnitudes of the impurity scattering rates and ranges, are as follows. As a function of temperature, one expects the nodal scattering rate to collapse at $`T_c`$, and to obey a $`T^3`$ dependence at the lowest temperatures, whose coefficient is determined by both the elastic and inelastic processes. On the other hand, the low energy dependence of the scattering rate at low temperatures is dominated by the quasi-linear $`\omega `$ dependence of the elastic scattering. The generic result appears to be a quasilinear low-energy behavior for weak out of plane impurities, where the slope of this result is unrelated to the slope of the high-energy quasilinear behavior determined by spin-fluctuation scattering. A very low-energy $`(\omega 12`$meV) flattening or upturn due to unitary scatters ($`_𝐤G\omega `$) at low energies in Eq. (3), is expected to be present in principle. This contribution is likely to be negligible in BSCCO due to the small concentration of native planar defects in current samples, and the fact that the expected impurity bandwidth is comparable to the energy resolution in ARPES experiments.
## V Conclusions
Understanding the lifetime of nodal quasiparticles is a fundamental problem of cuprate physics which must be solved if we are to agree that a $`d`$-wave BCS theory describes the optimally to over-doped superconducting state, the only part of the phase diagram apparently susceptible to semiquantitative description at this time. In addition, the fact that nodal quasiparticles appear to be robust even for strongly underdoped samples may be an important clue to the physics of the pseudogap, and a theory of the nodal states within BCS which works at optimal doping may help us to understand this clue. Finally, the lifetime of the nodal quasiparticles determines thermal and microwave conductivity as well as photoemission lineshapes and one would like to have a unified description of these quantities within a single model. Here, we have calculated the on-shell $`\omega =\epsilon _k`$ nodal quasiparticle scattering rate which enters in determining the MDC linewidth. Using a simple model which parameterizes the forward elastic scattering in terms of a range $`\kappa ^1`$, a normal state scattering rate $`\mathrm{\Gamma }_0`$, and an inelastic spin-fluctuation scattering parameterized by an RPA form with an effective Coulomb coupling $`U`$, we have studied the temperature and $`\omega `$-dependence of the scattering rate.
For a quasiparticle at the nodal momentum $`k_N`$, as $`T`$ decreases below $`T_c`$, the elastic scattering rate decreases from its normal state value $`\mathrm{\Gamma }_0`$. As the gap opens, it reaches a smaller value $`\mathrm{\Gamma }_{\mathrm{el}}(T)`$ determined by the forward scattering parameter $`\kappa `$ and the renormalized gap velocity $`\stackrel{~}{v}_2`$. This temperature dependence of the elastic scattering rate arises from the temperature dependence of the gap amplitude which controls the available phase space for scattering, suppressing it as the gap opens. Because the scattering rate depends exponentially on the gap, which opens rapidly in BSCCO, the elastic scattering rate has a cusp-like onset at $`T_c`$, decreasing rapidly to its low temperature value. The inelastic contribution to scattering rate also decreases as the gap opens, with a cusp, and then decreases as $`(T/\mathrm{\Delta }_0)^3`$ at low temperatures.
At low reduced temperatures, the $`\omega `$-dependence of the elastic scattering rate can exhibit a quasi-linear behavior, varying as $`\frac{\omega }{6}\left(\frac{\stackrel{~}{v}_2\kappa }{\mathrm{\Gamma }_0}\right)`$ if the scattering is forward or as $`\mathrm{\Gamma }_0\frac{\omega }{\mathrm{\Delta }_0}`$ if the scattering is more isotropic. The inelastic scattering rate initially increases as $`\omega ^3/\mathrm{\Delta }_0^2`$ so that there is an energy beyond which the inelastic scattering becomes dominant. For energies greater than of order $`3\mathrm{\Delta }_0`$, the inelastic scattering rate crosses over to a quasi-linear $`\omega `$-dependence with a slope of order one.
Although the calculations we have presented here are straightforward, at the present time the experimental situation regarding the direct measurement of the nodal scattering rate by ARPES is somewhat uncertain. There are some claims in the literature that the rate collapses in the SC state Kam00 , as found theoretically here, and some that marginal Fermi liquid linear behavior consistent with a quantum critical point persists down to the lowest temperatures Val00 . As samples and resolution of the ARPES technique improve, we expect this discrepancy to be resolved and our prediction for the nodal quasiparticle MDC width to be testable.
###### Acknowledgements.
The authors are grateful for discussions with D. Dessau, A. Fujimori, A. Ino, P. Johnson, and Z.-X. Shen. We would like to thank A. Chubukov for his comments regarding the second order on-shell inelastic scattering rate. Work was partially supported by ONR grant N00014-04-0060 and NSF DMR02-11166. |
warning/0507/math0507548.html | ar5iv | text | # Group actions and invariants in algebras of generic matrices
## 1. Introduction
Let $`k`$ be a field of characteristic zero, $`m`$ and $`n`$ be integers $`2`$, and $`G_{m,n}=k\{X_1,\mathrm{},X_m\}`$ be the $`k`$-algebra of $`m`$ generic $`n\times n`$-matrices. That is, $`G_{m,n}`$ is the $`k`$-subalgebra of $`\mathrm{M}_n(k[x_{ij}^{(h)}])`$ generated by
$$X_1=(x_{ij}^{(1)}),\mathrm{},X_m=(x_{ij}^{(m)}),$$
where the $`x_{ij}^{(h)}`$ are $`mn^2`$ independent commuting variables. By a theorem of Amitsur, $`G_{m,n}`$ is a domain of PI-degree $`n`$. There is a natural action of the general linear group $`\mathrm{GL}_m`$ on $`G_{m,n}`$ given by
(1.1)
$$g=(g_{ij}):X_j\underset{i=1}{\overset{m}{}}g_{ij}X_i.$$
In this paper we will prove the following theorem.
###### 1.2 Theorem.
For $`2mn^22`$, the domain $`(G_{m,n})^{\mathrm{SL}_m}`$ has PI-degree $`n`$.
The trace ring $`T_{m,n}`$ of $`G_{m,n}`$ is defined as the subring of $`\mathrm{M}_n(k[x_{ij}^{(h)}])`$ generated by elements of the form $`Y`$ and $`\mathrm{tr}(Y)`$, as $`Y`$ ranges over $`G_{m,n}`$. The action (1.1) on $`G_{m,n}`$ naturally extends to $`T_{m,n}`$. Note that the algebras $`G_{m,n}`$ and $`T_{m,n}`$, and their centers $`\mathrm{Z}(G_{m,n})`$ and $`\mathrm{Z}(T_{m,n})`$ have a natural $``$-grading inherited from $`\mathrm{M}_n(k[x_{ij}^{(h)}])`$ (each variable $`x_{ij}^{(h)}`$ has degree $`1`$) and that this grading is preserved by the action (1.1). As a consequence of Theorem 1.2 we obtain the following result.
###### 1.3 Theorem.
Let $`2mn^22`$, and let $`R`$ be one of the rings $`G_{m,n}`$, $`T_{m,n}`$, $`\mathrm{Z}(G_{m,n})`$, or $`\mathrm{Z}(T_{m,n})`$. Denote the degree $`d`$ homogeneous component of $`R`$ by $`R[d]`$. Then
$$\underset{d\mathrm{}}{lim\; sup}\frac{dim_kR^{\mathrm{SL}_m}[d]}{d^{(m1)n^2m^2+1}}$$
is a finite nonzero number.
One can think of the center of $`G_{m,n}`$ as consisting of the $`m`$-variable central polynomials for $`n\times n`$-matrices (over commutative $`k`$-algebras). Theorem 1.3 thus describes, for $`R=\mathrm{Z}(G_{m,n})`$, the asymptotic behavior of the dimension of the space of $`\mathrm{SL}_m`$-invariant homogeneous central polynomials $`p(X_1,\mathrm{},X_m)`$ for $`n\times n`$-matrices.
The $`\mathrm{GL}_m`$-representations on $`G_{m,n}`$, $`\mathrm{Z}(G_{m,n})`$, $`T_{m,n}`$ and $`\mathrm{Z}(T_{m,n})`$ have been extensively studied; see, e.g., . Once again, let $`R`$ be one of these rings. Recall that the irreducible polynomial representations of $`\mathrm{GL}_m`$ are indexed by partitions $`\lambda =(\lambda _1,\mathrm{},\lambda _s)`$ with $`sm`$ parts; cf., e.g., \[9, Section 2\]. Denote the multiplicity of the irreducible $`\mathrm{GL}_m`$-representation corresponding to $`\lambda `$ in $`R`$ by $`\mathrm{mult}_\lambda (R)`$. If $`(r^m)`$ is the partition $`(r,\mathrm{},r)`$ ($`m`$ times) then it is easy to show that
$$dimR^{\mathrm{SL}_m}[d]=\{\begin{array}{cc}\mathrm{mult}_{(r^m)}(R)\hfill & \text{if }d=rm\text{,}\hfill \\ \mathrm{\hspace{0.17em}0}\hfill & \text{if }d\text{ is not a multiple of }m\text{;}\hfill \end{array}$$
see Remark 9.2. The conclusion of Theorem 1.3 can thus be rephrased by saying that
$$\underset{r\mathrm{}}{lim\; sup}\frac{\mathrm{mult}_{(r^m)}(R)}{r^{(m1)n^2m^2+1}}$$
is a finite nonzero number. We also note that by the Berele-Drensky-Formanek correspondence, $`\mathrm{mult}_{(r^m)}(R)`$ equals the multiplicity of the $`S_m`$-character $`\chi ^{(d^m)}`$ in the cocharacter sequence of $`R`$; see \[9, Section 4\].
The division algebra of quotients of $`G_{m,n}`$ (or equivalently, of $`T_{m,n}`$) is called the universal division algebra of $`m`$ generic $`n\times n`$-matrices; we shall denote it by $`\mathrm{𝑈𝐷}(m,n)`$. Note that the $`\mathrm{GL}_m`$-action (1.1) on $`G_{m,n}`$ naturally extends to $`\mathrm{𝑈𝐷}(m,n)`$. We shall deduce Theorem 1.2 from the following related result.
###### 1.4 Theorem.
If $`2mn^22`$ and $`n3`$, then $`\mathrm{𝑈𝐷}(m,n)^{\mathrm{GL}_m}`$ is a division algebra of degree n.
For all other values of $`m,n2`$, $`\mathrm{𝑈𝐷}(m,n)^{\mathrm{GL}_m}`$ is a field; see Propositions 8.1 and 8.3. A brief summary of the properties of $`\mathrm{𝑈𝐷}(m,n)^{\mathrm{GL}_m}`$ and $`\mathrm{𝑈𝐷}(m,n)^{\mathrm{SL}_m}`$ is given in the two tables below.
The assertions of the tables in the cases where $`mn^22`$ and $`n3`$ are based on Theorems 1.4 and 5.1, the case where $`m=n=2`$ is considered in \[20, Section 14\], and the cases where $`mn^21`$ are treated in detail in Section 8.
It appears likely that Theorems 1.2 – 1.4 remain valid in prime characteristic (perhaps, not dividing $`n`$); we have not attempted to extend them in this direction. Our arguments rely on the work of Richardson and on our own prior papers , all of which make the characteristic zero assumption.<sup>1</sup><sup>1</sup>1We remark that Richardson worked over $`k=`$, and his proofs are based on analytic techniques. The results we need (in particular, \[21, Theorem 9.3.1\]), remain valid over any algebraically closed field of characteristic zero by the Lefschetz principle. Extending \[21, Theorem 9.3.1\] to prime characteristic is an open problem of independent interest.
### Conventions and Terminology.
All central simple algebras in this paper are assumed to be finite-dimensional over their centers. All algebraic varieties, algebraic groups, group actions, morphisms, rational maps, etc., are assumed to be defined over the base field $`k`$ (which we always assume to be of characteristic zero). By a point of an algebraic variety $`X`$ we shall always mean a $`k`$-point. Throughout, $`G`$ will denote a linear algebraic group. We shall refer to an algebraic variety $`X`$ endowed with a regular $`G`$-action as a $`G`$-variety. We will say that a $`G`$-variety $`X`$ (or the $`G`$-action on $`X`$) is generically free if $`\mathrm{Stab}_G(x)=\{1\}`$ for $`xX`$ in general position. Finally, unless otherwise specified, $`m`$ and $`n`$ are integers $`2`$.
## 2. Preliminaries
### Concomitants
Let $`\mathrm{\Gamma }`$ be an algebraic group and $`V`$ and $`W`$ be $`\mathrm{\Gamma }`$-varieties. Then we shall denote the set of $`\mathrm{\Gamma }`$-equivariant morphisms $`VW`$ (also known as concomitants) by $`\mathrm{Morph}_\mathrm{\Gamma }(V,W)`$ and the set of $`\mathrm{\Gamma }`$-equivariant rational maps $`VW`$ (also known as rational concomitants) by $`\mathrm{RMaps}_\mathrm{\Gamma }(V,W)`$.
In the case where $`W`$ is a finite-dimensional linear representation of $`\mathrm{\Gamma }`$, we also define a relative concomitant as a morphism $`f:VW`$ satisfying the following condition (which is slightly weaker than $`\mathrm{\Gamma }`$-equivariance): there is a character $`\chi :\mathrm{\Gamma }k^{}`$ such that
$$f(gv)=\chi (g)\left(gf(v)\right)$$
for all $`vV`$ and $`g\mathrm{\Gamma }`$. For a rational map $`f:VW`$ the notion of a relative rational concomitant is defined in a similar manner. If $`W=k`$, with trivial $`\mathrm{\Gamma }`$-action, then the term “invariant” is used in place of “concomitant”. For future reference we record the following:
###### 2.1 Lemma.
Suppose $`V`$ and $`W`$ are finite-dimensional linear representations of $`\mathrm{\Gamma }`$. Every rational concomitant $`f:VW`$ can be written as $`\frac{a}{b}`$, where $`a`$ is a relative concomitant and $`b`$ is a relative invariant.
###### Proof.
See the proof of \[5, Chapter 1, Proposition 1\]. Note that the characters associated to $`a`$ and $`b`$ are necessarily equal. ∎
If $`W`$ is a $`k`$-algebra and $`\mathrm{\Gamma }`$ acts on $`W`$ by $`k`$-algebra automorphisms, then the algebra structure of $`W`$ induces algebra structures on $`\mathrm{Morph}_\mathrm{\Gamma }(V,W)`$ and $`\mathrm{RMaps}_\mathrm{\Gamma }(V,W)`$ in a natural way. Namely, given $`a,b:VW`$ (or $`a,b:VW`$), one defines $`a+b`$ and $`ab`$ by $`(a+b)(v)=a(v)+b(v)`$ and $`ab(v)=a(v)b(v)`$ for $`vV`$.
###### 2.2 Theorem.
(Procesi \[16, Theorem 2.1\]; cf. also \[9, Theorem 10\], or \[24, Theorem 14.16\]) Let $`(\mathrm{M}_n)^m`$ be the space of $`m`$-tuples of $`n\times n`$-matrices; the group $`\mathrm{PGL}_n`$ acts on it by simultaneous conjugation. Then
* $`\mathrm{Morph}_{\mathrm{PGL}_n}((\mathrm{M}_n)^m,\mathrm{M}_n)T_{m,n}`$
* $`\mathrm{RMaps}_{\mathrm{PGL}_n}((\mathrm{M}_n)^m,\mathrm{M}_n)\mathrm{𝑈𝐷}(m,n)`$
Moreover, the two isomorphisms identify the $`i`$-th projection $`(\mathrm{M}_n)^m\mathrm{M}_n`$ with the $`i`$-th generic matrix $`X_i`$. ∎
Here $`T_{m,n}`$ and $`\mathrm{𝑈𝐷}(m,n)`$ are, respectively, the trace ring and the universal division algebra of $`m`$ generic $`n\times n`$-matrices, defined in the introduction. Note that part (b) of Theorem 2.2 follows from part (a) by Lemma 2.1, since the simple group $`\mathrm{PGL}_n`$ does not have nontrivial characters (so that relative concomitants and invariants are actually concomitants and invariants, respectively).
We also remark that the construction of $`T_{m,n}`$ remains well-defined if $`m=1`$. Theorem 2.2 also holds in this case, provided that one defines $`\mathrm{𝑈𝐷}(1,n)`$ to be the field of quotients of $`T_{1,n}`$, rather than $`G_{1,n}`$. (For $`m2`$, $`T_{m,n}`$ and $`G_{m,n}`$ have the same division algebra of quotients, but this is not the case for $`m=1`$.).
### Geometric actions
For the rest of this section we will assume that $`k`$ is algebraically closed. If $`X`$ is a $`\mathrm{PGL}_n`$-variety, then, as we mentioned above, $`\mathrm{RMaps}_{\mathrm{PGL}_n}(X,\mathrm{M}_n)`$ has an algebra structure naturally induced from $`\mathrm{M}_n`$. If the $`\mathrm{PGL}_n`$-action on $`X`$ is generically free then $`\mathrm{RMaps}_{\mathrm{PGL}_n}(X,\mathrm{M}_n)`$ is a central simple algebra of degree $`n`$, with center $`k(X)^{\mathrm{PGL}_n}`$; cf. \[18, Lemmas 8.5 and 9.1\].
Suppose that $`X`$ is a $`G\times \mathrm{PGL}_n`$-variety, and that the $`\mathrm{PGL}_n`$-action on $`X`$ is generically free. Then the $`G`$-action on $`X`$ naturally induces a $`G`$-action on $`\mathrm{RMaps}_{\mathrm{PGL}_n}(X,\mathrm{M}_n)`$. Following we define the action of an algebraic group $`G`$ on a central simple algebra $`A`$ to be geometric if $`A`$ is $`G`$-equivariantly isomorphic to $`\mathrm{RMaps}_{\mathrm{PGL}_n}(X,\mathrm{M}_n)`$ for some $`G\times \mathrm{PGL}_n`$-variety $`X`$ as above. The $`G\times \mathrm{PGL}_n`$-variety $`X`$ is then called the associated variety for the $`G`$-action on $`A`$; this associated variety is unique (as a $`G\times \mathrm{PGL}_n`$-variety), up to birational isomorphism; cf. \[20, Corollary 3.2\].
Note that we defined geometric actions only if $`k`$ is algebraically closed. Also note that if an algebraic group acts geometrically on a central simple algebra $`A`$, then the center of $`A`$ is necessarily a finitely generated field extension of $`k`$.
Of particular interest to us will be the case where $`X=(\mathrm{M}_n)^m`$ is the space of $`m`$-tuples of $`n\times n`$-matrices. Here $`\mathrm{PGL}_n`$ acts on $`(\mathrm{M}_n)^m`$ by simultaneous conjugation (since $`m2`$, this action is generically free) and $`G=\mathrm{GL}_m`$ acts on $`(A_1,\mathrm{},A_m)(\mathrm{M}_n)^m`$ by sending $`(A_1,\mathrm{},A_m)`$ to $`(B_1,\mathrm{},B_m)`$ where $`B_j=_{i=1}^mc_{ij}A_i`$ and $`g^1=(c_{ij})`$. The actions of $`\mathrm{GL}_m`$ and $`\mathrm{PGL}_n`$ commute, and the $`\mathrm{GL}_m`$-action on $`(\mathrm{M}_n)^m`$ induces the $`\mathrm{GL}_m`$-action (1.1) on $`\mathrm{𝑈𝐷}(m,n)`$. So $`(\mathrm{M}_n)^m`$ is the associated variety for the $`\mathrm{GL}_m`$-action on $`\mathrm{𝑈𝐷}(m,n)`$; see Theorem 2.2 (cf. also \[20, Example 3.4\]).
We conclude this section with a simple result which we will use repeatedly.
###### 2.3 Lemma.
Assume $`k`$ is algebraically closed. Let $`X`$ be a $`G\times \mathrm{PGL}_n`$-variety which is $`\mathrm{PGL}_n`$-generically free. Denote by $`\pi :XX/G`$ the rational quotient map for the $`G`$-action. Then for $`xX`$ in general position, the projection $`\mathrm{pr}_2:G\times \mathrm{PGL}_n\mathrm{PGL}_n`$ onto the second factor induces an isomorphism from $`\mathrm{Stab}_{G\times \mathrm{PGL}_n}(x)/\mathrm{Stab}_G(x)`$ onto $`\mathrm{Stab}_{\mathrm{PGL}_n}(\pi (x))`$. ∎
###### Proof.
Recall that by a theorem of Rosenlicht, $`\pi ^1(\overline{x})`$ is a single $`G`$-orbit for $`\overline{x}X/G`$ in general position; see \[22, Theorem 2\] or \[14, Section 2.3\]. Consequently, for $`xX`$ in general position the projection $`\mathrm{pr}_2`$ restricts to a surjective morphism $`\mathrm{Stab}_{G\times \mathrm{PGL}_n}(x)\mathrm{Stab}_{\mathrm{PGL}_n}(\pi (x))`$ of algebraic groups. The kernel of this morphism is clearly $`\mathrm{Stab}_G(x)`$, and the lemma follows. ∎
## 3. Geometric actions on division algebras
Throughout this section we will assume that $`k`$ is algebraically closed. The main result of this section is the following theorem; after its proof, we will deduce several corollaries.
###### 3.1 Theorem.
Assume $`k`$ is algebraically closed. Let $`G`$ be a linear algebraic group acting geometrically on a division algebra $`D`$ of degree $`n`$. Let $`X`$ be the associated $`G\times \mathrm{PGL}_n`$-variety. Then for $`xX`$ in general position,
$$S_x:=\mathrm{Stab}_{G\times \mathrm{PGL}_n}(x)/\mathrm{Stab}_G(x)$$
is reductive.
###### Proof.
Let $`X`$ be the associated $`G\times \mathrm{PGL}_n`$-variety for the $`G`$-action on $`D`$. Recall that the $`\mathrm{PGL}_n`$-action on $`X`$ is generically free. We want to show that the group $`S_x=\mathrm{Stab}_{G\times \mathrm{PGL}_n}(x)/\mathrm{Stab}_G(x)`$ is reductive for $`xX`$ in general position. Assume the contrary. Denoting the unipotent radical by $`R_u`$, this means that $`R_u(\mathrm{Stab}_{G\times \mathrm{PGL}_n}(x))`$ is not contained in $`G`$. Since unipotent groups are connected, this is equivalent to
(3.2)
$$\mathrm{𝐿𝑖𝑒}\left(R_u(\mathrm{Stab}_{G\times \mathrm{PGL}_n}(x))\right)\mathrm{𝐿𝑖𝑒}(G)$$
for $`xX`$ in general position. Here and in the sequel $`\mathrm{𝐿𝑖𝑒}`$ stands for the Lie algebra. To simplify notation, set $`H=G\times \mathrm{PGL}_n`$, and for $`xX`$, set $`H_x=R_u(\mathrm{Stab}_H(x))`$. Now define $`U_XX\times \mathrm{𝐿𝑖𝑒}(H)`$ by
$$U_X=\{(x,a)|xX\text{ and }a\mathrm{𝐿𝑖𝑒}(H_x)\}.$$
We first show that $`U_X`$ is a vector bundle over a dense open subset $`X_0X`$. By \[21, 6.2.1, 9.2.1, and 6.5.3\], there is an $`H`$-stable dense open subset $`X_0`$ of $`X`$ such that $`\{H_xxX_0\}`$ is an algebraic family of algebraic subgroups of $`H`$. Moreover, $`dim(H_x)`$ is constant for $`xX_0`$, say equal to $`d`$. Replacing $`X`$ by $`X_0`$, we may assume that $`\{H_xxX\}`$ is an algebraic family of algebraic subgroups of $`H`$. By \[21, 6.2.2\], $`x\mathrm{𝐿𝑖𝑒}(H_x)`$ defines a morphism of algebraic varieties from $`X`$ to the Grassmannian of $`d`$-dimensional subspaces of $`\mathrm{𝐿𝑖𝑒}(H)`$. Since the universal bundle over this Grassmannian is a vector bundle (see, e.g., \[28, 3.3.1\]), its pull-back $`U_X`$ is a vector bundle over $`X`$.
Note also that $`U_X`$ is, by definition, an $`H`$-invariant subbundle of the trivial bundle $`X\times \mathrm{𝐿𝑖𝑒}(H)X`$; here $`H`$ acts on its Lie algebra by the adjoint action. Since the $`\mathrm{PGL}_n`$-action on $`X`$ is generically free, the no-name lemma tells us that there is a $`\mathrm{PGL}_n`$-equivariant birational isomorphism $`U_XX\times k^d`$ such that the following diagram commutes
(For a proof and a brief discussion of the history of the no-name lemma, see \[4, Section 4.3\].) In other words, the vector bundle $`U_XX`$ has $`d`$ $`\mathrm{PGL}_n`$-equivariant rational sections $`\beta _1,\mathrm{},\beta _d:XU_X`$ such that $`\beta _1(x)`$, …, $`\beta _d(x)`$ are linearly independent for $`xX`$ in general position. We identify here $`\beta _i(x)`$ with $`a`$ if $`\beta _i(x)=(x,a)\{x\}\times \mathrm{𝐿𝑖𝑒}(H_x)`$. In view of (3.2), some $`k`$-linear combination $`\beta =c_1\beta _1+\mathrm{}+c_d\beta _d`$ has the property that $`\beta (x)\mathrm{𝐿𝑖𝑒}(G)`$ for $`xX`$ in general position.
Now recall that the natural projection $`\mathrm{SL}_n\mathrm{PGL}_n`$ induces a Lie algebra isomorphism $`\text{sl}_n\mathrm{𝐿𝑖𝑒}(\mathrm{PGL}_n)`$, allowing us to identify the two Lie algebras. Hence
$$U_XX\times \mathrm{𝐿𝑖𝑒}(G)\times \text{sl}_n.$$
Let $`f=pr\beta :X\text{sl}_n`$, where $`pr:U_X\text{sl}_n`$ denotes the natural projection. Note that $`\text{sl}_n\text{gl}_n=\mathrm{M}_n`$, so that $`f`$ may be viewed as a $`\mathrm{PGL}_n`$-equivariant rational map $`X\mathrm{M}_n`$, i.e., as an element of $`D`$. The condition that $`\beta (x)\mathrm{𝐿𝑖𝑒}(G)`$ ensures that $`f0`$. On the other hand, we will show below that for $`xX`$ in general position, $`f(x)`$ is a nilpotent $`n\times n`$-matrix, so that $`f^n=0`$. This means that $`D`$ contains a non-zero nilpotent element $`f`$, contradicting our assumption that $`D`$ is a division algebra.
It remains to be shown that for $`xX`$ in general position, $`f(x)`$ is a nilpotent matrix. The natural projection $`G\times \mathrm{PGL}_n\mathrm{PGL}_n`$ maps the unipotent group $`H_x`$ to a unipotent subgroup $`U`$ of $`\mathrm{PGL}_n`$. Denote by $`K`$ the preimage of $`U`$ in $`\mathrm{SL}_n`$. It is a solvable group, so its subset $`K_u`$ of unipotent elements is a closed subgroup. The surjection $`K_uU`$ is finite-to-one, so their Lie algebras are isomorphic. In particular, $`f(x)`$ belongs to $`\mathrm{𝐿𝑖𝑒}(K_u)\text{sl}_n\text{gl}_n=\mathrm{M}_n`$. Finally, since $`K_u`$ is a unipotent subgroup of $`\mathrm{GL}_n`$, its Lie algebra in $`\mathrm{M}_n`$ consists of nilpotent matrices, see, e.g., \[3, I.4.8\]. This completes the proof of Theorem 3.1. ∎
We now proceed with the corollaries. Recall that a subgroup $`S\mathrm{\Gamma }`$ is said to be a stabilizer in general position for a $`\mathrm{\Gamma }`$-variety $`X`$ if there exists a dense $`\mathrm{\Gamma }`$-invariant subset $`UX`$ such that $`\mathrm{Stab}(x)`$ is conjugate to $`S`$ for any $`xU`$. For a detailed discussion of this notion, see \[14, Section 7\].
###### 3.3 Corollary.
Assume $`k`$ is algebraically closed. Let $`G`$ be a linear algebraic group acting geometrically on a division algebra $`D`$ of degree $`n`$. Let $`X`$ be the associated $`G\times \mathrm{PGL}_n`$-variety.
* The induced $`\mathrm{PGL}_n`$-action on the rational quotient $`X/G`$ has a stabilizer $`S`$ in general position. Moreover, $`S`$ is reductive, and $`SS_x=\mathrm{Stab}_{G\times \mathrm{PGL}_n}(x)/\mathrm{Stab}_G(x)`$ for $`xX`$ in general position.
* If the $`G`$-action on $`X`$ is generically free, then
$`\mathrm{trdeg}_k(\mathrm{Z}(D^G))`$ $`=\mathrm{trdeg}_k(\mathrm{Z}(D)^G)`$
$`=dim(X)dim(G)+dim(S)n^2+1.`$
###### Proof.
(a) It follows from Theorem 3.1 and Lemma 2.3 that points in $`X/G`$ in general position have a reductive stabilizer. A theorem of Richardson (see \[21, Theorem 9.3.1\] or \[14, Theorem 7.1\]) now implies that the $`\mathrm{PGL}_n`$-action on $`X/G`$ has a stabilizer $`S`$ in general position. By Lemma 2.3, $`SS_x=\mathrm{Stab}_{G\times \mathrm{PGL}_n}(x)/\mathrm{Stab}_G(x)`$ for $`xX`$ in general position.
(b) The first equality follows from the fact that $`\mathrm{Z}(D^G)`$ is an algebraic extension of $`\mathrm{Z}(D)^G`$. Indeed, the minimal polynomial of any element of $`D^G`$ over $`\mathrm{Z}(D)`$ is unique and must therefore have coefficients in $`\mathrm{Z}(D)^G`$.
To prove the second equality, note that, $`\mathrm{Z}(D)=k(X/\mathrm{PGL}_n)=k(X)^{\mathrm{PGL}_n}`$ and thus
$$\mathrm{Z}(D)^G(k(X)^{\mathrm{PGL}_n})^G=k(X)^{G\times \mathrm{PGL}_n}=k(X/(G\times \mathrm{PGL}_n)).$$
Since we are assuming that $`G`$ acts generically freely on $`X`$, part (a) implies that $`S\mathrm{Stab}_{G\times \mathrm{PGL}_n}(x)`$ for $`xX`$ in general position. Hence the dimension of the general fiber of the rational quotient map $`XX/(G\times \mathrm{PGL}_n)`$ is equal to the dimension of $`(G\times \mathrm{PGL}_n)/S`$. The fiber dimension theorem now tells us that the transcendence degree of $`\mathrm{Z}(D)^G`$ is
$$dimX/(G\times \mathrm{PGL}_n)=dim(X)dim(G)dim(\mathrm{PGL}_n)+dim(S).\mathit{}$$
###### 3.4 Corollary.
Assume $`k`$ is algebraically closed. Let $`G`$ be a unipotent linear algebraic group acting geometrically on a division algebra $`D`$ of degree $`n`$. Then $`D^G`$ is a division algebra of degree $`n`$.
This was proved for algebraic actions in \[20, Proposition 12.1\].
###### Proof.
By \[20, Lemma 7.1\], for $`xX`$ in general position, $`\mathrm{Stab}_{G\times \mathrm{PGL}_n}(x)`$ is isomorphic to a subgroup of $`G`$, so is unipotent. On the other hand, by Theorem 3.1, the projection $`S_x`$ of this group to $`\mathrm{PGL}_n`$ is reductive. Thus $`S_x`$ is both unipotent and reductive, which is only possible if $`S_x=\{1\}`$. In other words,
$$\mathrm{Stab}_{G\times \mathrm{PGL}_n}(x)G\times \{1\}.$$
The desired conclusion now follows from \[20, Theorem 1.4\]. ∎
## 4. Dimension counting in the Grassmannian
In preparation for the proof of Theorem 1.4 in the next section, we will now establish the following:
###### 4.1 Proposition.
Assume $`k`$ is algebraically closed. Let $`V`$ be an $`N`$-dimensional $`k`$-vector space and let $`V=V_1\mathrm{}V_r`$, where $`dim(V_i)=N_i1`$. Let $`Z`$ be the subset of the Grassmannian $`\mathrm{Gr}(m,N)`$ consisting of $`m`$-dimensional subspaces $`W`$ of $`V`$ of the form
$$W=W_1\mathrm{}W_r,$$
where $`W_iV_i`$. (Here we allow $`W_i=(0)`$ for some $`i`$.) Then $`Z`$ is a closed subvariety of $`\mathrm{Gr}(m,N)`$. If $`2mN2`$, then each irreducible component of $`Z`$ has codimension $`N\mathrm{max}_{i=1,\mathrm{},r}(N_i)`$ in $`\mathrm{Gr}(m,N)`$. Moreover, equality holds (for some irreducible component of $`Z`$) only if (i) $`r=1`$ or (ii) $`r=2`$, $`m=2`$ and $`N=4`$.
###### Proof.
Let $`m_1,\mathrm{}m_r`$ be non-negative integers such that $`m_1+\mathrm{}+m_r=m`$ and such that $`m_iN_i`$ for all $`i`$. Let $`Z_{m_1,\mathrm{},m_r}`$ be the image of the map
$$\varphi _{m_1,\mathrm{},m_r}:\mathrm{Gr}(m_1,N_1)\times \mathrm{}\times \mathrm{Gr}(m_r,N_r)\mathrm{Gr}(m,N)$$
given by $`(W_1,\mathrm{},W_r)W_1\mathrm{}W_r`$. (Here $`\mathrm{Gr}(m_i,N_i)`$ is the Grassmannian of $`m_i`$-dimensional vector subspaces of $`V_i`$.) Since the domain of the map $`\varphi _{m_1,\mathrm{},m_r}`$ is projective, its image is closed in $`\mathrm{Gr}(m,N)`$. Thus each $`Z_{m_1,\mathrm{},m_r}`$ is a closed irreducible subvariety of $`\mathrm{Gr}(m,N)`$ birationally isomorphic to
$$\mathrm{Gr}(m_1,N_1)\times \mathrm{}\times \mathrm{Gr}(m_r,N_r)$$
and $`Z`$ is the union of the $`Z_{m_1,\mathrm{},m_r}`$. It remains to show that
(4.2)
$$dim\mathrm{Gr}(m,N)\underset{i=1}{\overset{r}{}}dim\mathrm{Gr}(m_i,N_i)N\underset{i=1,\mathrm{},r}{\mathrm{max}}(N_i),$$
and that equality is only possible if $`r=1`$ or $`r=2`$, $`N_1=N_2=2`$ and $`m_1=m_2=1`$ (and thus $`N=N_1+N_2=4`$ and $`m=m_1+m_2=2`$). Recall that $`dim\mathrm{Gr}(m,N)=(Nm)m`$. Letting $`l_i=N_im_i`$ and $`l=Nm=l_1+\mathrm{}+l_r`$, we can rewrite (4.2) as
$$lm\underset{i=1}{\overset{r}{}}l_im_il+m\underset{i=1,\mathrm{},r}{\mathrm{max}}(l_i+m_i)$$
or, equivalently,
$$(l1)(m1)1\underset{i=1}{\overset{r}{}}l_im_i\underset{i=1,\mathrm{},r}{\mathrm{max}}(l_i+m_i).$$
Proposition 4.1 is thus a consequence of the following elementary lemma. ∎
###### 4.3 Lemma.
Let $`(l_1,m_1),\mathrm{},(l_r,m_r)`$ be $`r`$ pairs of non-negative integers and let $`l=_{i=1}^rl_i`$ and $`m=_{i=1}^rm_i`$. Assume that $`l_i+m_i1`$ for every $`i=1,\mathrm{},r`$ and $`l,m2`$. Then
(4.4)
$$(l1)(m1)1\underset{i=1}{\overset{r}{}}l_im_i\underset{i=1,\mathrm{},r}{\mathrm{max}}(l_i+m_i).$$
Moreover, equality holds if and only if either (i) $`r=1`$ or (ii) $`r=2`$ and $`(l_1,m_1)=(l_2,m_2)=(1,1)`$.
###### Proof.
We consider two cases.
Case 1: Suppose that for every $`i=1,\mathrm{},r`$, either $`l_i=0`$ or $`m_i=0`$. Since $`l,m2`$, we have $`(l1)(m1)10`$. On the other hand, $`_{i=1}^rl_im_i\mathrm{max}_{i=1,\mathrm{},r}(l_i+m_i)=\mathrm{max}_{i=1,\mathrm{},r}(l_i+m_i)<0`$. Hence, in this case (4.4) holds and is a strict inequality.
Case 2: Now suppose that $`l_i,m_i1`$ for some $`i1,\mathrm{},r`$. After renumbering the pairs $`(l_1,m_1),\mathrm{},(l_r,m_r)`$, we may assume $`i=1`$. Now set
$$l_j^{}=\{\begin{array}{cc}l_11\text{, if }j=1\text{,}\hfill & \\ l_j\text{, otherwise;}\hfill & \end{array}\text{and}m_j^{}=\{\begin{array}{cc}m_11\text{, if }j=1\hfill & \\ m_j\text{, otherwise.}\hfill & \end{array}$$
Note that $`l_j^{},m_j^{}0`$ for every $`j=1,\mathrm{},r`$. Thus
$`(l1)(m1)1`$ $`=({\displaystyle \underset{i=1}{\overset{r}{}}}l_i^{})({\displaystyle \underset{j=1}{\overset{r}{}}}m_j^{})1={\displaystyle \underset{i=1}{\overset{r}{}}}l_i^{}m_i^{}+{\displaystyle \underset{ij}{}}l_i^{}m_j^{}1`$
$`{\displaystyle \underset{i=1}{\overset{r}{}}}l_i^{}m_i^{}1={\displaystyle \underset{i=1}{\overset{r}{}}}l_im_i(l_1+m_1)`$
$`{\displaystyle \underset{i=1}{\overset{r}{}}}l_im_i\underset{i=1,\mathrm{},r}{\mathrm{max}}(l_i+m_i).`$
This completes the proof of the inequality (4.4).
It is easy to see that equality holds in cases (i) and (ii). It remains to show that the inequality (4.4) is strict for all other choices of $`(l_1,m_1),\mathrm{},(l_r,m_r)`$. Indeed, a closer look at the above argument shows that equality in (4.4) can hold if and only if we are in Case 2 and
(a) $`l_i^{}m_j^{}=0`$ whenever $`ij`$ and
(b) $`l_1+m_1=\mathrm{max}_{i=1,\mathrm{},r}(l_i+m_i)`$.
Assume that conditions (a) and (b) are satisfied. Since $`_{i=1}^rl_i^{}=l11`$, we cannot have $`l_i^{}=0`$ for all $`i=1,\mathrm{},r`$. In other words, $`l_{i_0}^{}1`$ for some $`i_0\{1,\mathrm{},r\}`$. Then condition (a) says that $`m_j^{}=0`$ for every $`ji_0`$. On the other hand, $`m_{i_0}^{}=_{j=1}^rm_j^{}=m11`$, and applying condition (a) once again, we conclude that that $`l_i^{}=0`$ for every $`ii_0`$. To sum up, there exists an $`i_0\{1,\mathrm{},r\}`$ such that $`l_{i_0}^{}=l1`$, $`m_{i_0}=m1`$ and $`l_i^{}=m_i^{}=0`$ for every $`ii_0`$.
In particular, for every $`i1,i_0`$, we have $`l_i=l_i^{}=0`$ and $`m_i=m_i^{}=0`$, contradicting our assumption that $`l_i+m_i1`$. Thus $`i\{1,i_0\}`$ for every $`i=1,\mathrm{},r`$. In other words, either $`i_0=1`$ and $`r=1`$ (in which case (i) holds, and we are done) or $`i_0=2`$ and $`r=2`$. In the latter case $`l_1^{}=m_1^{}=0`$, $`l_2^{}=l1`$ and $`m_2^{}=m1`$, i.e., $`(l_1,m_1)=(1,1)`$ and $`(l_2,m_2)=(l1,m1)`$. Condition (b) now tells us that $`l=m=2`$, so that (ii) holds.
This completes the proof of Lemma 4.3 and thus of Proposition 4.1. ∎
## 5. Proof of Theorem 1.4 over an algebraically closed field
Recall from Section 2 that $`X=(\mathrm{M}_n)^m`$ is the associated $`\mathrm{GL}_m\times \mathrm{PGL}_n`$-variety for the $`\mathrm{GL}_m`$-action on $`\mathrm{𝑈𝐷}(m,n)`$. Here $`\mathrm{PGL}_n`$ acts on $`(\mathrm{M}_n)^m`$ by simultaneous conjugation (since $`m2`$, this action is generically free) and $`\mathrm{GL}_m`$ acts on $`(\mathrm{M}_n)^m`$ by sending $`(A_1,\mathrm{},A_m)`$ to $`(B_1,\mathrm{},B_m)`$, where $`B_j=_{i=1}^mc_{ij}A_i`$ and $`g^1=(c_{ij})`$.
We shall assume throughout this section that $`k`$ is an algebraically closed field of characteristic zero. Our goal is to prove Theorem 1.4 over such $`k`$. In view of \[20, Theorem 1.4(a)\], we only need to establish the following.
###### 5.1 Theorem.
Assume $`k`$ is algebraically closed. If $`n3`$ and $`2mn^22`$, then the $`\mathrm{GL}_m\times \mathrm{PGL}_n`$-action on $`(\mathrm{M}_n)^m`$ is generically free.
###### Proof.
The linear action of $`\mathrm{GL}_m`$ on $`(\mathrm{M}_n)^m`$ is easily seen to be the direct sum of $`n^2`$ copies of the natural $`m`$-dimensional representation of $`\mathrm{GL}_m`$, i.e., to be isomorphic to the $`\mathrm{GL}_m`$-action on $`n^2`$-tuples of vectors in $`k^m`$. Since $`n^2>m`$, this action is generically free. Corollary 3.3(a) with $`G=\mathrm{GL}_m`$ and $`X=(\mathrm{M}_n)^m`$ tells us that the $`\mathrm{PGL}_n`$-action on $`(\mathrm{M}_n)^m/\mathrm{GL}_m`$ has a reductive stabilizer $`S`$ in general position, and that $`S\mathrm{Stab}_{\mathrm{GL}_m\times \mathrm{PGL}_n}(x)/\mathrm{Stab}_{\mathrm{GL}_m}(x)=\mathrm{Stab}_{\mathrm{GL}_m\times \mathrm{PGL}_n}(x)`$ for $`xX`$ in general position.
Recall that $`(\mathrm{M}_n)^m/\mathrm{GL}_m`$ is $`\mathrm{PGL}_n`$-equivariantly birationally isomorphic to the Grassmannian $`\mathrm{Gr}(m,n^2)`$ of $`m`$-dimensional subspaces of $`\mathrm{M}_n`$. Thus the $`\mathrm{PGL}_n`$-action on $`\mathrm{Gr}(m,n^2)`$ has a stabilizer $`S`$ in general position, where $`S`$ is a reductive subgroup of $`\mathrm{PGL}_n`$. (Recall that $`S`$ is only well-defined up to conjugacy in $`\mathrm{PGL}_n`$). To prove Theorem 5.1, it suffices to show that $`S`$ is trivial.
Assume the contrary. Since $`S`$ is reductive, it contains a non-trivial element $`g`$ of finite order. Then every $`L\mathrm{Gr}(m,n^2)`$ in general position is fixed by some conjugate of $`g`$. In other words, the map
(5.2)
$$\begin{array}{ccc}\mathrm{PGL}_n\times \mathrm{Gr}(m,n^2)^g& & \mathrm{Gr}(m,n^2)\\ (h,L)& & h(L)\end{array}$$
is dominant; here $`\mathrm{Gr}(m,n^2)^g`$ denotes the fixed points of $`g`$ in $`\mathrm{Gr}(m,n^2)`$. Denote by $`C(g)`$ the centralizer of $`g`$ in $`\mathrm{PGL}_n`$. Note that $`\mathrm{Gr}(m,n^2)^g`$ is $`C(g)`$-stable. Hence the fiber of (5.2) over $`h(L)`$ contains $`(hc,c^1(L))`$ for every $`cC(g)`$. So by the fiber dimension theorem,
(5.3)
$$dim\mathrm{Gr}(m,n^2)+dimC(g)dim\mathrm{PGL}_n+dim\mathrm{Gr}(m,n^2)^g.$$
Since $`g`$ has finite order, it is diagonalizable. So we may assume that
$$g=\mathrm{diag}(\underset{l_1\text{ times}}{\underset{}{\lambda _1,\mathrm{},\lambda _1}},\mathrm{},\underset{l_s\text{ times}}{\underset{}{\lambda _s,\mathrm{},\lambda _s}})=\mathrm{diag}(\alpha _1,\mathrm{},\alpha _n),$$
where $`\lambda _1,\mathrm{},\lambda _s`$ are the (distinct) eigenvalues of $`g`$. Note that $`s2`$, because $`g1`$ in $`\mathrm{PGL}_n`$ and that $`g`$ acts on the matrix units $`E_{ij}`$ by $`gE_{ij}=\alpha _i\alpha _j^1E_{ij}`$. The matrix algebra $`\mathrm{M}_n`$ naturally decomposes as a direct sum of character spaces
$$V_\mu =\mathrm{Span}(E_{ij}|\alpha _i\alpha _j^1=\mu ).$$
In particular, $`V_1`$ is the commutator of $`g`$ in $`\mathrm{M}_n`$. Now (5.3) implies
$`dim\mathrm{Gr}(m,n^2)dim\mathrm{Gr}(m,n^2)^g`$ $`dim(\mathrm{PGL}_n)dimC(g)`$
$`=n^2dim(V_1).`$
So if $`n3`$, part (b) of the following lemma gives the desired contradiction, which completes the proof of Theorems 5.1, and thus of Theorem 1.4 in the case that $`k`$ is algebraically closed. ∎
###### 5.4 Lemma.
Let $`n2`$, and $`2mn^22`$.
* $`dimV_1dimV_\mu `$ for any $`\mu 1`$.
* If $`n3`$ (or $`n=2`$ and there are more than two distinct nonzero $`V_\mu `$), then $`dim\mathrm{Gr}(m,n^2)dim\mathrm{Gr}(m,n^2)^g>n^2dim(V_1)`$.
###### Proof.
(a) Note that $`dimV_1=l_1^2+\mathrm{}+l_s^2`$ and
$$dimV_\mu =\underset{\lambda _i\lambda _j^1=\mu }{}l_il_j$$
Since the eigenvalues $`\lambda _1,\mathrm{},\lambda _s`$ of $`g`$ are distinct, the last sum has at most one term for each $`i=1,\mathrm{},s`$. Thus there is a permutation $`\sigma `$ of $`\{1,\mathrm{},s\}`$ such that
$$dimV_\mu l_1l_{\sigma (1)}+\mathrm{}+l_sl_{\sigma (s)}.$$
So for $`v=(l_1,\mathrm{},l_s)`$ and $`w=(l_{\sigma (1)},\mathrm{},l_{\sigma (s)})`$, $`dimV_\mu vw`$. Hence by the Cauchy-Schwarz inequality,
$$dimV_\mu vw|v||w|=|v|^2=l_1^2+\mathrm{}+l_s^2=dim(V_1).$$
(b) Since $`g`$ is semisimple, every $`L\mathrm{Gr}(m,n^2)^g`$ is a direct sum of its character subspaces spaces, i.e., a direct sum of vector subspaces of the $`V_\mu `$. Part (b) now follows from Proposition 4.1 with $`V=\mathrm{M}_n`$, $`N=n^2`$, $`N_\mu =dim(V_\mu )`$, $`Z=\mathrm{Gr}(m,n^2)^g`$, and $`r`$ the number of distinct nonzero $`V_\mu `$. ∎
###### 5.5 Remark.
We assumed throughout this section that $`n3`$. If $`n=2`$ then the above argument still goes through provided there are more than two distinct non-zero character subspaces $`V_\mu `$; see Lemma 5.4(b). In particular, this will always be the case if $`g^21`$ in $`\mathrm{PGL}_n`$; indeed, in this case $`g=(\lambda _1,\lambda _2)`$, where $`\mu =\lambda _1/\lambda _2\pm 1`$ and the three spaces $`V_1`$, $`V_\mu `$ and $`V_{\mu ^1}`$ are distinct. Thus the above argument also shows that for $`n=m=2`$, either $`|S|=1`$ or $`S`$ has exponent $`2`$. It turns out that, in fact, in this case $`|S|=2`$; see \[20, Lemma 14.2\].
###### 5.6 Remark.
An alternative approach to proving Theorem 5.1 would be to appeal to the classification, due to A. G. Elashvili and A. M. Popov , of pairs $`(G,\varphi )`$, where $`G`$ is a semisimple algebraic group and $`\varphi :G\mathrm{GL}(V)`$ is an irreducible linear representation of $`G`$ such that the $`G`$-action on $`V`$ has a non-trivial stabilizer in general position. Since this classification is rather involved, and since additional work would be required to apply it in our situation (note that the group $`\mathrm{GL}_m\times \mathrm{PGL}_n`$ is not semisimple, and that its representation on $`(\mathrm{M}_n)^m`$ is not irreducible), we opted instead for the self-contained direct proof presented in this section.
## 6. $`\mathrm{SL}_m`$-invariant generic matrices
The goal of this section is to relate the rings of $`\mathrm{SL}_m`$-invariants in $`G_{m,n}`$ and $`\mathrm{𝑈𝐷}(m,n)`$.
###### 6.1 Lemma.
* Every element of $`\mathrm{𝑈𝐷}(m,n)^{\mathrm{GL}_m}`$ can be written in the form $`\frac{a}{b}`$, where $`a`$ is a homogeneous element of $`(T_{m,n})^{\mathrm{SL}_m}`$, and $`b`$ is a non-zero homogeneous element of $`\mathrm{Z}(T_{m,n})^{\mathrm{SL}_m}`$ of the same degree as $`a`$.
* Assume that a subgroup $`G`$ of $`\mathrm{GL}_m`$ has no non-trivial characters. Then every element of $`\mathrm{𝑈𝐷}(m,n)^G`$ can be written as $`\frac{a}{b}`$ where $`a(T_{m,n})^G`$ and $`0b\mathrm{Z}(T_{m,n})^G`$.
###### Proof.
Both parts follows from Lemma 2.1. In part (a), we take $`\mathrm{\Gamma }=\mathrm{GL}_m\times \mathrm{PGL}_n`$, $`V=(\mathrm{M}_n)^m`$ (with the $`\mathrm{\Gamma }`$-action defined in the beginning of Section 5) and $`W=\mathrm{M}_n`$ (where $`\mathrm{GL}_m`$ acts trivially on $`W`$ and $`\mathrm{PGL}_n`$ acts by conjugation). Here the relative concomitants $`(\mathrm{M}_n)^m\mathrm{M}_n`$ are the homogeneous elements of $`(T_{m,n})^{\mathrm{SL}_m}`$ and the relative invariants $`(\mathrm{M}_n)^mk`$ are the homogeneous elements of $`\mathrm{Z}(T_{m,n})^{\mathrm{SL}_m}`$; cf. Theorem 2.2.
If $`G`$ has no non-trivial characters then relative concomitants are (absolute) concomitants, i.e., elements of $`(T_{m,n})^G`$. Similarly, relative invariants are elements of $`\mathrm{Z}(T_{m,n})^G`$, and part (b) is thus simply a restatement of Lemma 2.1 in this special case. ∎
###### 6.2 Proposition.
Let $`G`$ be a subgroup of $`\mathrm{GL}_m`$ such that $`G`$ has no non-trivial characters. Then the following conditions are equivalent:
* $`\mathrm{𝑈𝐷}(m,n)^G`$ has PI-degree $`n`$.
* $`(T_{m,n})^G`$ has PI-degree $`n`$.
* $`(G_{m,n})^G`$ has PI-degree $`n`$.
###### Proof.
The equivalence of (a) and (b) follows from Lemma 6.1(b). The implication (c) $``$ (b) is obvious, since $`G_{m,n}T_{m,n}`$. It thus remains to prove that (b) $``$ (c).
Let $`g_n`$ be the multilinear central polynomial for $`n\times n`$-matrices in \[12, 13.5.11\] (or \[23, p. 26\]). If $`R`$ is a prime PI-algebra of PI-degree $`n`$, denote by $`g_n(R)`$ the set of all evaluations of $`g_n`$ in $`R`$, and denote by $`Rg_n(R)`$ the nonzero ideal of $`R`$ generated by $`g_n(R)`$. Denote by $`T`$ the trace ring of $`R`$. (Since we are working in characteristic zero, the (noncommutative) trace ring in \[12, 13.9.2\] is the same as the one we are using, see \[12, 13.9.4\].) Then $`Rg_n(R)`$ is a common ideal of $`R`$ and $`T`$, see \[12, 13.9.6\] (or \[23, 4.3.1\]).
Now let $`R=G_{m,n}`$. Then its trace ring is $`T=T_{m,n}`$. Recall that we are assuming (b) holds, i.e., $`T^G`$ has PI-degree $`n`$. Let $`s`$ be a non-zero evaluation of $`g_n`$ on $`T^G`$. Then $`s`$ is a nonzero $`G`$-invariant, and a central element of $`T`$ (since it is also an evaluation of $`g_n`$ on $`T`$). Since $`g_n`$ is multilinear, and since $`T`$ is generated as an $`R`$-module by central elements, $`s`$ belongs to the ideal of $`T`$ generated by $`g_n(R)`$, so that $`sTRg_n(R)R`$. Since $`s`$ is a $`G`$-invariant, it follows that $`sT^GR^G`$. Consider the central localization $`R^G[s^1]\mathrm{𝑈𝐷}(m,n)`$. Since it contains $`T^G`$, $`R^G[s^1]`$ must have PI-degree $`n`$, implying that also $`R^G`$ must have PI-degree $`n`$. This completes the proof of the implication (b) $``$ (c) and thus of Proposition 6.2. ∎
###### 6.3 Remark.
The same argument also shows that if the three equivalent conditions in Proposition 6.2 are true, then the division algebras of fractions of $`(G_{m,n})^G`$ and $`(T_{m,n})^G`$ are both equal to $`\mathrm{𝑈𝐷}(m,n)^G`$.
## 7. Proof of Theorems 1.2 and 1.4
### Proof of Theorem 1.2
Proposition 6.2 tells us that $`(G_{m,n})^{\mathrm{SL}_m}`$ has PI-degree $`n`$ if and only if so does $`\mathrm{𝑈𝐷}(m,n)^{\mathrm{SL}_m}`$. Thus in order to prove Theorem 1.2 it suffices to show that $`\mathrm{𝑈𝐷}(m,n)^{\mathrm{SL}_m}`$ has PI-degree $`n`$ whenever $`2mn^22`$.
For $`n=m=2`$ we showed this in \[20, Remark 14.4\] (in fact, the argument we gave there remains valid over any base field $`k`$ of characteristic $`2`$). For $`n3`$ and $`2mn^22`$, Theorem 1.4 tells us that $`\mathrm{𝑈𝐷}(m,n)^{\mathrm{GL}_m}`$ has PI-degree $`n`$ (and consequently, so does $`\mathrm{𝑈𝐷}(m,n)^{\mathrm{SL}_m}`$). In summary, we have shown that Theorem 1.2 follows from Theorem 1.4. ∎
### Proof of Theorem 1.4
We have already proved Theorem 1.4 in the case where the base field $`k`$ is algebraically closed; see Section 5. We will now reduce the general case to this one by using Lemma 6.1(a).
We begin with a simple lemma.
###### 7.1 Lemma.
Let $`K`$ be an extension field of $`k`$, let $`V`$ be a finite-dimensional $`k`$-vector space, and $`V_K=V_kK`$. Given a linear representation of $`\mathrm{SL}_m(k)`$ on $`V`$, we have
$$(V_K)^{\mathrm{SL}_m(K)}=V^{\mathrm{SL}_m(k)}_kK.$$
###### Proof.
Since $`\mathrm{SL}_m(k)`$ is dense in $`\mathrm{SL}_m(K)`$, the subspace $`(V_K)^{\mathrm{SL}_m(K)}`$ is defined inside $`V_K`$ by a system of homogeneous linear equations with coefficients in $`k`$. Clearly finitely many of these equations suffice. Since the dimension of the solution space of such a system is the rank of the corresponding matrix (which has coefficients in $`k`$), $`(V_K)^{\mathrm{SL}_m(K)}`$ has a $`K`$-basis consisting of elements of $`V^{\mathrm{SL}_m(k)}`$. ∎
For the remainder of this section, we will write $`G_{m,n}(K)`$, $`T_{m,n}(K)`$ and $`\mathrm{𝑈𝐷}(m,n)(K)`$ to denote the generic matrix algebra, trace ring and universal division algebra defined over the field $`K`$. Denote the algebraic closure of $`k`$ by $`\overline{k}`$. Since the $`\mathrm{SL}_m`$-action on $`\mathrm{𝑈𝐷}(m,n)`$ preserves degree, Lemma 7.1 immediately implies the following fact, which we record for later use.
###### 7.2 Corollary.
$`G_{m,n}(\overline{k})^{\mathrm{SL}_m(\overline{k})}=G_{m,n}(k)^{\mathrm{SL}_m(k)}_k\overline{k}`$, and $`T_{m,n}(\overline{k})^{\mathrm{SL}_m(\overline{k})}=T_{m,n}(k)^{\mathrm{SL}_m(k)}_k\overline{k}`$. ∎
We are now ready to complete the proof of Theorem 1.4 over an arbitrary field $`k`$ of characteristic zero. In Section 5 we showed that Theorem 1.4 holds over the algebraic closure $`\overline{k}`$ of $`k`$. That is, if $`2mn^22`$ then there exist elements $`c_1,\mathrm{},c_r\mathrm{𝑈𝐷}(m,n)(\overline{k})^{\mathrm{GL}_m(\overline{k})}`$ which span $`\mathrm{𝑈𝐷}(m,n)(\overline{k})`$ as a vector space over its center. By Lemma 6.1 we can write $`c_i=a_i/b_i`$, where $`a_iT_{m,n}(\overline{k})[d_i]^{\mathrm{SL}_m}`$ and $`0b_i\mathrm{Z}(T_{m,n}(\overline{k}))[d_i]^{\mathrm{SL}_m}`$ for some $`d_i0`$, $`i=1,\mathrm{},r`$. By Lemma 7.1, with $`K=\overline{k}`$ and $`V=\mathrm{Z}(T_{m,n}(k))[d_i]`$, we have $`\mathrm{Z}(T_{m,n}(k))[d_i]^{\mathrm{SL}_m}0`$. We may now replace $`b_i`$ by a non-zero element of $`\mathrm{Z}(T_{m,n}(k))[d_i]^{\mathrm{SL}_m}`$. The new $`c_i=a_i/b_i`$ are still $`\mathrm{GL}_m`$-invariant elements of $`\mathrm{𝑈𝐷}(m,n)(\overline{k})`$, and they still generate $`\mathrm{𝑈𝐷}(m,n)(\overline{k})`$ as a vector space over its center.
We now apply Lemma 7.1 once again (this time with $`V=T_{m,n}(k)[d_i]`$) to write each $`a_i`$ as a finite sum $`\gamma _{ij}a_{ij}`$, where each $`\gamma _{ij}\overline{k}`$ and each $`a_{ij}T_{m,n}(k)[d_i]^{\mathrm{SL}_m}`$. Now replace our collection of $`\mathrm{GL}_m`$-invariant elements $`\{c_i=a_i/b_i\}`$ in $`\mathrm{𝑈𝐷}(m,n)(\overline{k})`$ by $`\{c_{ij}=a_{ij}/b_i\}`$. By construction, the elements $`c_{ij}`$ lie in $`\mathrm{𝑈𝐷}(m,n)(k)^{\mathrm{GL}_m}`$ and span $`\mathrm{𝑈𝐷}(m,n)(\overline{k})`$ over its center. Hence, these elements generate a $`k`$-subalgebra of $`\mathrm{𝑈𝐷}(m,n)(k)^{\mathrm{GL}_m}`$ of PI-degree $`n`$. Consequently, $`\mathrm{𝑈𝐷}(m,n)(k)^{\mathrm{GL}_m}`$ itself has PI-degree $`n`$. This completes the proof of Theorem 1.4 (and of Theorem 1.2). ∎
## 8. The case $`mn^21`$
Theorems 1.2 and 1.4 assume that $`mn^22`$. We will now describe $`\mathrm{𝑈𝐷}(m,n)^{\mathrm{GL}_m}`$ and $`\mathrm{𝑈𝐷}(m,n)^{\mathrm{SL}_m}`$ for $`mn^21`$.
Recall the definition of the discriminant of $`n^2`$ matrices of size $`n\times n`$, say $`A_1,\mathrm{},A_{n^2}`$: it is the determinant of the $`n^2\times n^2`$-matrix whose $`i`$-th row consists of the entries of $`A_i`$, cf. (8.5). When viewed as a function $`(\mathrm{M}_n)^{n^2}k`$, $`\mathrm{\Delta }`$ is the unique multilinear alternating function such that $`\mathrm{\Delta }(e_{11},e_{12},\mathrm{},e_{nn})=1`$; cf., e.g., \[10, Lemma 3\]. Here the $`e_{ij}`$ are the matrix units.
###### 8.1 Proposition.
(a) If $`m>n^2`$, then $`\mathrm{𝑈𝐷}(m,n)^{\mathrm{SL}_m}=\mathrm{𝑈𝐷}(m,n)^{\mathrm{GL}_m}=k`$.
Now let $`m=n^2`$, and denote by $`\mathrm{\Delta }`$ the discriminant of the generic matrices $`X_1,\mathrm{},X_m`$.
* $`\mathrm{𝑈𝐷}(m,n)^{\mathrm{GL}_m}=k`$.
* $`(T_{m,n})^{\mathrm{SL}_m}=k[\mathrm{\Delta }]`$.
* $`\mathrm{𝑈𝐷}(m,n)^{\mathrm{SL}_m}=k(\mathrm{\Delta })`$.
###### Proof.
(a) We may clearly assume that $`k`$ is algebraically closed. In this case $`\mathrm{SL}_m`$ has a dense orbit in the associated variety $`X=(\mathrm{M}_n)^m`$. Thus the rational quotient $`X/\mathrm{SL}_m`$ is a single point (with trivial $`\mathrm{PGL}_n`$-action), and
$$\mathrm{𝑈𝐷}(m,n)^{\mathrm{SL}_m}=\mathrm{RMaps}_{\mathrm{PGL}_n}(pt,\mathrm{M}_n)=k.$$
Now suppose $`m=n^2`$. Then $`\mathrm{GL}_m`$ has a dense orbit in $`X=(\mathrm{M}_n)^m`$. Arguing as in part (a), we prove (b); cf. \[20, Proposition 13.1(a)\]. (c) is proved in \[9, p. 210\], and (d) follows from (c) by Lemma 6.1(b). ∎
###### 8.2 Remark.
Let $`m=n^2`$. Formanek showed that $`\mathrm{\Delta }(G_{m,n})^{\mathrm{SL}_m}`$ (\[9, p. 214\]) but $`\mathrm{\Delta }^iG_{m,n}`$ for every integer $`i2`$ (this follows from \[10, Theorem 16\]). Consequently for $`m=n^2`$,
$$(G_{m,n})^{\mathrm{SL}_m}=k[\mathrm{\Delta }^2,\mathrm{\Delta }^3].$$
###### 8.3 Proposition.
Suppose $`m=n^21`$, and let
$$Y=\underset{i,j=1}{\overset{n}{}}\mathrm{\Delta }(X_1,\mathrm{},X_m,e_{ji})e_{ij},$$
where the $`e_{ij}`$ are the matrix units.
* $`Y(T_{m,n})^{\mathrm{SL}_m}`$.
* The eigenvalues of $`Y`$ are algebraically independent over $`k`$ (and, in particular, distinct).
* $`(T_{m,n})^{\mathrm{SL}_m}=k[c_1,\mathrm{},c_{n1},Y]`$ is a polynomial ring in $`n`$ independent variables over $`k`$. Here $`c_1=\mathrm{tr}(Y),\mathrm{},c_n=(1)^ndet(Y)`$.
* $`\mathrm{𝑈𝐷}(m,n)^{\mathrm{SL}_m}=k(c_1,\mathrm{},c_{n1},Y)`$.
* $`\mathrm{𝑈𝐷}(m,n)^{\mathrm{GL}_m}=k(\frac{c_2}{(c_1)^2},\mathrm{},\frac{c_{n1}}{(c_1)^{n1}},\frac{1}{c_1}Y)`$.
###### Proof.
For the proof of (a)—(c), we may assume that $`k`$ is algebraically closed (cf. Corollary 7.2). (a) We view $`Y`$ as a regular map $`(\mathrm{M}_n)^m\mathrm{M}_n`$. We want to show that this map is $`\mathrm{PGL}_n`$-equivariant, i.e., $`YT_{m,n}`$. Since $`Y`$ is clearly $`\mathrm{SL}_m`$-equivariant (recall that $`\mathrm{SL}_m`$ acts trivially on $`\mathrm{M}_n`$), this will imply part (a).
We begin by observing that for any $`(A_1,\mathrm{},A_m)(\mathrm{M}_n)^m`$, and $`Z\mathrm{M}_n`$,
(8.4)
$$\mathrm{tr}(Y(A_1,\mathrm{},A_m)Z)=\mathrm{\Delta }(A_1,\mathrm{},A_m,Z).$$
Indeed, both sides are linear in $`Z`$, so we only need to check (8.4) for the elementary matrices $`Z=e_{ij}`$, where it is easy to do directly from the definition of $`Y`$.
Fix an $`m`$-tuple $`(A_1,\mathrm{},A_m)(\mathrm{M}_n)^m`$ of $`n\times n`$-matrices. Since the trace form on $`\mathrm{M}_n`$ is non-singular, $`Y(A_1,\mathrm{},A_m)`$ is the unique matrix satisfying (8.4) for every $`Z\mathrm{M}_n`$. The $`\mathrm{PGL}_n`$-equivariance of $`Y:(\mathrm{M}_n)^m\mathrm{M}_n`$ is an easy consequence of this and the fact that $`\mathrm{\Delta }`$ is $`\mathrm{PGL}_n`$-invariant (see \[9, p. 209\]). This concludes the proof of (a).
Our proof of parts (b) and (c) relies on the following claim: $`Y:(\mathrm{M}_n)^m\mathrm{M}_n`$ is the categorical quotient map for the $`\mathrm{SL}_m`$-action on $`(\mathrm{M}_n)^m`$. In other words, we claim that the $`n^2`$ elements $`\mathrm{\Delta }(X_1,\mathrm{},X_m,e_{ij})`$ ($`i,j=1,\mathrm{},n`$) generate $`k[(\mathrm{M}_n)^m]^{\mathrm{SL}_m}`$ as a $`k`$-algebra. To prove this claim we will temporarily write $`(A_1,\mathrm{},A_m)(\mathrm{M}_n)^m`$ as an $`m\times n^2`$-matrix
(8.5)
$$A=\left(\begin{array}{cccccc}a_{11}^{(1)}& a_{12}^{(1)}& \mathrm{}& a_{ij}^{(1)}& \mathrm{}& a_{nn}^{(1)}\\ \mathrm{}& \mathrm{}& & \mathrm{}& & \mathrm{}\\ a_{11}^{(m)}& a_{12}^{(m)}& \mathrm{}& a_{ij}^{(m)}& \mathrm{}& a_{nn}^{(m)}\end{array}\right).$$
That is, we write each $`n\times n`$ matrix $`A_h=(a_{ij}^{(h)})`$ as a single row of $`A`$. In this notation, $`g\mathrm{SL}_m`$ acts on $`(\mathrm{M}_n)^m`$ by multiplication by the transpose of $`g^1`$ on the left; that is, $`g(A)=(g^1)^{\text{transpose}}A`$ for every $`g\mathrm{SL}_m`$. Let $`\delta _{ij}(A_1,\mathrm{},A_m)`$ be the $`m\times m`$-minor of this matrix obtained by removing the $`ij`$-column from $`A`$ and taking the determinant of the resulting $`m\times m`$-matrix. The first theorem of classical invariant theory (see \[27, Theorem II.6.A\] or \[6, Theorem 2.1\]) says that the elements $`\delta _{ij}(X_1,\mathrm{},X_m)`$ generate $`k[(\mathrm{M}_n)^m]^{\mathrm{SL}_m}`$ as a $`k`$-algebra. On the other hand, it is easy to see that $`\delta _{ij}(X_1,\mathrm{},X_m)=\pm \mathrm{\Delta }(X_1,\mathrm{},X_m,e_{ij})`$. This proves the claim.
Now observe that since $`m=n^21`$,
$$dim((\mathrm{M}_n)^m//\mathrm{SL}_m)=mn^2(m^21)=n^2=dim(\mathrm{M}_n).$$
This means that the $`n^2`$ $`\mathrm{SL}_m`$-invariant functions
$$\mathrm{\Delta }(X_1,\mathrm{},X_m,e_{ji}):(\mathrm{M}_n)^mk$$
are algebraically independent over $`k`$. In other words, $`Y`$ (viewed as a matrix in $`T_{m,n}\mathrm{M}_n(k[x_{ij}^{(h)}]`$) has algebraically independent entries. Part (b) easily follows from this assertion; cf. \[15, Lemma II.1.4\].
Furthermore,
$`(T_{m,n})^{\mathrm{SL}_m}`$ $`=\mathrm{Morph}_{\mathrm{SL}_m\times \mathrm{PGL}_n}((\mathrm{M}_n)^m,\mathrm{M}_n)`$
$`\mathrm{Morph}_{\mathrm{PGL}_n}((\mathrm{M}_n)^m//\mathrm{SL}_m,\mathrm{M}_n)`$
$`\mathrm{Morph}_{\mathrm{PGL}_n}(\mathrm{M}_n,\mathrm{M}_n)=T_{1,n},`$
where $`T_{1,n}`$ is the trace ring of one generic $`n\times n`$-matrix. Here the last equality is a special case of Procesi’s Theorem 2.2(a) (with $`m=1`$). Since the chain of isomorphisms identifies $`Y`$ with the identity map $`\mathrm{M}_n\mathrm{M}_n`$, we conclude that
$$(T_{m,n})^{\mathrm{SL}_m}=k[c_1,\mathrm{},c_n,Y].$$
Since $`Y^n+c_1Y^{n1}+\mathrm{}+c_n=0`$, $`k[c_1,\mathrm{},c_{n1},Y]=k[c_1,\mathrm{},c_n,Y]`$. This proves the first assertion in (c).
To show that $`c_1,\mathrm{},c_{n1},Y`$ are algebraically independent over $`k`$, denote the eigenvalues of $`Y`$ by $`\lambda _1,\mathrm{},\lambda _n`$. By part (b), $`\lambda _1,\mathrm{},\lambda _n`$ are algebraically independent over $`k`$. Since $`Y`$ is algebraic over $`k(c_1,\mathrm{},c_n)`$, we have
$`\mathrm{trdeg}_k`$ $`k(c_1,\mathrm{},c_{n1},Y)=\mathrm{trdeg}_kk(c_1,\mathrm{},c_n,Y)`$
$`=\mathrm{trdeg}_kk(c_1,\mathrm{},c_n)=\mathrm{trdeg}_kk(\lambda _1,\mathrm{},\lambda _n)=n.`$
This shows that $`c_1,\mathrm{},c_{n1},Y`$ are algebraically independent over $`k`$, thus completing the proof of (c).
(d) is an immediate consequence of (c) and Lemma 6.1(b). To prove (e), denote the central torus of $`\mathrm{GL}_m`$ by $`𝔾_m`$. Then
$$\mathrm{𝑈𝐷}(m,n)^{\mathrm{GL}_m}=(\mathrm{𝑈𝐷}(m,n)^{\mathrm{SL}_m})^{𝔾_m}=k(c_1,\mathrm{},c_{n1},Y)^{𝔾_m},$$
where $`𝔾_m`$ acts on the purely transcendental extension $`k(c_1,\mathrm{},c_{n1},Y)`$ as follows: $`tc_it^{im}c_i`$ for $`i=1,\mathrm{},n1`$, and and $`tYt^mY`$. Part (e) easily follows from this description. ∎
###### 8.6 Remark.
Note that $`c_1=\mathrm{\Delta }(X_1,\mathrm{},X_{n^21},I_n)`$, where $`I_n`$ is the $`n\times n`$ identity matrix. By a theorem of Formanek, $`(c_1)^2`$ is an element of $`\mathrm{Z}(G_{m,n})^{\mathrm{SL}_m}`$ for $`m=n^21`$, see \[10, Theorem 16\].
## 9. Proof of Theorem 1.3
By Corollary 7.2, we may assume that $`k`$ is algebraically closed. Set $`A=(G_{m,n})^{\mathrm{SL}_m}`$ and $`B=(T_{m,n})^{\mathrm{SL}_m}`$. By Theorem 1.2, $`A`$ and $`B`$ both have PI-degree $`n`$. Thus $`\mathrm{Z}(A)=(\mathrm{Z}(G_{m,n}))^{\mathrm{SL}_m}`$ and $`\mathrm{Z}(B)=(\mathrm{Z}(T_{m,n}))^{\mathrm{SL}_m}`$. Since $`\mathrm{SL}_m`$ is a reductive group, and since $`T_{m,n}`$ is a finitely generated $`k`$-algebra and a finite module over its center, $`B`$ is a finite $`\mathrm{Z}(B)`$-module, and both $`B`$ and $`\mathrm{Z}(B)`$ are finitely generated Noetherian $`k`$-algebras, cf. \[26, Proposition 4.2\]. Moreover, $`B`$ is an FBN ring, cf. \[12, 13.6.6\].
By Corollary 3.3(b) and Remark 6.3, the transcendence degrees of both $`B`$ and $`\mathrm{Z}(B)`$ are $`t=(m1)n^2m^2+2`$. For notational simplicity, set
$$\mu (S)=\underset{d\mathrm{}}{lim\; sup}\frac{dim_kS[d]}{d^{t1}}$$
for any graded $`k`$-algebra $`S=_{d0}S[d]`$. By \[25, Lemma 6.1\] (cf. also \[11, §12.6\]), $`f(d)=dim_kB[d]`$ is eventually periodically polynomial, i.e., there are polynomials $`f_1,\mathrm{},f_s`$ with rational coefficients such that $`f(d)=f_i(d)`$ whenever $`d`$ is large enough and congruent to $`i`$ modulo $`s`$; moreover, the maximum of the degrees of the $`f_i`$ is $`t1`$. Consequently $`\mu (B)`$ exists and is equal to the largest among the leading coefficients of those $`f_i`$ of degree $`t1`$. A similar argument shows that $`\mu (Z(B))`$ exists and is a nonzero number.
Consider the multilinear central polynomial $`g_n`$ for $`n\times n`$ matrices used in the proof of Proposition 6.2. Since it is multilinear and nonzero on $`A`$, we can find a nonzero evaluation $`c`$ of $`g_n`$ at homogeneous elements of $`A`$; this $`c`$ is homogeneous. Since $`c`$ is also an evaluation of $`g_n`$ on $`G_{m,n}`$, $`cT_{m,n}G_{m,n}`$, so that $`cBA`$ and $`c\mathrm{Z}(B)\mathrm{Z}(A)`$. Then for all integers $`dj`$, $`cB[dj]A[d]B[d]`$, where $`j=\mathrm{deg}c`$. Replacing $`c`$ by $`c^s`$ if necessary, we may assume that $`s`$ divides $`j`$. Consequently, whenever $`d`$ is large enough and congruent to $`i`$ modulo $`s`$,
$$f_i(dj)dim_kA[d]f_i(d).$$
It follows easily that $`\mu (A)`$ exists and is equal to the largest among the leading coefficients of those $`f_i`$ of degree $`t1`$. A similar argument shows that $`\mu (Z(A))`$ exists and is a nonzero number. ∎
###### 9.1 Remark.
The above proof shows that $`\mu ((G_{m,n})^{\mathrm{SL}_m})=\mu ((T_{m,n})^{\mathrm{SL}_m})`$, and that $`\mu (\mathrm{Z}((G_{m,n})^{\mathrm{SL}_m}))=\mu (\mathrm{Z}((T_{m,n})^{\mathrm{SL}_m}))`$.
###### 9.2 Remark.
Consider the $`\mathrm{GL}_m`$-representation on $`R`$, where $`R=G_{m,n}`$, $`T_{m,n}`$, $`\mathrm{Z}(G_{m,n})`$ or $`\mathrm{Z}(T_{m,n})`$. Recall that irreducible polynomial representations of $`\mathrm{GL}_m`$ are indexed by partitions $`\lambda =(\lambda _1,\mathrm{},\lambda _s)`$ with $`sm`$ parts; cf. \[9, Section 2\]. Denote the multiplicity of the irreducible representation corresponding to $`\lambda `$ in $`R`$ by $`\mathrm{mult}_\lambda (R)`$. Writing $`(r^m)`$ for the partition $`\lambda =(r,\mathrm{},r)`$ ($`m`$ times), we have
(a) $`dim_kR^{\mathrm{SL}_m}[d]=0`$ if $`d`$ is not a multiple of $`m`$, and
(b) $`dim_kR^{\mathrm{SL}_m}[rm]=\mathrm{mult}_{(r^m)}(R)`$ for any integer $`r1`$.
###### Proof.
(a) Assume $`R^{\mathrm{SL}_m}[d]`$ is nonzero. Then it is a direct sum of one-dimensional representations of $`\mathrm{GL}_m`$ of the form $`M=\mathrm{Span}_k(v)`$. Moreover, any such representation is given by $`g(v)=det(g)^rv`$ for some integer $`r`$; cf., e.g., \[9, Theorem 3(a)\]. On the other hand, substituting $`g=tI_m`$, where $`tk`$ and $`I_m`$ is the $`m\times m`$ identity matrix, we obtain, $`g(v)=t^dv`$. Since $`det(tI_m)=t^m`$, we see that $`d=rm`$, as claimed.
(b) If $`d=rm`$ and $`0vR^{\mathrm{SL}_m}[rm]`$ then the partition associated to the 1-dimensional irreducible $`\mathrm{GL}_m`$-module $`M=\mathrm{Span}(v)`$ is $`(r^m)`$; cf. , e.g., \[9, Theorem 2\]. Now consider the direct sum decomposition $`R=R_\lambda `$, where $`R_\lambda `$ is the sum of all irreducible $`\mathrm{GL}_m`$-submodules of $`R`$ with associated partition $`\lambda `$. The argument of part (a) shows that $`R_{(r^m)}=R^{\mathrm{SL}_m}[rm]`$. Moreover, since $`dim(M)=1`$, we have
$$dim_kR^{\mathrm{SL}_m}[rm]=dim_kR_{(r^m)}=\mathrm{mult}_{(r^m)}(R),$$
as claimed. ∎
## 10. Standard polynomials
Let $`G_{m,n}`$ be the ring of $`m`$ generic $`n\times n`$-matrices. By Theorem 1.2, $`(G_{m,n})^{\mathrm{SL}_m}`$ is a PI domain of degree $`n`$, whenever $`2mn^22`$. We will now describe one particular element of this ring. Let
$$F_m(x_1,\mathrm{},x_m)=\underset{\sigma \mathrm{S}_n}{}(1)^\sigma x_{\sigma (1)}\mathrm{}x_{\sigma (n)}k\{x_1,\mathrm{},x_m\}$$
be the standard polynomial. Since $`F_m`$ is multilinear and alternating, one checks easily that for $`g\mathrm{GL}_m`$,
(10.1)
$$g(F_m)=det(g)F_m;$$
see, e.g., \[23, 1.4.12\]. Substituting $`m`$ generic $`n\times n`$-matrices $`X_1,\mathrm{},X_m`$ into $`F_m`$, we obtain
$$f_{m,n}=F_m(X_1,\mathrm{},X_m)=\underset{\sigma \mathrm{S}_n}{}(1)^\sigma X_{\sigma (1)}\mathrm{}X_{\sigma (n)}G_{m,n}.$$
From (10.1), we see that $`f_{m,n}(G_{m,n})^{\mathrm{SL}_m}`$. By the Amitsur-Levitzki Theorem, $`f_{m,n}=0`$ iff $`m2n`$.
Fix $`m,n2`$ and let $`K`$ be the center of $`\mathrm{𝑈𝐷}(m,n)`$.
###### 10.2 Proposition.
For $`2m<2n`$, $`K(f_{m,n})`$ generates a $`\mathrm{GL}_m`$-stable maximal subfield of $`\mathrm{𝑈𝐷}(m,n)`$.
The proof is algebraic in nature and works in characteristic $`2`$.
###### Proof.
The fact that $`K(f_{m,n})`$ is a $`\mathrm{GL}_m`$-stable subfield follows from (10.1). In order to prove that this subfield is maximal, it suffices to verify that $`f_{m,n}`$ has an eigenvalue of multiplicity $`1`$. (Indeed, if, say, $`n=d[K(f_{m,n}):K]`$, then the characteristic polynomial $`p(t)`$ of $`f_{m,n}`$ in $`\mathrm{𝑈𝐷}(m,n)`$ has the form $`p(t)=q(t)^d`$, where $`q(t)`$ is the minimal polynomial of $`f_{m,n}`$ over $`K`$. This shows that the multiplicity of each eigenvalue of $`f_{m,n}`$ is divisible by $`d`$.)
Since the multiplicity of eigenvalues cannot decrease when evaluating $`f_{m,n}`$ in $`\mathrm{M}_n`$, it suffices now to show that $`f_{m,n}`$ (or equivalently, $`F_m`$) has some evaluation in $`\mathrm{M}_n`$ with an eigenvalue of multiplicity one. We now proceed to construct such an evaluation. Since
$$F_m(1,x_2,\mathrm{},x_m)=F_{m1}(x_2,\mathrm{},x_m)$$
for $`m`$ odd (cf. \[23, Exercise 1.2.3\]), we may assume that $`m`$ is even, say $`m=2r2`$, with $`1<rn`$. In $`\mathrm{M}_n`$, consider the sequence of $`m`$ matrix units
$$e_{1,2},e_{2,2},e_{2,3},e_{3,3},\mathrm{},e_{r2,r1},e_{r1,r1},e_{r1,r},e_{r,1}.$$
When permuting these matrix units cyclically, their product is nonzero; for any other permutation, their product is zero. Since an $`m`$-cycle is odd, it follows that $`F_m`$ evaluated at these matrix units is
$$e_{1,1}e_{2,2}+e_{2,2}+\mathrm{}e_{r1,r1}+e_{r1,r1}e_{r,r}=e_{1,1}e_{r,r},$$
which has $`1`$ as an eigenvalue of multiplicity one (since $`\mathrm{char}(k)2`$). ∎
We do not know an explicit expression for any non-constant element of $`(G_{m,n})^{\mathrm{SL}_m}`$ (as a polynomial in the generic $`n\times n`$-matrices $`X_1,\mathrm{},X_m`$) in the case where $`2nmn^22`$; we leave this as an open question. Note that for $`m=n^2`$ and $`m=n^21`$, such elements are exhibited in Remarks 8.2 and 8.6.
## Acknowledgments
We are grateful to A. Berele and E. Formanek for helpful comments. |
warning/0507/gr-qc0507065.html | ar5iv | text | # Negative Energies and a Constantly Accelerating Flat Universe
## I Introduction
One of the most challenging tasks in contemporary physics is to understand the observational results indicating that we are living in a flat accelerating universe. The most popular interpretation is that we are dominated by a homogeneous component with negative pressure often called dark energy. The supernovae data from sn indicate that the equation of state of this dark energy is compatible with the ’concordance model’ ($`\mathrm{\Omega }_M=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }`$ = 0.7, w=p/$`\rho `$=-1) but a careful interpretation of the data nous shows that a component with w $`<`$ -1 is allowed. Models with very exotic $`w(z)`$ may come from modified gravity and have very different consequences for the fate of the UniverseLinderG . Such models are quite unsatisfactory since they inevitably lead to a generic violation of positive energy conditions resulting in vacuum quantum instabilities Lin1 Cald Fram Carr .
However, alternative proposals have been made fhc2 Lin2 where the local instability issue can be avoided because the interaction between the positive and negative energy universes is global. The model presented in fhc2 is particularly attractive since the conjugate universe and its negative energy content are not introduced by hand but emerge as a result of imposing new symmetries on the initial action. It is also very predictive: flatness is mandatory and a constant acceleration phase is one of the very few mathematical possibilities in such a tightly constrained theoretical framework. We confront in this paper its predictions to present observationnal data.
## II Motivations for a modified GR
Investigation of negative energies in Relativistic Quantum Field Theory (QFT) indicates that the correct theoretical framework should be found in a modification of General Relativity (GR) fred .
\- TheoreticaI motivations
In second quantification, all relativistic field equations admit genuine negative energy field solutions creating and annihilating negative energy quanta. Unitary time reversal links these fields to the positive energy ones. The unitary choice, usual for all other symmetries in physics, also allows us to avoid the well known paradoxes associated with time reversal. Positive and negative energy fields vacuum divergences encountered after second quantization, are unsurprisingly found to be exactly opposite. The negative energy fields action must be maximised. However, there is no way to reach a coherent theory involving negative energies in flat-spacetime. Indeed, if positive and negative energy scalar fields are time reversal conjugate, so must be their Hamiltonian densities and actions. This is only possible in the context of GR thanks to the metric transformation under discrete symmetries.
\- Phenomenological motivations
In a mirror negative energy world, whose fields remain non coupled to our world positive energy fields, stability is insured and the behavior of matter and radiation is as usual. Hence, it is just a matter of convention to define each one as a positive or negative energy world. Otherwise, if they interact gravitationally, promising phenomenology is expected. Indeed, many outstanding enigmas indicate that repelling gravity might play an important role in physics: flat galactic rotation curves, the Pioneer effect, the flatness of the universe, acceleration and its voids, etc… But negative energy states never manifested themselves up to now, suggesting that a barrier is at work preventing the two worlds to interact except through gravity.
\- A modified GR to circumvent the main issues
A trivial cancellation between vacuum divergences is not acceptable since the Casimir effect shows evidence for vacuum fluctuations. But the positive and negative energy worlds could be maximally gravitationally coupled in such a way as to produce at least exact cancellations of vacuum energies gravitational effects. Also, a generic catastrophic instability issue arises whenever quantum positive and negative energy fields are allowed to interact. If we restrict the stability issue to the modified gravity, this disastrous scenario is avoided. Finally, allowing both positive and negative energy virtual photons to propagate the electromagnetic interaction, simply makes it disappear. The local gravitational interaction is treated very differently in our modified GR .So this unpleasant feature is also avoided.
## III Conjugate Worlds Gravitational Coupling
Ref. Wein shows that time reversal does not affect a scalar action. However, if the inertial coordinates $`\xi ^\alpha `$ are transformed in a non-trivial way:
$$\xi ^\alpha \text{ }\stackrel{T}{}\stackrel{~}{\xi }_T^\alpha $$
(1)
where $`\stackrel{~}{\xi }^\alpha \xi _\alpha `$, metric terms are affected and the action is not expected to be invariant under $`T`$. Having two conjugate inertial coordinate systems, two time reversal conjugate metric tensors can be built:
$$g_{\mu \nu }=\eta _{\alpha \beta }\frac{\xi ^\alpha }{x^\mu }\frac{\xi ^\beta }{x^\nu },\stackrel{~}{g}_{\mu \nu }=\eta _{\alpha \beta }\frac{\stackrel{~}{\xi }^\alpha }{x^\mu }\frac{\stackrel{~}{\xi }^\beta }{x^\nu }$$
(2)
Then, a new set of fields couples to the new $`\stackrel{~}{g}_{\mu \nu }`$ metric field. The total action is the sum of $`I_M`$, the usual action for matter and radiation in the external gravitational field $`g_{\mu \nu }`$, $`\stackrel{~}{I}_M`$ the action for matter and radiation in the external gravitational field $`\stackrel{~}{g}_{\mu \nu }`$ and the actions $`I_G+\stackrel{~}{I}_G`$ for the gravitational fields alone. The conjugate actions are separately general coordinate scalars and adding the two pieces is necessary to obtain a discrete symmetry reversal invariant total action. $`g_{\mu \nu }`$ and $`\stackrel{~}{g}_{\mu \nu }`$ are linked since these are symmetry reversal conjugate objects, explicitly built out of space-time coordinates. We postulate that there exists a privileged general coordinate system such that $`\stackrel{~}{g}_{\mu \nu }`$ identifies with $`g^{\mu \nu }`$, where for instance a discrete time reversal transformation applies as $`x^0x^0`$. In this system, varying the action, applying the extremum action principle and making use of the relation $`\delta g^{\rho \kappa }\left(x\right)=g^{\rho \mu }\left(x\right)g^{\nu \kappa }\left(x\right)\delta g_{\mu \nu }\left(x\right)`$ leads to a modified Einstein equation only valid in the privileged coordinate system. This equation is not general covariant and not intended to be so.
## IV Modified cosmology
Following the method outlined in the previous section, local solutions satisfying the symmetry invariance requirements under Parity or space/time exchange transformations have been found and interpreted in fhc2 . In the case of cosmology and time reversal, there exists one global privileged coordinate system where a couple of purely time dependent time reversal conjugate solutions can be derived from the couple of conjugate actions. The existence of a time reversal conjugate universe was also suggested a long time ago in Ref. Sak . The only possible privileged coordinate system where both metrics are spatially homogeneous and isotropic is the flat Cartesian one:
$$d\tau ^2=B(t)dt^2A(t)dx^2$$
(3)
The privileged coordinate system is the conformal time system where $`B=A`$. Expressing in the polar coordinate system, the modified cosmological Einstein equations are:
$$3A\left(\frac{\ddot{A}}{A}+\frac{1}{2}\left(\frac{\dot{A}}{A}\right)^2\right)\frac{3}{A}\left(\frac{\ddot{A}}{A}\frac{3}{2}\left(\frac{\dot{A}}{A}\right)^2\right)=0$$
(4)
The purely time dependent scale factor evolution is then driven (a nondimensional time unit is used) by the following differential equations in the three particular domains:
$$\text{ }a<<1\ddot{a}\frac{3}{2}\frac{\dot{a}^2}{a}a1/t^2\text{ where }t<0,$$
(5)
$$\text{ }a1\ddot{a}\frac{\dot{a}^2}{a}ae^t,$$
(6)
$$\text{ }a>>1\ddot{a}\frac{1}{2}\frac{\dot{a}^2}{a}at^2\text{ where }t\text{ }>\text{ 0}\text{}.$$
(7)
We check that $`tt`$ implies $`1/t^2t^2`$ but also $`e^te^t`$, as required. Flatness is the main prediction of this model. A striking and very uncommon feature is that the evolution of the scale factor is completely independent of the matter and radiation content in the two universes. In particular, the observed flatness can no longer be translated into the usual estimation of $`\mathrm{\Omega }_m=1`$ from the WMAP data WMAP . According the $`t^2`$ evolution, we are most probably living in a constantly accelerating universe. Our and the conjugate universe either have crossed in the past or will cross each other and time reversal will occur in the future. Our universe is accelerated without any need for a cosmological constant or dark energy component No source enters our cosmological equation except at the crossing time where a small perturbation is needed to start the non stationary cosmological solutions.
## V Application to supernovae data
We have confronted this model with the existing supernovae data published by sn . The luminosity distance, for a flat constantly accelerated ($`t^2`$ evolution) universe is:
$$d_L=a_0r_l(1+z)=\frac{2}{H_0}(\sqrt{1+z}1)(1+z)$$
(8)
Figure 1 shows the fitted magnitude versus redshift curve from the model, where the only free parameter is the normalisation parameter $`\mathrm{m}_\mathrm{s}`$, compared to the observational data derived from the gold sample of sn . The quality of the fit is estimated with the computed $`\chi ^2`$ to be at a $`2\%`$ confidence level (CL). This should be compared to the $`11\%`$ confidence level found by using the standard $`\mathrm{\Lambda }\mathrm{CDM}`$ model with no prior. Fitting our model on the SCP data scp leads to a $`59\%`$ confidence level compatibility, to be compared with a $`56\%`$ confidence level using the $`\mathrm{\Lambda }\mathrm{CDM}`$ model.
To go further in this analysis, we parametrise the scale factor evolution as a simple power low $`\mathrm{a}\mathrm{t}^\alpha `$. The luminosity distance reads:
$$d_L=\frac{1}{H_0}\frac{\alpha }{\alpha 1}((1+z)^{11/\alpha }1)(1+z)$$
(9)
Fitting simultaniously $`\mathrm{m}_\mathrm{s}`$ and $`\alpha `$ on the Riess gold data sample gives $`\alpha =1.41\pm 0.13`$ which is $`4.5\sigma `$ away from the predicted value. Assuming a simple variation of the intrinsic supernovae magnitude versus the redshift $`\mathrm{m}_\mathrm{s}/\mathrm{z}=0.15`$ which is below the current statistical error, leads to a systematical error : $`+0.50.3(syst.)`$ for this model, both for Riess sn and SCP scp data. This shows that the power law parametrisation is much more sensitive to systematical errors than the standard fitting procedure including evolution of the equation of state sn . In conclusion, the current precision is not sufficient to discriminate between the $`\mathrm{\Lambda }\mathrm{CDM}`$ model and the model presented here.
We have then investigated the expected sensitivity with future SN projects. We simulate the SNLS snls1 snls2 experiment expecting about 700 SNIa with redshifts up to about 1, adding 300 simulated nearby SNIa from the future SN factory project and using $`\mathrm{\Lambda }\mathrm{CDM}`$ as fiducial model. We get $`\alpha =1.26\pm 0.04(\mathrm{stat}.)\pm 0.3(\mathrm{syst}.)`$, where the systematical error has been evaluated using a $`10\%`$ evolution on the intrinsic magnitude with redshift. Thus the SNLS experiment will be able to distinguish between the $`\mathrm{\Lambda }\mathrm{CDM}`$ model and this one at a 3 sigma level.
The SNAP/JDEM SNAPweb mission will observe about 2000 SNIa with redshifts up to 1.7. The intrinsic evolution of SNIa magnitudes is expected to be controlled at the percent level. Again, using $`\mathrm{\Lambda }\mathrm{CDM}`$ as a fiducial model and 300 SNIa from the SN factory model gives:$`\alpha =1.24\pm 0.02(\mathrm{stat}.)\pm 0.04(\mathrm{syst}.)`$, where the systematical error is coming from a $`2\%`$ remaining possible evolution of the intrinsic magnitude. The SNAP/JDEM mission will thus definitely answer the question of the compatibility of this new model with supernovae observations.
Finally, as explained in fhc2 , the constantly accelerated evolution could be affected by Pioneer like effets (as it is indeed in our neighborhood) resulting in locally inverted evolutions of space-space metric elements. This should result in a systematical drift proportional to the relative photons time of flight accross the regions with inverted cosmological regime along their path. The effect predominantly affects the higher redshift part of the path, when the matter structure (living in the inverted regime) occupied a relatively more important fraction of space, so a jerk behaviour is a natural outcome of the model. This could account for the difference between the constantly accelerated regime prediction of this model and the current data favoring a recent transition from a decelerating to an accelerating universe.
## VI Conclusion
Introducing discrete symmetries in the context of General Relativity not only allows to solve many long lasting theoretical issues such as negative energies and stability, QFT vacuum divergences and the cosmological constant but also leads to very remarkable phenomenological predictions. A constantly accelerating necessarily flat universe is a natural outcome of the model and cannot be excluded by present data The large scale structure formation and evolution need to be completely revisited in the new context where the rules of the game are significantly modified due to the interactions between conjugate density fluctuations. Last, as shown in fhc2 , the space/time exchange symmetry allows the derivation of a propagating solution and the clarification of the status of tachyonic representations. The published supernovae data are not in disagreement with this model. The one is very sensitive to systematical effects. Therefore, only the SNAP/JDEM mission should be able to distinguish without ambiguity between standard cosmology and a constantly accelerated regime as predicted by the new model in the case where no pioneer-like theoretical systematical effect is to be expected. |
warning/0507/hep-lat0507011.html | ar5iv | text | # The Low-Lying Dirac Spectrum of Staggered Quarks
## I Introduction
The low energy regime of Quantum Chromodynamics (QCD) exhibits a rich and interesting phenomenology, including the $`U_A(1)`$ axial anomaly, chiral symmetry breaking and the topological properties of the theory. A crucial step in elucidating these effects is understanding the low-lying eigenvalue spectrum of the Dirac operator. There are a number of detailed predictions of the properties of these low-lying modes, such as the existence of an Index Theorem, the Banks-Casher relation, and the distribution of the first few eigenvalues in fixed topological charge sectors. Such effects are, however, inherently non-perturbative and can only be studied fully using techniques like lattice Monte Carlo simulation.
Universality decrees that any correct discretization of QCD must be correct close enough to the continuum limit. If we want to make phenomenological predictions, however, it is crucial that the QCD lattice Dirac operator also has the correct low-lying spectrum at the lattice spacings typically simulated. By this we mean that the deformations to the eigenvalue spectrum brought about by the discretization must be small. In this case the relevant scale for comparison is the light sea quark mass, $`m_{u,d}=𝒪(5\text{ MeV})`$.
In this paper we show that these deformations are small enough for one particular class of lattice fermions that are widely used in large scale lattice simulations: improved staggered quarks. Working at similar lattice spacings in the quenched theory (i.e. pure Yang–Mills gluodynamics) we show that the low-lying QCD Dirac eigenvalue spectrum very closely reproduces the continuum features. In particular, we show that improved staggered fermions do respond correctly to the gluonic topological charge, and discuss why some confusion on this issue exists in the literature.
To reach these conclusions we test a number of predictions, quantitative and qualitative, for the low-lying modes: the Atiyah–Singer Index Theorem, that relates the zero modes to the topological charge of the gauge field; the Banks–Casher relation, linking the chiral condensate to the density of eigenvalues at the origin; the universality of the low-lying eigenvalue spectrum in the $`\epsilon `$-regime, and its connection to Random Matrix Theory. In addition, we change both the lattice spacing and volume and show that the variation of the spectrum matches theoretical expectations.
This paper expands and refines the work in Follana:2004sz ; Follana:2004wp . We include much more extensive data, both on larger volumes and finer lattices. Related work has been presented in Wong1 ; Wong2 .
The structure of the paper is as follows. In Sec. II we describe our expectations for the Dirac eigenvalue spectrum, both in the continuum and with lattice staggered fermions. We also review the existing literature. We describe our methodology in Sec. III, and discuss the features of the spectrum in Sec. IV. In Sec. V we compare our results for the low-lying eigenvalues with the universal predictions from random matrix theory. We summarise and conclude our study in Sec. VI.
## II Theoretical Background
In this section we review the theoretical expectations for the spectrum of the Dirac operator in the continuum and lattice staggered cases.
### II.1 Continuum QCD
A single, continuum fermion species is described by the massless, gauge covariant Dirac operator, which is anti-Hermitian with a purely imaginary eigenvalue spectrum:
$$\text{ / }Df_s=i\lambda _sf_s,\lambda _s.$$
(1)
Choosing orthonormalised eigenvectors, satisfying
$$f_s^{}f_t=\delta _{st},$$
(2)
the chirality of an eigenvector is defined by
$$\chi _sf_s^{}\gamma _5f_s.$$
(3)
The Dirac operator also anticommutes with $`\gamma _5`$,
$$\{\text{ / }D,\gamma _5\}=0,$$
(4)
implying that the spectrum is symmetric about zero:
$$sp(\text{ / }D)=\{\pm i\lambda _s,\lambda _s\}.$$
(5)
If $`\lambda _s0`$, then $`\gamma _5f_s`$ is also an eigenvector with eigenvalue $`i\lambda _s`$, and zero chirality, $`\chi _s=0`$. The zero modes, $`\lambda _s=0`$, can be chosen with definite chirality, $`\chi _s=\pm 1`$.
#### II.1.1 Topological charge and the Index Theorem
Smooth continuum gauge field configurations can be classified by their winding number, which is an integer given by the expression:
$$Q=\frac{1}{32\pi ^2}d^4xϵ_{\mu \nu \sigma \tau }\text{Tr}F_{\mu \nu }(x)F_{\sigma \tau }(x).$$
(6)
If we now consider the Dirac operator in the given gauge background, and we denote by $`n_\pm `$ the number of zero modes with a given chirality, their difference is related to the topological charge via the Atiyah–Singer Index Theorem Atiyah:1963 ; Atiyah:1968 :
$$Q=m\text{Tr}\frac{\gamma _5}{\text{ / }D(A)+m}=\underset{s}{}\chi _s=n_+n_{},$$
(7)
where $`m`$ is the quark mass (we keep the sign conventions of Ref. Follana:2004sz ).
#### II.1.2 The chiral condensate and the Banks–Casher relation
The bare chiral condensate in a given gauge field background, $`A`$, is given by the trace over the quark propagator:
$$\overline{\psi }\psi (m)=\frac{1}{V}\text{Tr}\frac{1}{\text{ / }D(A)+m}.$$
(8)
In the limit of $`m0`$ it acts as an order parameter for the spontaneous breaking of chiral symmetry. Expanding in eigenmodes of $`\text{ / }D`$ and using the $`\pm i\lambda _s`$ symmetry,
$$\overline{\psi }\psi (m)=\frac{1}{V}\underset{s}{}\frac{1}{i\lambda _s+m}=\frac{m}{V}\underset{s}{}\frac{1}{\lambda _s^2+m^2}.$$
(9)
Averaging over a full ensemble of gauge fields, the sum over $`s`$ can be replaced by an integral over eigenvalues with measure given by the spectral density:
$$\rho (\lambda )\frac{1}{V}\underset{s}{}\delta (\lambda \lambda _s).$$
(10)
The integral has a smooth limit as $`m0`$:
$`\mathrm{\Sigma }\overline{\psi }\psi (m=0)`$ $`=`$ $`\underset{m0}{lim}m{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\lambda {\displaystyle \frac{\rho (\lambda )}{\lambda ^2+m^2}}`$ (11)
$`=`$ $`\pi \rho (0),`$
which is the Banks–Casher relation. We expect the chiral condensate $`\mathrm{\Sigma }`$ to be non-zero when the chiral symmetry is spontaneously broken. This requires that the large volume limit is taken before $`m0`$ to ensure that the discrete sum over eigenmodes can be replaced by the smooth integral.
#### II.1.3 The $`\epsilon `$-regime and spectral universality
The $`\epsilon `$-regime of QCD occurs when the theory is regularised in a finite volume such that the typical system size $`L`$ satisfies
$`\epsilon `$-regime: $`(\mathrm{\Lambda }_{\text{QCD}})^1L(M_\pi )^1.`$ (12)
In such systems, the volume is large enough that we may derive an effective chiral field theory for pions by integrating out all other hadronic states. The typical Compton wavelength of the particles so removed is $`(\mathrm{\Lambda }_{\text{QCD}})^1`$. At the same time, the system is very small compared to the Compton wavelength of the pions, so we can further integrate out all their kinetic modes to leave us with a theory describing just static pions.
It can then be rigorously shown that the remaining finite volume partition function exactly corresponds to that of a random matrix theory (RMT) with the same chiral symmetries as QCD Damgaard:2001ep ; Akemann:2003tv . This universality Leutwyler:1992yt ; Shuryak:1993pi allows us to predict the distributions of low-lying non-zero eigenvalues of the Dirac operator using RMT Nishigaki:1998is ; Damgaard:2000ah . The relevant universality class is determined by the chiral symmetries of the finite volume partition function, and hence of QCD; for fermions in the fundamental representation of the gauge group with $`N_c3`$ the appropriate class is the chiral Unitary Ensemble (chUE) with Dyson index $`\beta =2`$.
Rather than the absolute eigenvalue spacings, RMT predicts only the relative spacings of the eigenvalues in the form of microscopic (or unfolded) spectral densities $`\rho _s(z)`$. These detail the universal distribution of one or more eigenvalues. There are separate predictions for each sector of topological charge. The absolute spectral densities include a QCD-specific scale factor based on the chiral condensate:
$$\rho (\lambda )=\mathrm{\Sigma }V\rho _s(\lambda \mathrm{\Sigma }V).$$
(13)
Comparison of the Dirac spectrum with RMT therefore provides a method of determining the chiral symmetries of QCD, of probing the role of topology, and of measuring the chiral condensate.
For the combined spectral density of all eigenvalues in a given topological charge sector RMT predicts Verbaarschot:1993pm
$$\rho _{\text{all}}(z)=\frac{z}{2}\left[J_\nu ^2(z)J_{\nu +1}(z)J_{\nu 1}(z)\right],$$
(14)
where $`J_\nu (z)`$ are Bessel functions and $`\nu =|Q|`$. As these and all other predictions depend only on $`|Q|`$, we improve the statistical accuracy of the comparison with the Dirac spectrum by combining the $`\pm |Q|`$ sub-ensembles.
There are also separate predictions for the $`k^{\text{th}}`$ eigenvalue Forrester1993 ; Damgaard:2000ah , e.g. $`k=1`$:
$$\rho _1(z)=\{\begin{array}{cc}\frac{z}{2}e^{z^2/4},\hfill & (\nu =0),\hfill \\ \frac{z}{2}I_2(z)e^{z^2/4},\hfill & (\nu =1),\hfill \\ \frac{z}{2}\left(I_2^2(z)I_3(z)I_1(z)\right)e^{z^2/4},\hfill & (\nu =2).\hfill \end{array}$$
(15)
Expressions for higher $`k`$ are given in Ref. Damgaard:2000ah .
### II.2 Lattice staggered fermions
The massless, gauge-invariant, One-link staggered Dirac operator on a $`d=4`$ dimensional Euclidean lattice with spacing $`a`$ is
$`\text{}D(x,y)`$ $`=`$ $`{\displaystyle \frac{1}{2au_0}}{\displaystyle \underset{\mu =1}{\overset{d}{}}}\eta _\mu (x)[U_\mu (x)\delta _{x+\widehat{\mu },y}H.c.]`$ (16)
$`\eta _\nu (x)`$ $`=`$ $`(1)^{_{\mu <\nu }x_\mu }`$ (17)
with $`u_0`$ an optional tadpole-improvement factor given, in our case, by the fourth root of the mean plaquette. $`\text{ / }D`$ is antihermitian and obeys a remnant of the continuum $`\gamma _5`$ anticommutation relation:
$$\{\text{ / }D,ϵ\}=0,\mathrm{with}ϵ(x)=(1)^{_{\mu =1}^dx_\mu }.$$
(18)
As in the continuum case, its spectrum is therefore purely imaginary, with eigenvalues occurring in complex conjugate pairs:
$$sp(\text{ / }D)=\{\pm i\lambda _s,\lambda _s\}.$$
(19)
In fact, if $`f_s`$ is an eigenvector with eigenvalue $`i\lambda _s`$, then $`ϵf_s`$ is also an eigenvector with eigenvalue $`i\lambda _s`$.
The corresponding action describes $`N_t=2^{d/2}=4`$ “tastes” of fermions which interact via unphysical “taste breaking” interactions that vanish in the continuum limit as $`a^2`$. In the limit $`a0`$ there is an $`SU(N_t)SU(N_t)`$ chiral symmetry, and the spectrum is therefore an exact $`N_t`$-fold copy of that described in Sec. II.1.
At finite lattice spacing the chiral symmetry group is reduced to $`U(1)U(1)`$ and we do not expect to see this picture, for example there will not be an exact Index Theorem anymore. The relevant operator for the Index Theorem must be a taste-singlet, $`\gamma _5^{ts}`$, to couple to the vacuum correctly Golterman:1984 . It is a point split operator in all $`d`$-dimensions, which can be made gauge invariant by the insertion of the appropriate $`U`$ fields. It does not anticommute with the staggered Dirac operator, $`\{\text{ / }D,\gamma _5^{ts}\}0`$. This also implies that the chirality of the eigenmodes will not be exactly $`\pm 1`$ or $`0`$.
Exact zero eigenvalues only occur accidentally (and are completely excluded at non-zero fermion mass), and all the eigenmodes will in principle contribute in Eq. (7) Smit:1987fn . Due to the simple properties
$`[\gamma _5^{ts},ϵ]=0,`$ (20)
$`ϵ^{}ϵ=I,`$ (21)
each pair of eigenvectors {$`f_s`$, $`ϵf_s`$} has the same chirality $`\chi _s`$:
$$ϵf_s|\gamma _5^{ts}|ϵf_s=f_s|ϵ^{}\gamma _5^{ts}ϵ|f_s=f_s|\gamma _5^{ts}|f_s.$$
(22)
All of the preceding considerations depend only on simple symmetry properties, and apply equally well to the improved staggered operators we have used in our work.
Early studies of the One-link staggered operator and unimproved, Wilson gauge action for coarse, thermalized lattices were done in Smit:1987fn ; Hands:1990wc ; Venkataraman:1998yj ; Hasenfratz:2003 ; Follana:2003 . There was no clear separation in the low-lying modes, either in eigenvalue or chirality. The low-lying One-link eigenvalues also did not follow the universal predictions for continuum QCD. In particular they showed no dependence on the topological charge sector Berbenni-Bitsch:1998tx ; Damgaard:1998ie ; Gockeler:1998jj ; Damgaard:1999bq ; Damgaard:2000qt .
In this study, however, we are interested in whether improvement of the gauge action and Dirac operator allow the continuum–like spectrum to be seen.
Certainly, some continuum features emerge using very smooth gauge backgrounds (obtained by either repeated cooling or discretising classical instanton solutions) Smit:1987fn ; Hands:1990wc . By contrast, here we are interested in thermalized, unsmoothed configurations at lattice spacings (and gauge actions) typical of lattice simulations, i.e. $`a0.1`$ fm. We conjecture that these failures of continuum predictions follow from the use of coarse lattices and unimproved actions, resulting in large taste changing interactions that ruin the continuum chiral symmetries. We shall show that the improved staggered operators introduced below allow all the continuum features of the spectrum to emerge.
That improvement allows this to happen at such lattice spacings is not entirely unexpected. QCD predicts that the topological susceptibility is suppressed as we reduce the sea quark mass. This is seen in simulations using improved gauge and staggered fermion actions Hasenfratz:2001wd ; Bernard:2002sa ; Bernard:2003gq ; Aubin:2004qz . With One-link staggered sea quarks no such variation is seen Kogut:1990qd ; Bitar:1991wr ; Kuramashi:1993mv ; Alles:2000cg ; Hasenfratz:2001wd . Other evidence for the beneficial effect of improved staggered operators was seen in the Schwinger model Durr:2003xs ; Durr:2004as ; Durr:2004ab ; Durr:2004rz ; Durr:2004ta , as well as in QCD Follana:2004sz ; Follana:2004wp ; Durr:2004as ; Durr:2004rz .
#### II.2.1 Improved staggered fermions
The continuum chiral symmetry is broken by the (unphysical) taste-changing interactions. At leading order these involve exchange of an ultraviolet gluon of momentum $`q\pi /a`$. Such interactions are perturbative for typical values of the lattice spacing. They can therefore be removed systematically using the method of Symanzik. This gives lattice operators with better taste symmetry and smaller scaling violations. We do this by smearing the gauge field to apply a form factor to the quark-gluon vertex and suppress the coupling between quarks and high-momentum gluons Lepage:1998vj ; Bernard:1999xx ; Orginos:1999cr . This amounts to replacing the gauge link $`U`$ in Eq. (17) by a “fat” link: a weighted sum of staples as shown graphically in Fig. 1. The choice of weight factors, $`c_i`$, determines the degree of improvement. In this paper we compare three different, improved operators known as Asqtad, Fat7xAsq and Hyp.
The first two are based on the Fat7 operator which includes staples of length up to $`7a`$. We can remove all $`𝒪(a^2)`$ discretization errors (at tree level) by adding two more terms to the Fat7 prescription, including a Naik three-link term Naik:1986bn and a five-link Lepage staple Lepage:1998vj . This operator is called Asq, whereas Asqtad is the tadpole improved version which aims to also reduce the radiative corrections to the scaling.
The improvement procedure may be iterated; for our most highly improved operator the original links are first Fat7 smeared and projected back onto the $`SU(3)`$ gauge group. These fattened links are then used to construct the Asq operator. We call this operator Fat7xAsq. We stress that this is the same operator that was used in Follana:2004sz . Regrettably, there it was misnamed Hisq, a name already used to describe a slightly different Dirac operator.
The final improved staggered Dirac operator is Hyp Knechtli:2000ku . Its construction is motivated by perfect action ideas, and involves three levels of (restricted) APE smearing with projection onto $`SU(3)`$ at each level. The restrictions are such that each fat link includes contributions only from thin links belonging to hypercubes attached to the original link.
The spectrum of all improved Dirac operators is of the form in Eq. (19). Improved operators, however, show very small taste-changing effects both in the hadronic splittings Follana:2003 ; Follana:2004 and, as we shall see, the eigenvalue spectrum.
## III Details of the simulation
In this section we provide details of the gauge backgrounds used in this study, and the methods used to determine the eigenvalues and the topological charge. We also introduce notation for the eigenvalues to be used in the following sections.
### III.1 The gauge action
We use a quenched, $`SU(3)`$ gluonic action that is both tree level Symanzik and tadpole improved Curci ; Curcierratum ; Luscher ; Luscher\_erratum ; Alford :
$$S=\beta \underset{\genfrac{}{}{0pt}{}{x}{\mu <\nu }}{}\left(\frac{5}{3}P_{\mu \nu }(x)\frac{1}{12}\frac{R_{\mu \mu \nu }(x)}{u_{0}^{}{}_{}{}^{2}}\frac{1}{12}\frac{R_{\mu \nu \nu }(x)}{u_{0}^{}{}_{}{}^{2}}\right),$$
(23)
where $`P`$, $`R`$ are $`1\times 1`$ and $`2\times 1`$ Wilson loops respectively. The tadpole improvement coefficient $`u_0`$ is defined as the fourth root of the mean plaquette.
We have generated five ensembles of around 1000 configurations, with parameters for these ensembles as in Table 1 (taken from Refs. Bonnet:2001rc ; Zhang:2001fk ). The physical spatial size of the lattice is $`aL`$, and the volume $`a^4Va^4L^3T`$. We have also used a few configurations from a Wilson gauge ensemble with lattice spacing $`a0.1`$ fm and volume $`16^3\times 32`$.
Three ensembles have been tuned to study the effect of varying the lattice spacing at fixed physical volume ($`aL1.5`$ fm). Another three vary the volume at fixed $`a=0.093`$ fm.
The two coarser lattice spacings are representative of current, large scale unquenched simulations using improved staggered quarks, with the smallest $`a`$ being indicative of future calculations. Our gauge action differs from the one used by the MILC Collaboration Davies:2003ik only in small one loop radiative corrections which do not affect the relevance of the results in this paper.
The main objective of this paper is to compare our results with the expectations from the continuum limit and with the predictions from RMT. For this purpose any systematic uncertainty in the determination of the lattice spacings (for example due to finite-volume effects) is not important. It would only become important for obtaining a precise determination of some physical quantity. Therefore in general we ignore this source of errors in our measurements, using only the statistical errors based on a jack-knife analysis of the data.
### III.2 Measuring the topological charge
We measure the topological charge $`Q`$ using a number of gluonic methods to look for consistency. The topological charge of a gluonic configuration is not uniquely defined at finite lattice spacing, and different methods will sometimes disagree. Over an ensemble, the number of configurations that disagree on the value of $`Q`$ vanishes as $`a^2`$ in the continuum limit.
We first cool the gauge fields and then apply a lattice topological charge operator. The cooling is (separately) done with two different gauge actions (the Wilson and the 5Li deForcrand:1995qq ; deForcrand:1997sq ). We measure the topological charge with the highly accurate “5 Loop improved” (5Li) operator deForcrand:1995qq ; deForcrand:1997sq . We then check that both charges are very close to integers and stable under cooling. If the topological charges (rounded to the nearest integer) defined by this procedure are the same for both actions, then we assign the configuration to a definite topological charge sector.
Does this consistency criterion bias our fixed-$`Q`$ subensemble averages? No; at $`a=0.093`$ fm ensembles fewer than $`10\%`$ of the configurations have an ambiguous topological charge. For the finest ensemble, $`a=0.077`$ fm, this number goes down to around $`2\%`$.
We have repeated parts of the analysis by ignoring the consistency requirement and using only one of the cooling actions in the definition of $`Q`$, assigning every configuration to a topological charge sector. The results are indistinguishable from those presented here, and therefore this robustness criterion neither introduces a systematic uncertainty in our results, nor changes the conclusions.
We stress that we only use cooling to determine the topological charge. All the Dirac spectrum measurements are done on the original thermalized (“hot”) configurations.
### III.3 Determining the eigenvalues
To calculate the low-lying spectrum of the Dirac operators we use an implementation of the Cullum-Willoughby Lanczos-based algorithm Cullum\_book . Due to the structure of the staggered Dirac operator the calculation of the spectrum can be done on an hermitian, positive semi-definite operator, defined only on half of the lattice (either the odd or the even sublattice). The algorithm can be extended to the calculation of the corresponding eigenvectors, which are needed to obtain the chirality.
On the ensembles used in this study we found no significant numerical difficulties with this algorithm. In practice, the absence of exact zero modes simplifies the calculation from a numerical point of view, as the linear operator we are dealing with is then positive definite.
Due to the exact $`\pm `$ symmetry, we only need to consider half of the spectrum, say $`\{i\lambda _s,\lambda _s>0\}`$ (which we call the positive sector). The other half is a mirror image which we call the negative sector. From now on we dispense with the imaginary unit $`i`$, and discuss and show on the plots the imaginary part of the positive half spectrum $`\{\lambda _s,\lambda _s>0\}`$.
For the analysis we divide the (positive) eigenmodes on each configuration into two sets based on the gluonic topological charge. Inspired by our continuum limit expectations, we expect the $`4|Q|`$ near-zero modes to divide into $`2|Q|`$ modes each side of zero. We call the arithmetic mean of the positive near-zero modes $`\mathrm{\Lambda }_0`$.
The remaining positive modes are divided into quartets. We call the means of successive sets $`\mathrm{\Lambda }_{1,2,\mathrm{}}`$. We define the intra-quartet splitting of the $`s^{\text{th}}`$ set, $`\delta \mathrm{\Lambda }_s`$, as the difference between the largest and smallest eigenvalues in the quartet. We also define the inter-quartet gap between the $`s^{\text{th}}`$ and $`(s+1)^{\text{th}}`$ quartets, $`\mathrm{\Delta }\mathrm{\Lambda }_s`$, as the difference between the smallest eigenvalue of the higher quartet and the largest eigenvalue of the lower quartet.
We use the notation $`_Q`$ to denote the expectation value over sub-ensembles with fixed gluonic topological charge $`\pm Q`$.
Unless otherwise specified, the eigenvalues are expressed in physical units of MeV, with the scale set using the quoted lattice spacing.
## IV Spectrum and Index Theorem
In this section we study the Dirac spectra, and show how the continuum results emerge for the improved operators.
### IV.1 The Index Theorem
We begin our analysis by qualitatively comparing the low-lying modes of various staggered quark operators on typical gauge backgrounds selected from those with $`Q=2`$. Working at approximately fixed lattice volume ($`aL1.50`$ fm), we compare three different gauge field actions: unimproved Wilson at $`a=0.1`$ fm; Symanzik improved, Eq. (23), at a comparable $`a=0.093`$ fm, and at the finer $`a=0.077`$ fm. In Fig. 2 we show both the value and the chirality of the first sixteen positive eigenvalues.
Near the continuum limit we expect to see first $`2|Q|=4`$ near-zero modes with their chirality renormalised slightly away from unity. The remaining modes should have chirality near zero and divide into almost degenerate quartets with intra-quartet splitting much smaller than the inter-quartet gap.
For the improved gauge action on the finer lattice the spectrum looks quite continuum-like for all Dirac operators and we see a clear Index Theorem. For the improved operators we also see a very clear quadruple degeneracy in the non-zero modes. The renormalisation of the chirality away from 1 is small for the improved Dirac operators (around $`Z=1.2`$), becoming as large as $`Z4`$ for the One-link fermions. This result is fairly typical of operator renormalisation for staggered quarks, with the improved actions showing significantly better renormalisation factors Hein ; Trottier:2002 ; Lee:2002 . We will discuss later the renormalisation of the chiral condensate and see the same effect.
On the coarser improved gauge background, we still have a good approximation to the Index Theorem using the Fat7xAsq and Hyp operators: there is a separation between near-zero modes of high chirality and the rest of the spectrum; the number of near-zero modes is $`2|Q|`$; and all of them have very similar chiralities. The rest of the spectrum have almost zero chirality and cluster into quartets. These properties can be seen to a lesser extent for the Asqtad operator, but are absent in the One-link case.
Nothing clear can be seen in the unimproved, Wilson gauge case. Improving the Dirac operator makes a difference, increasing the chirality of the low modes. Nonetheless the chiralities do not show a clear Index Theorem, and the eigenvalues do not gather into quartets.
To give an idea of the behaviour over a large ensemble, Fig. 3 shows a scatter plot of the absolute value of the chirality versus the (positive) eigenvalues. Results for different operators are plotted, using the same $`a=0.093`$ fm ensemble of gauge fields. We can see the formation of a gap between modes of small and large chirality as we improve the staggered operators, as well as the overall increase in chirality of the near-zero modes. As seen in Fig. 4 for the most improved operator on the finest ensemble, the separation gets sharper as we approach the continuum limit.
The separation is not strict even for the most improved action, with a few configurations showing intermediate values of the chirality. At least some of these cases are associated with configurations where the $`𝒪(a^2)`$ ambiguity in the topological charge is large and/or where the total number of near-zero modes is greater than the minimum prescribed by the Index Theorem, i.e. $`n_++n_{}>N_t|Q|`$.
As the near-zero modes for the more improved actions clearly separate and have well defined chirality, we may define a fermionic index $`\overline{Q}=(n_+n_{})/N_t`$ on each configuration. $`\overline{Q}`$ is then strongly correlated with $`Q`$, and we have a very good approximation to the Index Theorem. We expect the number of configurations where $`Q\overline{Q}`$ to vary as $`𝒪(a^2)`$. Although we don’t have enough statistics to test this scaling quantitatively, we can get some estimates from our results. For our $`a=0.093`$ fm ensemble, the disagreement between $`Q`$ and $`\overline{Q}`$ is around $`10\%`$, and it goes down to about $`2\%`$ on the finest $`a=0.077`$ fm ensemble. This ambiguity is therefore of the same order as the one intrinsic in the gluonic definition, as discussed in Sec. III.2. We stress here that $`\overline{Q}`$ played no part in the assignment of configurations into sectors of fixed topological charge.
We may worry that when $`Q\overline{Q}`$ the division of the eigenvalues into near-zero and non-zero modes, and the division of the latter into quartets, will be skewed. For instance, on such configurations the intra-quartet splitting will be as large as the inter-quartet gap. The examination of the histograms of $`\delta \mathrm{\Lambda }_s`$ and $`\mathrm{\Delta }\mathrm{\Lambda }_s`$ show that in all cases any slight bias on the mean is swamped by the statistical spread.
### IV.2 The Banks–Casher relation
We begin our quantitative analysis with the Banks–Casher relation (11), which links the density of non-zero eigenvalues at the origin with the chiral condensate, and thus provides a model independent determination of $`\mathrm{\Sigma }`$. For $`N_t`$ degenerate flavours, the Banks-Casher relation would be modified to
$$\mathrm{\Sigma }=\pi \rho (0)/N_t.$$
(24)
This is the relevant form for staggered fermions, with $`N_t=4`$. Note that the appropriate $`\overline{\psi }\psi `$ operator for the chiral condensate is the taste-singlet, but this is the local operator and so there are no other modifications to the Banks-Casher relation for staggered quarks.
In Fig. 5 we plot a histogram of the spectral density averaged over all configurations in the ensemble. The errors come from a jack-knife analysis of the data.
The spectrum divides clearly into two separate sections: a sharply peaked spike near $`\lambda =0`$, coming from the near-zero modes, and a broader distribution at larger eigenvalue from the remaining modes. This separation is much sharper for the improved actions. The spectral densities corresponding to the improved actions also show an excellent scaling behaviour with the lattice spacing.
For larger $`\lambda `$, a $`\lambda ^3`$ behaviour is expected to set in. Having only calculated the first 100 eigenvalues, however, we instead find the spectral density cut off at larger $`\lambda `$. It would be necessary to include many more modes to see this behaviour.
Using Eq. (24), we can extract the single taste chiral condensate by extrapolating the density of non-zero modes to zero, $`\rho (\lambda )\rho (0)`$. The chiral condensate is the order parameter for chiral symmetry breaking, and we see clear evidence from Fig. 5 that all our lattices are large enough that $`\mathrm{\Sigma }0`$.
To estimate $`\rho (0)`$ we carry out uncorrelated fits to $`\rho (\lambda )`$, excluding the near-zero mode peak from the fit range. We use a polynomial fit function
$$\rho (\lambda )=\rho _0+\lambda \rho _1+\lambda ^2\rho _2+\lambda ^3\rho _3.$$
(25)
We carry out a number of fits to each data set, constraining different combinations of $`\rho _i`$ to zero and varying the fit range. The intercept $`\rho _0`$ is then used to obtain the chiral condensate values given in the upper half of Table 2. The quoted errors cover the fit function and range uncertainties, which dominate the statistical and fitting errors.
Note that the One-link density at the origin $`\rho (0)`$ is much larger than the one for the improved case. This is clearly noticeable when inverting the operator, with the improved actions needing fewer iterations than the One-link action (also seen in Orginos:1999cr ).
The larger $`\rho (0)`$ gives a much higher value for $`\mathrm{\Sigma }`$ for the One-link action, but it is important to note that this is the bare lattice chiral condensate and does not have to agree between actions in the same way that the bare lattice quark mass does not. What should agree is the chiral condensate defined in some physical way, for example in the $`\overline{MS}`$ scheme at a fixed scale. The appropriate renormalisation factor for $`\overline{\psi }\psi `$ is the inverse of that for the quark mass and this has been calculated to $`𝒪`$$`(\alpha _s)`$ for the staggered quark actions used here and the improved gluon action (at $`𝒪`$$`(\alpha _s)`$ there is no effect from one-loop improvements to the gluon action or the presence of sea quarks). The quark mass renormalisation was quoted in Aubin:2004mas as:
$`m^{\overline{MS}}(\mu )`$ $`=`$
$`{\displaystyle \frac{am_0}{a}}`$ $`\left(1+\alpha _V(q^{})Z_m^{(2)}(a\mu ,(am)_0)+𝒪(\alpha ^2)\right),`$
$`Z_m^{(2)}(a\mu ,am_0)`$ $`=`$ $`\left(b(am_0){\displaystyle \frac{4}{3\pi }}{\displaystyle \frac{2}{\pi }}\mathrm{ln}(a\mu )\right).`$ (26)
where $`b(am_0)`$ is the finite piece of the lattice mass renormalisation and has small $`am_0`$ dependence. The $`\frac{4}{3\pi }`$ comes from the continuum mass renormalisation. With tadpole-improvement $`b(0)`$ takes the value 0.543 for the Asqtad quark action and 3.6 for the One-link action (denoted naive staggered quarks in Aubin:2004mas ). For Fat7xAsq $`b(0)`$ is 0.375. For the chiral condensate renormalisation, $`Z_S`$, we need to take the inverse of the renormalisation above. We take $`\mu `$ to be 2 GeV, as a standard reference point, and take the scale of $`\alpha _s`$ to be 2/a, as determined for the Asqtad action at comparable values of $`a`$ in Aubin:2004mas . The values of $`\alpha _s`$ in the $`V`$ scheme in the quenched approximation are taken from Mason . It is then clear that $`Z_S`$ for the Asqtad and Fat7xAsq actions is very close to 1 for all lattice spacing values, but significantly smaller than 1 for the One-link action. The renormalisation brings the value of the physical chiral condensate for the One-link action from the Banks-Casher relation down to 290 MeV, much closer to that from the improved actions. The renormalisation constant is so far from 1 in this case, however, that significant higher order corrections are to be expected. For the improved actions, however, the renormalisation is very small indeed and numbers change by less than 2% from the bare results of table 2. The estimated uncertainty from higher order perturbative corrections is 4% i.e. $`1\times \alpha _s^2`$. From table 2 we also have a statistical error of 4%. In addition there are statistical and systematic errors from the determination of the lattice spacing that are relevant here. All of these errors are small compared to the overall 10-20% error from the inability to determine the lattice spacing uniquely in the quenched approximation. This means that any physical value for the chiral condensate that we quote will have a significant error from this source alone. Our results for the Fat7xAsq action give a result for $`\left(\overline{\psi }\psi ^{1/3}\right)^{\overline{MS}}(2\mathrm{G}\mathrm{e}\mathrm{V})`$ of 260(10)(10)(30) MeV, where the first error is statistical, the second perturbative and the third a quenching error. This is in agreement with other results in the literature in the quenched approximation (for example, Giusti:1998 ). An improvement on this value requires unquenched configurations, but our result shows that an accurate calculation is certainly possible with improved staggered quarks.
### IV.3 Scaling of the near-zero modes
The near-zero eigenvalues for a given Dirac operator vary both with the lattice spacing and volume.
In Fig. 6 we show the variation of the mean near-zero modes $`\mathrm{\Lambda }_0`$ with lattice spacing at fixed physical volume ($`aL1.5`$ fm). The panels show the results for sectors of topological charge $`|Q|=1`$ and $`2`$.
For a sufficiently small lattice spacing, we expect the eigenvalues to go to zero as $`(a^2)`$ plus $`𝒪(a^4)`$ corrections. The improved staggered results show this trend across the full range of lattice spacings studied, with small $`𝒪(a^4)`$ corrections, slightly positive for Fat7xAsq but negative for Asqtad. By contrast, the One-link near-zero modes show very strong deviations from the leading order scaling behaviour.
In almost all quantities that we measure the Hyp operator gives results that agree with Fat7xAsq within one standard error. Unless there is a significant difference we do not discuss the Hyp data and for clarity they are not shown on the figures, but we have checked that all statements regarding Fat7xAsq apply equally to Hyp.
The Fat7xAsq near-zero modes are numerically smaller than those of the One-link action for $`a0.13`$ fm. The Asqtad near-zero modes are smaller than the One-link ones for $`a0.08`$ fm. Even when the modes are similarly sized, however, we have seen earlier that there is a qualitative difference in the chirality as we improve the action. This is a clear signal that the operator improvement is working.
The variation of the near-zero eigenvalues with volume at fixed lattice spacing is shown in Fig. 7. $`\mathrm{\Lambda }_0`$ is expected to vary as the inverse square root of the volume. This can be justified in a semiclassical picture, where the topological charge density is dominated by a dilute gas of instantons. The finite volume corrections to the zero modes come from the lack of large instantons, excluded from the lattice by their cores sizes $`raL`$. Such corrections should vary as $`1/\rho ^2`$, or $`1/\sqrt{a^4V}`$ Smith:1998wt .
We see from Fig. 7 that the data at fixed lattice spacing follows this approximate volume scaling quite closely. We might ask whether $`\mathrm{\Lambda }_0`$ really goes to zero in the large volume limit at finite lattice spacing. Fits to our data allowing a non-zero intercept are not stable, but a large reduction in statistical errors, as well as simulations in larger volumes, would be needed to answer this question definitively.
### IV.4 Scaling of the low-lying non-zero modes
We turn now to the low-lying non-zero modes. From the requirement of a physical spectral density in Eq. (10) and the Banks-Casher relation, Eq. (11), it is clear that the size of these is governed, up to a constant, by the product of the bare chiral condensate and the volume $`\mathrm{\Sigma }V`$. Specifically, the number of modes in a small interval $`(0,\delta )`$ should scale as $`\delta \rho (0)V\delta \mathrm{\Sigma }V`$, and therefore we expect each mode to scale as $`(\mathrm{\Sigma }V)^1`$. The variation of these modes with changing lattice spacing and volume is therefore fixed by the response of the bare chiral condensate.
In Fig. 8 we show the variation of the first two non-zero mode quartets $`\mathrm{\Lambda }_{1,2}`$ with lattice spacing at fixed physical volume ($`aL1.5`$ fm) for several $`Q`$. For all the improved actions the changes are very small, indicating only small discretization errors, as for the chiral condensate. For the near-zero modes, the deviations from the leading order scaling behaviour seem to be of opposite sign for Asqtad and Fat7xAsq. The One-link data show much larger deviations, of the same sign as for the Asqtad data.
The different actions appear to lack a common continuum limit. This is not surprising: as discussed above, the plot is inversely related to the density and therefore the bare chiral condensate. As explained earlier, this receives a large renormalisation for the One-link action compared to the improved actions that means that the bare results are not comparable.
From the arguments above, we expect $`\mathrm{\Lambda }_s`$ to vary as the inverse volume at fixed lattice spacing. We find this to be so for both One-link and improved actions. Fig. 9 shows the results for $`s=1,2`$.
Given the different volume scaling of the two sets of eigenmodes, we cannot really talk about a gap between near-zero and non-zero modes: we can always choose a volume large enough that the “non-zero” modes (suppressed as $`1/V`$) are numerically smaller than the “near-zero” modes (which fall only as $`1/\sqrt{V}`$). The modes are still qualitatively different however, and can be separated by their chirality.
Taste breaking interactions disappear in the continuum limit and the eigenmodes should recover a 4-fold degeneracy. We plot the intra-quartet splitting of the first quartet, $`\delta \mathrm{\Lambda }_1`$ for different sectors of topological charge in Fig. 10 (higher quartets have smaller splittings). The lattice volume is fixed at $`aL1.5`$ fm. The splitting falls with lattice spacing, as expected, and in a way consistent with the $`𝒪(a^2)`$ dependence of the taste breaking interactions. The splitting depends only weakly on the topological charge.
This is the only quantity we have studied where there is a difference of more than one standard error between Hyp and Fat7xAsq results; the Fat7xAsq results are marginally closer to our continuum expectations. (This was also true in Fig. 6, but the difference was not statistically significant).
There are strong similarities between this plot and Fig. 6: the improved fermion actions all show only very weak deviations from the naive scaling behaviour; the $`𝒪(a^4)`$ corrections are positive for Fat7xAsq and negative for Asqtad; the One-link results however show very large deviations from the expected behaviour. Numerically, the Fat7xAsq splittings are smaller than the One-link for $`a0.12`$ fm; the Asqtad numbers are less for $`a0.08`$ fm.
In Fig. 11 we show the variation of the splitting with volume at fixed lattice spacing. Once more the splitting is independent of the topological charge. There is a striking difference in the volume dependence: the One-link splittings vary markedly whilst the improved action results are nearly insensitive.
To satisfy ourselves that there is a clear restoration of the taste symmetry, we need the intra-quartet splitting $`\delta \mathrm{\Lambda }_s`$ to be much less than the gap between neighbouring quartets $`\mathrm{\Delta }\mathrm{\Lambda }_s`$. In Fig. 12 we show the ratio of these as a function of lattice spacing for $`s=1`$ (where it is largest).
Again, the ratios show little dependence on the topological charge. In all cases the improved action results are better than the One-link. This is increasingly true at small lattice spacings, as the improved actions show a ratio falling clearly to zero as $`a^2`$.
We can understand these results quite easily. As the quartets become degenerate, $`\mathrm{\Delta }\mathrm{\Lambda }_s`$ will tend towards the difference between the quartet means. Fig. 8 shows this is controlled by the chiral condensate and has little lattice spacing dependence for the improved actions. The intra-quartet splitting falls as $`a^2`$ and so, therefore, will the ratio.
## V RMT Predictions
As discussed previously, the chiral symmetries of staggered quarks are more complicated than for continuum QCD. At finite lattice spacing the $`N_t^2=16`$ pions split into 5 multiplets containing $`(1,4,6,4,1)`$ states, and only one of these states becomes massless in the chiral limit. This is known as the “Goldstone pion”, with mass $`M_G`$. The fifteen remaining states have masses $`M_{\text{NG}}`$ arising from the taste breaking interactions. The (chirally extrapolated) masses are $`𝒪(a^2)`$ and therefore zero only in the continuum limit.
There are potentially two universal regions for different parameters of the system:
$`\epsilon `$-regime: $`(\mathrm{\Lambda }_{\text{QCD}})^1L(M_{\text{NG}})^1,`$
$`\epsilon ^{}`$-regime: $`(M_{\text{NG}})^1L(M_\text{G})^1.`$
The first corresponds to Eq. (12), with the same chiral symmetries that we expect in the continuum. The low-lying non-zero modes have a near $`N_t`$-fold degeneracy, and follow the predictions of Eqs. (14,15). These are exactly the same distributions as for continuum QCD (and chiral lattice fermions Edwards:1999ra ; Edwards:1999zm ; Damgaard:1999tk ; Hasenfratz:2002rp ; Bietenholz:2003mi ; Giusti:2003gf ).
At finite lattice spacing there is a second, $`\epsilon ^{}`$-regime, where the finite volume partition function describes only the static mode of the single Goldstone pion. There is not even approximate restoration of the continuum symmetries, and the associated RMT has only $`U(1)U(1)`$ chiral symmetry. The universal predictions for the eigenvalues are therefore strikingly different from those for continuum QCD. In particular, the predictions are the same for all sectors of topological charge Damgaard:PrivComm . Presumably it is this regime that was studied in Berbenni-Bitsch:1998tx ; Damgaard:1998ie ; Gockeler:1998jj ; Damgaard:1999bq ; Damgaard:2000qt . With the coarse lattices and unimproved gauge ensembles used in these studies, they observed no sensitivity of the eigenvalue spectra to $`Q`$, leading to the incorrect folklore that staggered quarks are “blind to the topology”.
To see the continuum chiral symmetries, there must be a sufficiently large mass gap between the heaviest pion and the lightest of the other hadrons. This requires the use of improved gauge and fermion actions and a sufficiently small lattice spacing. In particular if we increase the size of the lattice $`L`$, the lattice spacing must be reduced accordingly in order to satisfy $`L(M_{\text{NG}})^1`$. Otherwise we would start to see the effects of the mass of the non-Goldstone pions. At coarse lattice spacings, with unimproved fermions, $`M_{\text{NG}}\mathrm{\Lambda }_{\text{QCD}}`$ and there will be no $`\epsilon `$-regime.
If the lattice spacing is keep fixed and we increase the volume, we expect there will be a non-universal crossover until the lattice reaches the $`\epsilon ^{}`$-regime when universal behaviour will resume. Of course, it would be interesting to simulate a range of lattice sizes and witness the crossover between the two universal regions. The range of lattice spacings is, however, well beyond the resources of this study. We instead concentrate on the continuum-like region.
### V.1 Comparison
We begin by comparing the ratios of the mean non-zero eigenvalues $`\mathrm{\Lambda }_s_Q/\mathrm{\Lambda }_t_Q`$ (which we denote “s/t”) to the ratios from RMT. Factors of the chiral condensate scale away, so this provides a parameter–free test of the $`\epsilon `$-regime predictions.
We show the eigenvalue ratios in Figs. 13-15. The ratios differ markedly between different topological charge sectors, showing clearly that the staggered fermions are in no way “topology blind”. Also shown are the ratios predicted by RMT. Fig. 13 shows that all the staggered actions agree quite well with the expected ratios at $`a=0.093`$ fm. We therefore have good evidence for the existence of a continuum $`\epsilon `$-regime. It also implies that the continuum chiral symmetries are present, with the correct pattern of spontaneous breaking.
There are some deviations from the universal predictions. They are not discretization effects: Fig. 14 shows the Fat7xAsq ratios for a range of different lattice spacings at fixed volume $`aL1.5`$ fm and discrepancies exist even on the finest lattice. Increasing the system size from $`1.5`$ to $`1.8`$ fm in Fig. 15 does reduce the size of the disagreement. Indeed, on the larger volume the ratios agree with RMT at least as well as for chiral fermions on comparable lattices Giusti:2003gf .
Fig. 13 appears to suggest that the One-link Dirac operator fits the universal predictions as well as the improved operators. This appears so even on the coarsest lattices despite the large pion splitting and lack of obvious spectral quartets in the One-link spectrum.
In this respect the “ratios plots” can be misleading. Firstly they only compare the mean eigenvalues with RMT and provide no information as to whether the shape of the individual distributions match. Also, when comparing with RMT we must rescale all the eigenvalues using the chiral condensate. Agreement with the ratio plot does not imply that eigenvalues in different topological charge sectors are scaled by the same factor.
To address the first point we separately fit the individual spectral densities to the predictions from RMT. These one parameter fits based on Eq. (13) yield a prediction for the chiral condensate. To avoid the ambiguities of histogram bin-size we use the cumulative eigenvalue distributions,
$$p(\lambda )=_0^\lambda 𝑑\lambda ^{}\rho (\lambda ^{}).$$
(27)
We show examples of the fits in Figs. 16 and 17.
For all operators the fits are worst for low $`s`$ and for small $`|Q|`$. We can qualitatively understand this in terms of eigenvalue repulsion. The taste changing interactions have split the near-zero modes, which has a knock-on effect on the rest of the spectrum as the eigenvalues avoid level crossings. The lower quartets are closest to the near-zero modes and are thus most strongly affected. As the quartets get numerically larger as we increase the topological charge, so the repulsion should be less.
It is clear that the improved Dirac operators match the predictions of RMT very closely over the full range of lattice spacings. By contrast, even at $`a=0.077`$ fm the One-link Dirac operator still shows significant discrepancies.
Separately fitting the $`(Q=0,k=\mathrm{1..4})`$ and $`(|Q|=\mathrm{0..3},k=1)`$ eigenvalue distributions to the RMT predictions gives 7 values for the bare chiral condensate for each Dirac operator on each ensemble. We find these numbers to be broadly consistent. The spread in the numbers is greatest for the One-link action and on the smallest volumes. This is not surprising, when such eigenvalue distributions showed the greatest deviation from the RMT predictions
The lower half of Table 2 gives estimates for the chiral condensate obtained from the median of the seven individual numbers. The quoted errors are half the range, which forms the dominant uncertainty.
We can compare these with the model-independent estimates of $`\mathrm{\Sigma }`$ obtained using the Banks–Casher relation. There is very good agreement between the numbers. The dominant source of error is different in each case: for the Banks-Casher relation it comes from the choice of fit range and function used to extrapolate the spectral density to zero eigenvalue. For the RMT comparisons, it comes from the variation in fit parameters for different eigenvalues, and is therefore related to exactly how well the effective partition function describes the lattice theory. The good agreement of the numbers suggests both systematic biases are under control. See section IV.2 for a discussion of how the bare chiral condensate must be renormalised to extract a physical result.
## VI Conclusions and Outlook
It is widely believed that any “problems” with a given lattice Dirac operator will show themselves most clearly in the topologically sensitive, low-lying eigenmodes. In this study we have studied these modes for the staggered lattice Dirac operator, using a range of improved and unimproved versions. The improved operators pass all the tests. In particular, and contrary to previous accepted wisdom, such fermions do respond exactly as expected to the gluonic topological charge, and in a way identical to continuum QCD and other lattice formulations of the Dirac operator.
We have seen that for improved operators the eigenvalue spectrum divides cleanly into near-zero and non-zero modes. The near-zero modes are characterised by a uniformly high chirality, with their relative number fixed by the Atiyah-Singer Index Theorem. Indeed they could be used to define the topological charge $`Q`$ of a configuration. Their eigenvalues scale with lattice spacing and volume as expected.
The non-zero modes, by contrast, have chirality that is near zero. They divide into near degenerate quartets, with the splitting reducing to zero as $`a^2`$ in common with the taste breaking interactions. The quartet means are controlled by the chiral condensate and scale with lattice spacing and volume accordingly. In addition, the low-lying non-zero modes follow closely the universal distributions predicted by random matrix theory.
Since the near-zero and non-zero modes scale differently with the volume, it is clear that the “gap” between the highest near-zero and the lowest non-zero modes is a volume dependent quantity. At finite lattice spacing we can always go to a volume large enough that the two sets of eigenvalues have comparable magnitude. This is not really an issue because the gap is not a physically observable quantity: what matters is the size of the near-zero modes compared to the light quark masses and whether there is a well–defined Index Theorem. As we have seen the two sets of modes are distinguishable by their chirality, even in the large volume limit.
A quantitative study of axial anomaly physics will require the near-zero modes to be smaller than the light quark masses, i.e. $`\lambda <𝒪(5\mathrm{M}\mathrm{e}\mathrm{V})`$, or, at least, smaller than the sea and valence quark masses used in the simulation. Then the eigenvalue spectrum is cut off at the low end by the quark mass and the fact that the near-zero modes are not actually zero is irrelevant.
To understand what this means for our quenched ensembles, it is necessary to define the size of the typical near-zero mode. The mean zero mode $`\lambda _0_Q`$ increases as we increase $`Q`$, due to eigenvalue repulsion. For our $`\beta =4.8`$, $`L=20`$ lattice we find that the mean zero-mode is fitted quite well by $`\lambda _0_Q(3+Q)`$ MeV for $`Q=\mathrm{1..3}`$. If we define the typical topological charge $`\overline{Q}`$ to be the RMS value from the topological susceptibility, $`\overline{Q}=\sqrt{\chi V}=2.6`$, then the typical near-zero mode is $`\overline{\lambda }_0=(3+2.6)6`$ MeV. This suggests that we are already able to study the quenched axial anomaly.
We can attempt to extrapolate our quenched results to ask whether this condition is true on existing configurations from the MILC collaboration Aubin:2004wf . This requires various unsupported, if reasonable, assumptions, not least that sea quark effects do not strongly affect the size of individual near-zero modes and that $`\lambda _0_Q`$ continues to be linear in $`Q`$.
Typical dynamical simulation volumes are larger than we have used in this quenched study, to avoid finite volume effects in hadron spectroscopy. As well as an overall $`\sqrt{V}`$ suppression in the size of the near-zero modes, this will also affect the typical value $`\overline{Q}`$ and therefore the amount of eigenvalue repulsion. There are two competing effects: the sea quark induced screening of topological charge reduces $`\chi `$, whilst the larger volume increases the width of the histogram of $`Q`$.
Consider the $`a=0.09`$ fm MILC ensemble with sea quark masses $`m_s/m_{u,d}=5`$ and volume $`V=28^3\times 96`$. Using the data in Bernard:2003gq , we obtain $`\overline{Q}=7.3`$. Assuming the quenched relation above holds, we have $`\overline{\lambda }_0=(3+7.3)/\sqrt{28^4/20^4}6`$ MeV (to be conservative we assume that the volume scaling of the near-zero modes goes with the smallest length of the lattice, therefore $`28^4`$ instead of $`28^3\times 96`$). This is very encouraging. The lightest sea quark mass presently used in these dynamical simulations is $`m_0=16`$ MeV, which is several times larger than this projected near-zero mode.
Even given the caveats surrounding this rough estimate, it does seem likely that a systematic study of the QCD axial anomaly is possible on these configurations and various aspects of this work are now in progress.
An on-going debate surrounds the theoretical justification of the methods use to simulate $`N_f<N_t`$ flavours of sea quarks Bunk:2004br ; Bunk:2004kf ; Hart:2004sz . Recent work suggests there is no ambiguity Adams:2004mf ; Adams:2004wp ; Maresca:2004me . Our results (and related work elsewhere Durr:2004as ; Durr:2004rz ) indicates that with improved staggered quarks the spectrum matches very closely the continuum expectations with clear restoration of taste symmetry. It therefore paves the way for the full analysis of the staggered taste basis that is necessary to finally resolve this issue.
To clarify the role of topology and the issue of the $`\epsilon `$-regime for staggered quarks is important because improved staggered fermions are currently the only lattice formulation that offers the prospect of a comprehensive study of the QCD axial anomaly free from the systematic effects of lattice spacing and unphysically large sea quark mass.
###### Acknowledgements.
This research is part of the EU integrated infrastructure initiative hadron physics project under contract number RII3-CT-2004-506078. We thank: Ph. de Forcrand for his topological charge measurement code; A. Hasenfratz for help in implementing the Hyp operator; G.P. Lepage, P. Damgaard and G. Akemann for useful discussions. E.F., C.T.H.D. and Q.M. are supported by PPARC and A.H. by the U.K. Royal Society. The eigenvalue calculations were carried out on computer clusters at Scotgrid and the Dallas Southern Methodist University. We thank David Martin and Kent Hornbostel for assistance. |
warning/0507/astro-ph0507508.html | ar5iv | text | # The complex velocity distribution of galaxies in Abell 1689: implications for mass modelling
## 1 Introduction
The increasing amount of data on Abell 1689, a cluster of galaxies at z= 0.183, has recently motivated several detailed analyses of its dynamical status and mass distribution. As the largest known gravitational lensing object it has been studied in detail by Broadhurst et al. (2005a, 2005b) who found that the inferred mass distribution is much steeper compared to what is expected for dark matter haloes forming in currently available cosmological $`N`$-body simulations. In particular they find the concentration parameter of the best-fitting NFW (Navarro, Frenk & White 1997) profile of the cluster to be between $`c=8`$ and $`c=14`$ depending on projected radius to which the mass distribution was studied. The smaller value was obtained from the strong lensing results in the inner part of the cluster, while the larger value was found when the study was extended to a scale of 2 Mpc using the results from weak lensing. The larger value seems rather high for the estimated mass of $`2\times 10^{15}M_{\mathrm{}}`$ compared to the expected value of $`c=5`$ for haloes of this mass as found in cosmological $`N`$-body simulations (e.g. Bullock et al. 2001). The discrepancy has led Oguri et al. (2005) to claim that the mass distribution in A1689 may be in conflict with the standard cold dark matter (CDM) scenarios for structure formation. They have shown that the disagreement can be partially reduced if the uncertainties in the parameter estimation due to the possible triaxiality of the halo are properly taken into account.
The X-ray data for Abell 1689 obtained with the XMM-Newton telescope have been analyzed by Andersson & Madejski (2004). The X-ray gas surface brightness distribution appears rather regular and smooth. However, a closer inspection reveals that the gas temperature profile is highly asymmetric and the gas mass fraction is lower than usual, which may point towards a perturbed structure. Moreover, the total mass inferred from the X-ray analysis gives a value twice as small as that found from gravitational lensing.
Kinematical analysis of the cluster galaxies shows even larger discrepancies. The redshift survey performed by Teague, Carter & Gray (1990) led to the identification of a few structures along the line of sight and an estimate of 2355 km s<sup>-1</sup> for the velocity dispersion of the cluster members. Struble & Rood (1999) estimate the cluster dispersion to be 1989 km s<sup>-1</sup> using the same data. Girardi et al. (1997) applied the wavelet analysis to the same data and detected even more substructure. Although their estimate for the main cluster velocity dispersion was 1429 km s<sup>-1</sup>, they calculated the mass by adding the masses of two main substructures with low velocity dispersions of the order of 300-400 km s<sup>-1</sup> which led to the value of about $`2\times 10^{14}M_{\mathrm{}}`$, an order of magnitude lower than the mass obtained from the lensing studies.
In this Letter we reanalyze the velocity distribution of galaxies in the field of A1689 using the larger sample now available. We confirm that the cluster indeed has a complex structure in velocity space, strongly indicating the presence of dynamically independent structures along the line of sight. By imposing different cut-offs in velocity we show how the cluster mass estimate can change by a large factor, which illustrates the difficulty in inferring it from the kinematical data. The complicated mass distribution around the cluster may also affect mass estimates done with other methods. We refer to cosmological $`N`$-body simulations in order to demonstrate that distant haloes positioned along the line of observations can indeed produce line-of-sight velocity distributions similar to the one in A1689. Therefore any estimate of concentration for an object in such environment may be biased by an error not associated with the method of mass determination but due to the presence of foreground and background structures.
## 2 Velocity distribution of galaxies in A1689
We have searched the NED database for galaxies with redshifts $`z=0.1832\pm 0.05`$ and located at distances smaller than 2 Mpc from the cluster centre assumed to be at RA=$`13^\mathrm{h}11^\mathrm{m}30.3^\mathrm{s}`$, Dec=$`01^{}20^{}53^{\prime \prime }`$ (J2000). It corresponds to the position of the elliptical galaxy closest to the centre of the main structure detected by Girardi et al. (1997). It is also within 100 kpc from the centre of X-ray gas surface brightness distribution. The redshift data for galaxies thus chosen come mainly from surveys by Teague et al. (1990), Balogh et al. (2002) and Duc et al. (2002).
The line-of-sight velocities of 192 galaxies in the reference frame of the cluster as a function of distance from the cluster centre are shown in the upper left panel of Fig. 1. The colours code the probable membership of the galaxies in different groups separated in velocity space. The division has been made by a simple cut-off in constant $`v`$. Separating first galaxies with $`|v|>6000`$ km s<sup>-1</sup> we get a group of 15 galaxies marked by red dots. The other two groups with 3000 km s$`{}_{}{}^{1}<|v|<6000`$ km s<sup>-1</sup> are marked with green and blue respectively for the positive (35 galaxies) and negative velocities (12 galaxies). The remaining 130 galaxies with $`|v|<3000`$ km s<sup>-1</sup> are marked with black dots and correspond most probably to the main body of the cluster. The same colour coding applies to the velocity distribution histogram (number of galaxies per velocity bin of size 1000 km s<sup>-1</sup>) shown in the upper right panel of the Figure. The lower left panel of the Figure shows the position of the galaxies belonging to each group on the plane of the sky; the positions overlap indicating that the groups lie along the line of sight.
The histogram shown in the upper right panel of Fig. 1 is similar to that in Fig. 1 of Girardi et al. (1997), but the identification of structure is somewhat different. In particular, the green, red and blue peaks in our histogram are the same as those at $`cz=60`$, $`64`$ and $`52\times 10^3`$ km s<sup>-1</sup> respectively in their Figure, but we do not see the structures with velocities close to the cluster mean, which they identified as S2 and S3, as separate. Although some fluctuations in the number of galaxies can be seen in this region when we plot the histogram with better resolution, we do not think they are significant. The two upper panels of Fig. 1 show qualitatively that the cluster has a complicated structure in velocity space, at variance with what is expected for relaxed, isolated objects.
The lower right panel of Fig. 1 plots the velocity dispersion profiles calculated using different galaxy samples. The red profile was obtained from the total sample of 192 galaxies, for the blue one the galaxies with $`|v|>6000`$ km s<sup>-1</sup> with respect to the cluster mean were removed, and the black one is for galaxies with $`|v|<3000`$ km s<sup>-1</sup>. The data points were calculated with 30 galaxies per bin and assigned standard sampling errors (see Łokas & Mamon 2003). The profiles can be used as a quantitative measure of the mass of the structure. Assuming that the galaxies trace the overall NFW mass distribution in the cluster and have isotropic orbits we can estimate the parameters of the NFW profile, the virial mass $`M_v`$ and concentration $`c`$ by fitting the velocity dispersion data to the solutions of the Jeans equation
$$\sigma _{\mathrm{los}}^2(R)=\frac{2}{I(R)}_R^{\mathrm{}}\frac{\nu \sigma _r^2(r)r}{\sqrt{r^2R^2}}dr,$$
(1)
where $`\nu (r)`$ and $`I(R)`$ are the 3D and the surface distribution of the tracer as a function of a true ($`r`$) and projected ($`R`$) distance from the object centre respectively and $`\sigma _r`$ is the radial velocity dispersion related to the mass distribution in the object (see Łokas & Mamon 2001, 2003).
The best-fitting $`M_v`$ and $`c`$ values we obtain from the velocity dispersion profiles for the 3 samples we have considered are given in Table LABEL:parameters. The range of mass values illustrates well how much the estimated parameters depend on the sample of galaxies chosen. We find that for the whole sample, as well as for the intermediate sample the resulting masses are significantly larger than expected. Only the most restrictive sample gives a more reasonable value of $`2.6_1^{+2}\times 10^{15}M_{\mathrm{}}`$ (at 68 per cent confidence level), much more in agreement with the value deduced from recent studies based on lensing (Broadhurst et al. 2005a,b). Although this sample may still contain unbound galaxies which bias the result towards higher masses, it is clear that any further division of this sample into two parts of comparable size, as was done by Girardi et al (1997), would result in a mass estimate at least a factor of few lower, strongly at variance with the value estimated from lensing.
Although our best-fitting concentration for this sample is much higher than expected, as in the case of studies based on lensing, the data do not allow us to really constrain the concentration, i.e. all values in the range $`5<c<100`$ are consistent with the data at $`1\sigma `$ level. We emphasize, however, that our cut-offs in velocity were rather arbitrary and although they followed the gaps in the $`v(R)`$ diagram it would be difficult to justify them in a quantitative way. In particular, none of the galaxies in the $`v(R)`$ diagram in Fig. 1 would be removed by the application of standard methods for the rejection of outliers. This suggests that in agreement with visual impression from the $`v(R)`$ diagram, the galaxies with discrepant velocities are not just interlopers but belong to some neighbouring structures.
Our analysis illustrates the difficulties encountered when the standard Jeans approach is uncritically applied to clusters before considering all the possible indications on their dynamical status and/or their environment. As it happens for A1689, it could be that discrepant mass estimates are obtained depending on the sample selection criteria. We note that our sample is a compilation of a few surveys with substantial fraction of spiral galaxies (41 percent among those with known morphological type). While the sample of Teague et al. (1990) comes from a standard magnitude-limited survey, the selection criteria of those of Balogh et al. (2002) and Duc et al. (2002) were aimed at star-forming galaxies which may bias the sample towards outer regions with more substructure. The analysis of the cluster dynamics could be significantly improved if a survey of many galaxy redshifts complete up to a given limiting magnitude was available. This would allow for a proper comparison with other well studied clusters and a more accurate description of its velocity distribution.
## 3 Comparison with $`N`$-body simulations
In this section we make use of the $`N`$-body simulations to study the origin of complex velocity distributions like the one in A1689. For this work we used the results of a cosmological dark matter simulation described by Wojtak et al. (2005). The simulation was performed within a box of size 150 $`h^1`$ Mpc assuming the concordance cosmological model ($`\mathrm{\Lambda }`$CDM) with parameters $`\mathrm{\Omega }_M=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, $`h=0.7`$ and $`\sigma _8=0.9`$. The final output of the simulation contained a few tens of massive haloes found with standard FOF procedures. To mimic observations we place an imaginary observer at a given distance from the halo and project the particle velocities along the line of sight and their positions on the surface of the sky. Among the most massive haloes we chose one with a similar line-of-sight velocity distribution as is seen in A1689. The virial mass of the halo is $`5.4\times 10^{14}M_{\mathrm{}}`$ and it has about $`10^4`$ particles inside the virial radius which allows us to reliably measure the density profile. The properties of the halo are summarized in Table LABEL:par.
Out of all particles seen by our observer in the direction of the halo inside the projected radius $`R=r_v`$ and with velocities $`|v|<3000`$ km s<sup>-1</sup> with respect to the mean velocity of the halo we randomly select 200 particles to have a similar statistics as for A1689. (The velocity range corresponds to about $`4\sigma _{\mathrm{los}}`$ for a halo of this mass.) The summary of the observed properties of the halo is presented in Fig. 2 with the panels analogous to those in Fig. 1 for A1689. The $`v(R)`$ diagram in the upper left panel is highly irregular with particle velocities spread out over the whole velocity range and making it difficult to decide which of them should be treated as true members of the halo. A distinct structure is seen at about $`v=2500`$ km s<sup>-1</sup> which is even better visible in the upper right panel showing the line-of-sight velocity histogram. As for A1689 we separate the particles using the velocity criterion; those with $`|v|<1500`$ km s<sup>-1</sup> are marked with black dots in the plots while those with $`|v|>1500`$ km s<sup>-1</sup> with green or blue dots depending on the sign of the velocity. The lower left panel shows the positions of the particles belonging to different velocity bins on the surface of the sky.
The lower right panel of Fig. 1 plots the velocity dispersion profiles calculated from the data from all particles (blue) and only from those with $`|v|<1500`$ km s<sup>-1</sup> (black). As for A1689 we did a similar exercise of fitting these data with the solutions of the Jeans equation (1) assuming isotropic orbits and adjusting the virial mass and concentration. The results for different samples are listed in Table LABEL:par. As we can see, the estimated parameters differ dramatically depending on the sample. For the sample with $`|v|<1500`$ km s<sup>-1</sup> the virial mass is $`M_v=4.8_{1.6}^{+1.4}\times 10^{14}M_{\mathrm{}}`$ (at 68 per cent confidence level) which agrees well with the mass $`M_v=5.4\times 10^{14}M_{\mathrm{}}`$ known from the 3D information about the halo. The concentration proves more difficult to estimate with such a small sample and simple modelling since we get $`c=2.1_1^{+5}`$ (at 68 per cent c.l.) which does not include the true value $`c=9.2`$.
What is the reason behind the complicated velocity structure of the halo? The actual spatial distribution of the particles in a part of the observed region is shown in projection in Fig. 3. The centre of our halo is located at $`x=0`$, $`z=0`$. The line of sight of the observer is along the $`z`$ axis of the plots. The second, smaller halo of mass $`M_v=8\times 10^{13}M_{\mathrm{}}`$ contributing to the $`v(R)`$ diagram shown in Fig. 2 is located at the distance of 45 Mpc from the main halo. As before, we marked the particles with line-of-sight velocities $`v<1500`$ km s<sup>-1</sup> (with respect to the mean velocity of the main halo) with black dots and those with larger velocities with green ones. Left panel shows all particles in this region of the simulation box, while the right one about 200 chosen randomly to create the mock data (not all 200 particles are shown because the region in the direction of negative $`z`$ is not plotted).
The Figure demonstrates that in spite of their proximity in velocity space ($`2300`$ km s<sup>-1</sup> is of the order of $`3\sigma _{\mathrm{los}}`$ of the big halo) the two haloes are in fact very distant. The distance of the smaller halo, 45 Mpc, corresponds to about 20 virial radii of the big halo therefore the haloes do not affect each other dynamically and are not bound to each other, but still their projected velocity distributions are entangled (note that there are green particles in the centre of the bigger halo and black particles in the centre of the smaller halo). The reason for this is the rather low value of the Hubble velocity in comparison with velocity dispersion of bound structures, e.g. velocity dispersions of massive haloes or galaxy clusters are comparable to the Hubble flow at distances as large as 8 virial radii from their centres (see Fig. 3 of Wojtak et al. 2005). The smaller halo in our example is receding from the bigger one with a velocity of about $`2300`$ km s<sup>-1</sup> mainly due to the Hubble flow which at distance of 45 Mpc is of the order of $`3000`$ km s<sup>-1</sup>.
To further illustrate the point, we provide an example of a well-behaved $`v(R)`$ diagram. Fig. 4 shows again the $`v(R)`$ diagram of our simulated halo from Fig. 2 in the left panel, while in the right panel we present an analogous plot of line-of-sight velocities as a function of projected distance for 200 particles chosen from the same halo, but observed in a different direction. In this case the halo has no massive neighbours along the line of sight and the single particles with discrepant velocities can be easily dealt with using standard procedures for interloper removal.
Our purposely chosen example illustrates well the difficulties in interpreting the measured line-of-sight velocities of galaxy clusters in the case of presence of neighbouring structures. Whereas velocity differences amounting up to $`2300`$ km s<sup>-1</sup>, as in our simulated haloes, could easily be interpreted as orbital velocities of dynamically bound objects, they could also correspond to structures seen in projection but otherwise separated by distances much larger than their virial radii and, therefore, totally unrelated.
To make a connection with the studies based on lensing we note that the surface density distribution (the main lensing observable) measured along the line of sight in our simulations is increased by about 25 percent everywhere along the projected radius of the bigger halo due to the presence of the smaller halo. The significance of this effect will of course depend on the exact properties of the haloes and probability of their alignment, which will be studied elsewhere.
## 4 Conclusions
We have shown that cosmological structures quite distant from each other, when aligned with the direction of observation, can produce projected velocity distributions which are quite difficult to interpret. In particular, such extended distributions can lead to very different velocity dispersion, and therefore mass, estimates. The complicated structure of the velocity distribution, with many peaks, suggests however that we indeed deal with multiple objects situated along the line of sight. On the other hand the simulations show that close neighbours, within one virial radius from each other (like mergers or infalling subhalo), have similar velocities and would produce regular, one-peak velocity distributions.
Given the multi-peak velocity structure of A1689 we conclude that the groups of galaxies with $`\pm 4000`$ km s<sup>-1</sup> (and of course also the more discrepant group at $`+8000`$ km s<sup>-1</sup>) with respect to the cluster mean velocity are probably separate structures not associated with the cluster, but aligned along the line of sight. If the velocities $`\pm 4000`$ km s<sup>-1</sup> are due mainly to the Hubble flow these structures are located at about 60 Mpc from the cluster. For a cluster mass of $`2\times 10^{15}M_{\mathrm{}}`$ the distance corresponds to about 17 virial radii. This would mean that the structures do not affect the cluster dynamically and cannot be responsible for any departures from equilibrium. This rather complex structure in velocity does not necessarily translate itself into the X-ray gas distribution that can appear regular and smooth. This is indeed the case for A1689, in which the morphology of the X-ray data is commonly interpreted as a clear indication of the relaxed state of the cluster. Andersson & Madejski (2004) demonstrate however that in the case of two similar clusters aligned along the line of sight X-ray data can easily underestimate the mass by a factor of 2.
The presence of foreground and background structures in the line of sight of A1689 may affect the path of light coming from the lensed galaxies. In their study of strong lensing in A1689 Broadhurst et al. (2005a) managed to subtract the neighbouring structure from the main lensing signal. However, no such correction was made when the analysis was extended by Broadhurst et al. (2005b) to weak lensing and larger distances from the cluster centre. Hoekstra (2003) has shown that even the presence of distant large-scale structure in the Universe can affect the weak lensing signal, significantly increasing the uncertainty in the estimated parameters for a given cluster. A similar increase in the estimated errors was shown to be the case if the cluster departs from spherical symmetry (Oguri et al. 2005). It would be interesting to verify whether objects in the vicinity of the cluster could cause similar effect, thereby decreasing the claimed discrepancy with CDM structures, especially when the weak lensing signal is very low and the inferred surface mass distribution in the outer regions very uncertain, as in the case of A1689 (Broadhurst et al. 2005b).
## Acknowledgements
We wish to thank T. Broadhurst, H. Hoekstra and the anonymous referee for their comments on the paper. Computer simulations presented in this paper were performed at the Leibnizrechenzentrum (LRZ) in Munich. EŁ is grateful for the hospitality of Instituto de Astrofísica de Andalucia in Granada where part of this work was done. RW acknowledges the summer student program at Copernicus Center. This research has made use of the NASA/IPAC Extragalactic Database (NED) operated by the Jet Propulsion Laboratory. This work was partially supported by the Polish Ministry of Scientific Research and Information Technology under grant 1P03D02726 and the exchange program of CSIC/PAN. |
warning/0507/gr-qc0507008.html | ar5iv | text | # Cylindrically symmetric wormholes
## 1. Introduction
Wormholes may be defined as handles or tunnels linking different universes or widely separated regions of our own Universe. That such wormholes may be traversable by humanoid travelers was first conjectured by Morris and Thorne in 1988 and has led to a flurry of activity that has continued to the present. For a summary of the more recent developments see .
The wormholes discussed by Morris and Thorne (MT) are assumed to be spherically symmetric. It is implicitly assumed that the throat is a smooth surface. A generic feature of static wormholes, whether spherically symmetric or not, is the violation of the weak energy condition.
In this paper we propose a wormhole that is different from an MT wormhole, starting with the assumption of cylindrical symmetry. This assumption allows considerably more freedom in choosing the metric coefficients.
It will be shown that the “shape function” $`b=b(\rho ,z)`$, given in Eq. (2) below, must be a function of $`\rho `$ alone. For physical reasons $`b=b(\rho )`$, being related to the mass of the wormhole between any two $`z`$ values, has to be a continuous function of $`\rho `$. Again for physical reasons, moving in the $`z`$ direction instead, $`b(\rho )`$ can be changed abruply at some $`z`$, effectively replacing one layer of wormhole material by another. This results in a jump discontinuity at this $`z`$ value. We will therefore replace the partial derivative by a one-sided derivative. The main conclusion is that for this type of wormhole, together with the use of a one-sided derivative, the weak energy condition need not be violated.
## 2. The solution
Following Islam , p. 20, a cylindrically symmetric static metric can be put into the form
(1)
$$ds^2=fdt^2+l\rho ^2d\theta ^2+Ad\rho ^2+2Bd\rho dz+Cdz^2$$
(employing the Lorentzian signature), where $`f,l,A,B`$, and $`C`$ are all functions of $`\rho `$ and $`z`$. Using the transformation $`\rho ^{}=F(\rho ,z)`$ and $`z^{}=G(\rho ,z)`$, we get the following form by finding $`d\rho ^{}`$ and $`dz^{}`$, solving for $`d\rho `$ and $`dz`$, and substituting in Eq. (1):
$$\begin{array}{c}ds^2=fdt^2+l\rho ^2d\theta ^2+J^2\{(AG_2^22BG_1G_2+CG_1^2)(d\rho ^{})^2\hfill \\ \hfill +2\left[AG_2F_2+B(G_2F_1+G_1F_2)CG_1F_1\right]d\rho ^{}dz^{}\\ \hfill +(AF_2^22BF_1F_2+CF_1^2)(dz^{})^2\},\end{array}$$
where the subscripts denote partial derivatives (e.g., $`F_1=F/\rho `$) and $`J=F_1G_2F_2G_1`$ is the Jacobian of the transformation. Since the functions $`F`$ and $`G`$ are arbitrary, we can let
$$AG_2F_2+B(G_2F_1+G_1F_2)CG_1F_1=0$$
and
$$J^2(AF_2^22BF_1F_2+CF_1^2)=l.$$
Assuming there exists a nontrivial solution, the metric can be written in the form (omitting primes)
$$ds^2=fdt^2+H(\rho ,z)d\rho ^2+l(\rho ^2d\theta ^2+dz^2).$$
In the spirit of Morris and Thorne , we will write the metric in the form
(2)
$$ds^2=e^{2\mathrm{\Phi }(\rho ,z)}dt^2+\frac{1}{1b(\rho ,z)/\rho }d\rho ^2+\left[K(\rho ,z)\right]^2(\rho ^2d\theta ^2+dz^2),$$
where $`\mathrm{\Phi }=\mathrm{\Phi }(\rho ,z)`$ is called the *redshift function* and $`b=b(\rho ,z)`$ the *shape function*. By the usual assumption of asymptotic flatness we require that $`\mathrm{\Phi }=\mathrm{\Phi }(\rho ,z)`$ as well as its partial derivatives approach zero as $`\rho ,z\mathrm{}`$. The function $`K(\rho ,z)`$ is a positive nondecreasing function of $`\rho `$ that determines the proper radial distance from the origin, so that $`(/\rho )K(\rho ,z)>0`$. Thus $`2\pi \rho K(\rho ,z)`$ is the proper circumference of the circle passing through $`(\rho ,z)`$. (A similar function appears in the metric describing a rotating wormhole .) Apart from these requirements we assume that $`\mathrm{\Phi }`$, $`b`$, and $`K`$ can be freely assigned to meet the desired physical requirements of the wormhole.
To see how to best interpret the shape function $`b=b(\rho ,z)`$, we need to calculate the nonzero components of the Einstein tensor in the orthonormal frame. These are listed next:
(3)
$$\begin{array}{c}G_{\widehat{t}\widehat{t}}=\left(\frac{1}{\rho ^2K(\rho ,z)}\frac{K(\rho ,z)}{\rho }+\frac{1}{2\rho ^3}\right)\left(\rho \frac{b(\rho ,z)}{\rho }b(\rho ,z)\right)\hfill \\ \hfill +\frac{1}{[K(\rho ,z)]^4}\left(\frac{K(\rho ,z)}{z}\right)^2\frac{1}{[K(\rho ,z)]^3}\frac{^2K(\rho ,z)}{z^2}\\ \hfill +(1\frac{b(\rho ,z)}{\rho })\times \\ \hfill \left[\frac{3}{\rho K(\rho ,z)}\frac{K(\rho ,z)}{\rho }\frac{2}{K(\rho ,z)}\frac{^2K(\rho ,z)}{\rho ^2}\frac{1}{[K(\rho ,z)]^2}\left(\frac{K(\rho ,z)}{\rho }\right)^2\right]\\ \hfill \frac{3}{4\rho ^2[K(\rho ,z)]^2}\left(1\frac{b(\rho ,z)}{\rho }\right)^2\left(\frac{b(\rho ,z)}{z}\right)^2\\ \hfill \frac{1}{2\rho [K(\rho ,z)]^2}\left(1\frac{b(\rho ,z)}{\rho }\right)^1\frac{^2b(\rho ,z)}{z^2},\end{array}$$
(4)
$$\begin{array}{c}G_{\widehat{\rho }\widehat{\rho }}=\frac{1}{[K(\rho ,z)]^2}\left[\frac{^2\mathrm{\Phi }(\rho ,z)}{z^2}+\left(\frac{\mathrm{\Phi }(\rho ,z)}{z}\right)^2\right]\hfill \\ \hfill \frac{1}{[K(\rho ,z)]^4}\left(\frac{K(\rho ,z)}{z}\right)^2+\frac{1}{[K(\rho ,z)]^3}\frac{^2K(\rho ,z)}{z^2}\\ \hfill +(1\frac{b(\rho ,z)}{\rho })\times \\ \hfill [\frac{2}{K(\rho ,z)}\frac{\mathrm{\Phi }(\rho ,z)}{\rho }\frac{K(\rho ,z)}{\rho }+\frac{1}{\rho }\frac{\mathrm{\Phi }(\rho ,z)}{\rho }\\ \hfill +\frac{1}{[K(\rho ,z)]^2}\left(\frac{K(\rho ,z)}{\rho }\right)^2+\frac{1}{\rho K(\rho ,z)}\frac{K(\rho ,z)}{\rho }],\end{array}$$
(5)
$$\begin{array}{c}G_{\widehat{\theta }\widehat{\theta }}=\left(\frac{1}{2\rho ^2K(\rho ,z)}\frac{K(\rho ,z)}{\rho }\frac{1}{2\rho ^2}\frac{\mathrm{\Phi }(\rho ,z)}{\rho }\right)\left(\rho \frac{b(\rho ,z)}{\rho }b(\rho ,z)\right)\hfill \\ \hfill \frac{1}{[K(\rho ,z)]^3}\frac{\mathrm{\Phi }(\rho ,z)}{z}\frac{K(\rho ,z)}{z}+\frac{1}{[K(\rho ,z)]^2}\left[\frac{^2\mathrm{\Phi }(\rho ,z)}{z^2}+\left(\frac{\mathrm{\Phi }(\rho ,z}{z}\right)^2\right]\\ \hfill +(1\frac{b(\rho ,z)}{\rho })\times [\frac{^2\mathrm{\Phi }(\rho ,z)}{\rho ^2}+\left(\frac{\mathrm{\Phi }(\rho ,z)}{\rho }\right)^2\\ \hfill +\frac{1}{K(\rho ,z)}\frac{\mathrm{\Phi }(\rho ,z)}{\rho }\frac{K(\rho ,z)}{\rho }+\frac{1}{K(\rho ,z)}\frac{^2K(\rho ,z)}{\rho ^2}]\\ \hfill +\frac{3}{4\rho ^2[K(\rho ,z)]^2}\left(1\frac{b(\rho ,z)}{\rho }\right)^2\left(\frac{b(\rho ,z)}{z}\right)^2\\ \hfill +\frac{1}{2\rho [K(\rho ,z)]^2}\left(1\frac{b(\rho ,z)}{\rho }\right)^1\frac{\mathrm{\Phi }(\rho ,z)}{z}\frac{b(\rho ,z)}{z}\\ \hfill \frac{1}{2\rho [K(\rho ,z)]^3}\left(1\frac{b(\rho ,z)}{\rho }\right)^1\frac{K(\rho ,z)}{z}\frac{b(\rho ,z)}{z}\\ \hfill +\frac{1}{2\rho [K(\rho ,z)]^2}\left(1\frac{b(\rho ,z)}{\rho }\right)^1\frac{^2b(\rho ,z)}{z^2},\end{array}$$
(6)
$$\begin{array}{c}G_{\widehat{z}\widehat{z}}=\hfill \\ \hfill \left(\frac{1}{2\rho ^2}\frac{\mathrm{\Phi }(\rho ,z)}{\rho }\frac{1}{2\rho ^2K(\rho ,z)}\frac{K(\rho ,z)}{\rho }\frac{1}{2\rho ^3}\right)\left(\rho \frac{b(\rho ,z)}{\rho }b(\rho ,z)\right)\\ \hfill +\frac{1}{[K(\rho ,z)]^3}\frac{\mathrm{\Phi }(\rho ,z)}{z}\frac{K(\rho ,z)}{z}\\ \hfill +(1\frac{b(\rho ,z)}{\rho })\times [\frac{^2\mathrm{\Phi }(\rho ,z)}{\rho ^2}+\left(\frac{\mathrm{\Phi }(\rho ,z)}{\rho }\right)^2\\ \hfill +\frac{1}{K(\rho ,z)}\frac{\mathrm{\Phi }(\rho ,z)}{\rho }\frac{K(\rho ,z)}{\rho }+\frac{1}{\rho }\frac{\mathrm{\Phi }(\rho ,z)}{\rho }+\frac{2}{\rho K(\rho ,z)}\frac{K(\rho ,z)}{\rho }\\ \hfill +\frac{1}{K(\rho ,z)}\frac{^2K(\rho ,z)}{\rho ^2}]\\ \hfill +\frac{1}{2\rho [K(\rho ,z)]^3}\left(1\frac{b(\rho ,z)}{\rho }\right)^1\frac{K(\rho ,z)}{z}\frac{b(\rho ,z)}{z}\\ \hfill +\frac{1}{2\rho [K(\rho ,z)]^2}\left(1\frac{b(\rho ,z)}{\rho }\right)^1\frac{\mathrm{\Phi }(\rho ,z)}{z}\frac{b(\rho ,z)}{z},\end{array}$$
(7)
$$\begin{array}{c}G_{\widehat{\rho }\widehat{z}}=(1\frac{b(\rho ,z)}{\rho })^{1/2}\times \hfill \\ \hfill [\frac{1}{[K(\rho ,z)]^2}\frac{K(\rho ,z)}{\rho }\frac{\mathrm{\Phi }(\rho ,z)}{z}\frac{1}{K(\rho ,z)}\frac{^2\mathrm{\Phi }(\rho ,z)}{\rho z}\\ \hfill \frac{1}{K(\rho ,z)}\frac{\mathrm{\Phi }(\rho ,z)}{z}\frac{\mathrm{\Phi }(\rho ,z)}{\rho }\frac{1}{[K(\rho ,z)]^2}\frac{^2K(\rho ,z)}{\rho z}\\ \hfill \frac{1}{\rho [K(\rho ,z)]^2}\frac{K(\rho ,z)}{z}+\frac{1}{[K(\rho ,z)]^3}\frac{K(\rho ,z)}{\rho }\frac{K(\rho ,z)}{z}]\\ \hfill +(1\frac{b(\rho ,z)}{\rho })^{1/2}\times [\frac{1}{2\rho K(\rho ,z)}\frac{\mathrm{\Phi }(\rho ,z)}{\rho }\frac{b(\rho ,z)}{z}\\ \hfill +\frac{1}{2\rho [K(\rho ,z)]^2}\frac{K(\rho ,z)}{\rho }\frac{b(\rho ,z)}{z}+\frac{1}{2\rho ^2K(\rho ,z)}\frac{b(\rho ,z)}{z}].\end{array}$$
We would like $`b=b(\rho ,z)`$ to correspond to the throat of the wormhole for some $`\rho `$ and $`z`$. But in that case the fraction $`1/(1b(\rho ,z)/\rho )`$ is undefined at the throat. In particular, $`T_{\widehat{t}\widehat{t}}=\rho =(1/8\pi )G_{\widehat{t}\widehat{t}}`$ is undefined. A remarkable feature of the solution is the following: the expressions
$$\frac{1}{1b(\rho ,z)/\rho }\text{and}\frac{b(\rho ,z)}{z}$$
always occur together. So we must require that $`b(\rho ,z)/z=0`$, i.e., $`b`$ must be independent of $`z`$. Accordingly, we will write $`b=b(\rho )`$ from now on and omit all terms containing the factor $`b(\rho ,z)/z`$.
These observations allow us to interpret $`b=b(\rho )`$ in terms of the usual embedding diagram, such as Fig. 1 in MT , by letting $`t`$ be a fixed moment in time in Eq. (2) and then choosing the slice $`z=0`$. In MT every circle represents a sphere because of the assumption of spherical symmetry. In our case, every circle represents a cylinder, so that the throat, or part of the throat, is a cylindrical surface with minimal radius $`\rho =\rho _0`$ extending along the $`z`$ axis. Just how far the throat extends depends on the energy conditions, discussed in the next section.
## 3. The weak energy condition
Recall that the weak energy condition (WEC) can be stated as $`T_{\widehat{\alpha }\widehat{\beta }}\mu ^{\widehat{\alpha }}\mu ^{\widehat{\beta }}0`$ for all timelike and, by continuity, all null vectors and where $`T_{\widehat{\alpha }\widehat{\beta }}`$ are the components of the stress-energy tensor in the orthonormal frame.
Since the stress-energy tensor has the same algebraic structure as the Einstein tensor, we have (from the Einstein field equations $`8\pi T_{\widehat{\alpha }\widehat{\beta }}=G_{\widehat{\alpha }\widehat{\beta }}`$)
$$\begin{array}{ccc}T_{\widehat{t}\widehat{t}}=\rho =\frac{1}{8\pi }G_{\widehat{t}\widehat{t}},& T_{\widehat{\rho }\widehat{\rho }}=\tau =\frac{1}{8\pi }G_{\widehat{\rho }\widehat{\rho }},& \\ T_{\widehat{\theta }\widehat{\theta }}=\frac{1}{8\pi }G_{\widehat{\theta }\widehat{\theta }},& T_{\widehat{z}\widehat{z}}=\frac{1}{8\pi }G_{\widehat{z}\widehat{z}},& T_{\widehat{\rho }\widehat{z}}=\frac{1}{8\pi }G_{\widehat{\rho }\widehat{z}}.\end{array}$$
For the case of a diagonal stress-energy tensor the WEC can be written
(8)
$$G_{\widehat{t}\widehat{t}}0,G_{\widehat{t}\widehat{t}}+G_{\widehat{\rho }\widehat{\rho }}0,G_{\widehat{t}\widehat{t}}+G_{\widehat{\theta }\widehat{\theta }}0,G_{\widehat{t}\widehat{t}}+G_{\widehat{z}\widehat{z}}0$$
corresponding to the timelike vector $`(1,0,0,0)`$ and the null vectors $`(1,1,0,0)`$, $`(1,0,1,0)`$, and $`(1,0,0,1)`$, respectively. Because of the off-diagonal element $`T_{\widehat{\rho }\widehat{z}}`$, we also need to consider the null vector
$$(1,\frac{1}{\sqrt{2}},0,\frac{1}{\sqrt{2}}),$$
which yields
(9)
$$\frac{1}{2}(G_{\widehat{t}\widehat{t}}+G_{\widehat{\rho }\widehat{\rho }})+\frac{1}{2}(G_{\widehat{t}\widehat{t}}+G_{\widehat{z}\widehat{z}})+G_{\widehat{\rho }\widehat{z}}0.$$
Suppose we start the investigation with the second energy condition:
(10)
$$\begin{array}{c}G_{\widehat{t}\widehat{t}}+G_{\widehat{\rho }\widehat{\rho }}=\hfill \\ \hfill \left(\frac{1}{\rho ^2K(\rho ,z)}\frac{K(\rho ,z)}{\rho }\right)\left(\rho \frac{db(\rho )}{d\rho }b(\rho )\right)+\frac{1}{2\rho ^3}\left(\rho \frac{db(\rho )}{d\rho }b(\rho )\right)\\ \hfill +\frac{1}{[K(\rho ,z)]^2}\left[\frac{^2\mathrm{\Phi }(\rho ,z)}{z^2}+\left(\frac{\mathrm{\Phi }(\rho ,z)}{z}\right)^2\right]\\ \hfill +(1\frac{b(\rho )}{\rho })\times [\frac{2}{K(\rho ,z)}\frac{\mathrm{\Phi }(\rho ,z)}{\rho }\frac{K(\rho ,z)}{\rho }+\frac{1}{\rho }\frac{\mathrm{\Phi }(\rho ,z)}{\rho }\\ \hfill \frac{2}{\rho K(\rho ,z)}\frac{K(\rho ,z)}{\rho }\frac{2}{K(\rho ,z)}\frac{^2K(\rho ,z)}{\rho ^2}];\end{array}$$
(recall that $`b(\rho ,z)/z=0.`$)
The factor
(11)
$$\rho db(\rho )/d\rho b(\rho )$$
is negative at the throat, since $`b(\rho _0)=\rho _0`$ and $`b^{}(\rho _0)<1.`$ It follows that the second term
$$\frac{1}{2\rho ^3}\left(\rho \frac{db(\rho )}{d\rho }b(\rho )\right)$$
is negative with an absolute value equal to a multiple of $`1/\rho _0^2`$. The first term has a smaller absolute value, as we will see below (Expression (15)).
At the throat, $`1b(\rho )/\rho =0.`$ Since the partial derivatives of $`\mathrm{\Phi }=\mathrm{\Phi }(\rho ,z)`$ go to zero as $`z\mathrm{}`$, $`G_{\widehat{t}\widehat{t}}+G_{\widehat{\rho }\widehat{\rho }}`$ eventually becomes negative. So we need to cut off the wormhole region at some $`z`$ value above the equatorial plane $`z=0`$. (Since the analysis for the lower region would follow similar lines, we will concentrate exclusively on the upper region.) For proper choices of $`\mathrm{\Phi }(\rho ,z)`$ and $`K(\rho ,z)`$ it should be possible to make the sum of the first three terms in Eq. (10) nonnegative below this $`z`$ value. Even though we have great freedom in choosing $`\mathrm{\Phi }`$ and $`K`$, the existence of such functions cannot be taken for granted. So let us consider the following forms:
$$\mathrm{\Phi }(\rho ,z)=A(\rho )G(z)\text{and}K(\rho ,z)=1+\frac{F(z)}{\rho ^a},$$
where $`0<a1`$.
Qualitatively, the basic shapes are given in Figures 1 and 2: $`A(\rho )`$ is concave up everywhere with a vertical asymptote at $`\rho =ϵ>0`$, well to the left of $`\rho =\rho _0`$, $`F(z)`$ is concave down and $`G(z)`$ concave up on an interval that extends well beyond $`z=z_1`$, where $`F^{}(z_1)=G^{}(z_1)=0`$, to some $`z=z_2`$. Both $`ϵ`$ and $`z_2`$ are to be determined later.
The partial derivatives are:
$$\frac{\mathrm{\Phi }(\rho ,z)}{\rho }=A^{}(\rho )G(z)>0\text{for}0<z<z_2,$$
$$\frac{^2\mathrm{\Phi }(\rho ,z)}{\rho ^2}=A^{\prime \prime }(\rho )G(z)<0\text{for}0<z<z_2,$$
$$\frac{\mathrm{\Phi }(\rho ,z)}{z}=A(\rho )G^{}(z)<0\text{for}0<z<z_1,$$
$$\frac{^2\mathrm{\Phi }(\rho ,z)}{z^2}=A(\rho )G^{\prime \prime }(z)>0\text{for}0<z<z_2,$$
$$\frac{K(\rho ,z)}{\rho }=\frac{aF(z)}{\rho ^{a+1}}<0\text{for}0<z<z_2,$$
$$\frac{^2K(\rho ,z)}{\rho ^2}=\frac{a(a+1)F(z)}{\rho ^{a+2}}>0\text{for}0<z<z_2,$$
$$\frac{K(\rho ,z)}{z}=\frac{F^{}(z)}{\rho ^a}>0\text{for}0<z<z_1,$$
and
$$\frac{^2K(\rho ,z)}{z^2}=\frac{F^{\prime \prime }(z)}{\rho ^a}<0\text{for}0<z<z_2.$$
On the interval $`(z_1,z_2)`$, $`G^{}(z)`$ and $`F^{}(z)`$ change signs.
On the interval $`(0,z_2)`$ it seems sufficient for now to choose $`A`$ so that
(12)
$$\begin{array}{c}\frac{1}{[K(\rho ,z)]^2}\left[\frac{^2\mathrm{\Phi }(\rho ,z)}{z^2}+\left(\frac{\mathrm{\Phi }(\rho ,z)}{z}\right)^2\right]=\hfill \\ \hfill \frac{1}{\left(1+\frac{F(z)}{\rho ^a}\right)^2}\left\{A(\rho )G^{\prime \prime }(z)+\left[A(\rho )G^{}(z)\right]^2\right\}\end{array}$$
exceeds the first two terms (of order $`1/\rho _0^2`$) in Eq. (10). This can always be accomplished by “raising” $`A(\rho )`$ sufficiently, independently of $`G`$ and $`F`$.
As already noted, $`1b(\rho )/\rho `$ is zero at the throat. In some wormhole solutions the violation actually occurs near the throat, rather than at the throat. For the second part of the right-hand side of Eq. (10), we get
(13)
$$\begin{array}{c}(1\frac{b(\rho )}{\rho })[\frac{2a}{\rho }\frac{\frac{F(z)}{\rho ^a}}{1+\frac{F(z)}{\rho ^a}}A^{}(\rho )G(z)\hfill \\ \hfill +\frac{1}{\rho }A^{}(\rho )G(z)\frac{2a^2}{\rho ^2}\frac{\frac{F(z)}{\rho ^a}}{1+\frac{F(z)}{\rho ^a}}].\end{array}$$
As long as $`a`$ is small, the positive second term dominates on the interval $`(0,z_2)`$ for any $`F(z)`$ and for a wide variety of choices of $`A`$ and $`G`$. (We will see later that $`A^{}(\rho )G(z)`$ has to be relatively large to begin with.)
For the first energy condition, $`G_{\widehat{t}\widehat{t}}0`$, we need to specify $`K(\rho ,z)=1+F(z)/\rho ^a`$ more precisely: the first part of $`G_{\widehat{t}\widehat{t}}`$,
(14)
$$\begin{array}{c}\left(\frac{1}{\rho ^2K(\rho ,z)}\frac{K(\rho ,z)}{\rho }+\frac{1}{2\rho ^3}\right)\left(\rho \frac{db(\rho )}{d\rho }b(\rho )\right)\hfill \\ \hfill +\frac{1}{\left(1+\frac{F(z)}{\rho ^a}\right)^4}\left(\frac{F^{}(z)}{\rho ^a}\right)^2\frac{1}{\left(1+\frac{F(z)}{\rho ^a}\right)^3}\frac{F^{\prime \prime }(z)}{\rho ^a}\end{array}$$
must be nonnegative on the interval $`(0,z_2)`$. (Observe that the last two terms are positive on this interval.) That such a function can be constructed can be seen by starting with the rough straight-line approximation for $`F(z)`$, the left part extending from $`(0,b)`$ to a peak at $`(z_1,b+c)`$:
$$K(\rho ,z)=1+\frac{(c/z_1)z+b}{\rho ^a}.$$
Expression (14) becomes
(15)
$$\begin{array}{c}\left(\frac{a}{\rho ^3}\frac{\frac{F(z)}{\rho ^a}}{1+\frac{F(z)}{\rho ^a}}+\frac{1}{2\rho ^3}\right)\left(\rho \frac{db(\rho )}{d\rho }b(\rho )\right)+\frac{1}{\left(1+\frac{(c/z_1)z+b}{\rho ^a}\right)^4}\left(\frac{c/z_1}{\rho ^a}\right)^2\hfill \end{array}$$
The first term is of order $`1/\rho _0^2`$, as we saw earlier. The second term is smallest at $`z=z_1`$. To match the first term, we need (at least roughly)
$$\frac{1}{\left(1+\frac{c+b}{\rho ^a}\right)^2}\frac{c/z_1}{\rho ^a}=\frac{1}{\rho _0}.$$
This is a quadratic equation in $`c`$ with infinitely many real solutions in terms of the other parameters. Finally, the corner at $`z=z_1`$ can be replaced by a small arc with an arbitrarily large curvature $`\kappa `$, where $`\kappa =F^{\prime \prime }(z)/\left[1+(F^{}(z))^2\right]^{3/2}F^{\prime \prime }(z)`$ near $`z=z_1`$.
Away from the throat we obtain the remaining part of $`G_{\widehat{t}\widehat{t}}`$ by returning to $`K(\rho ,t)=1+F(z)/\rho ^a`$ and recalling that $`b(\rho ,z)/z=0`$:
$$\left(1\frac{b(\rho )}{\rho }\right)\frac{\frac{F(z)}{\rho ^a}}{1+\frac{F(z)}{\rho ^a}}\left(\frac{a2a^2}{\rho ^2}\frac{a^2\frac{F(z)}{\rho ^a}}{\left(1+\frac{F(z)}{\rho ^a}\right)\rho ^2}\right).$$
For small $`a`$, this expressions is positive for *any* $`F(z)`$.
## 4. The remaining energy conditions
The energy conditions $`G_{\widehat{t}\widehat{t}}+G_{\widehat{\theta }\widehat{\theta }}0`$ and $`G_{\widehat{t}\widehat{t}}+G_{\widehat{z}\widehat{z}}0`$ can be checked in a similar way with the understanding that some refinements may still have to be made. For example, in the condition $`G_{\widehat{t}\widehat{t}}+G_{\widehat{z}\widehat{z}}0`$ the terms involving the “negative energy expression”(11), $`\rho db(\rho )/d\rho b(\rho )`$, are actually positive. Unfortunately, there is also the dangerous negative term
$$\frac{1}{[K(\rho ,z)]^3}\frac{\mathrm{\Phi }(\rho ,z)}{z}\frac{K(\rho ,z)}{z}$$
on the entire interval $`(0,z_2)`$ since $`G^{}(z)`$ and $`F^{}(z)`$ have opposite signs. (In the condition $`G_{\widehat{t}\widehat{t}}+G_{\widehat{\theta }\widehat{\theta }}0`$ this term has the opposite sign, so that this problem does not arise.) The terms not involving $`1b(\rho )/\rho `$ are
(16)
$$\begin{array}{c}\left(\frac{1}{2\rho ^2K(\rho ,z)}\frac{K(\rho ,z)}{\rho }\frac{1}{2\rho ^2}\frac{\mathrm{\Phi }(\rho ,z)}{\rho }\right)\left(\rho \frac{db(\rho )}{d\rho }b(\rho )\right)\hfill \\ \hfill +\frac{1}{[K(\rho ,z)]^3}\frac{\mathrm{\Phi }(\rho ,z)}{z}\frac{K(\rho ,z)}{z}+\text{two positive terms}=\\ \hfill \left(\frac{1}{2\rho ^2}\frac{1}{1+\frac{F(z)}{\rho ^a}}\frac{aF(z)}{\rho ^{a+1}}\frac{1}{2\rho ^2}A^{}(\rho )G(z)\right)\left(\rho \frac{db(\rho )}{d\rho }b(\rho )\right)\\ \hfill +\frac{1}{\left(1+\frac{F(z)}{\rho ^a}\right)^3}A(\rho )G^{}(z)\frac{F^{}(z)}{\rho ^a}+\text{two positive terms}.\end{array}$$
Since $`\rho db(\rho )/d\rho b(\rho )`$ is negative, the first term on the right-hand side of Eq. (16) is indeed positive. The second term is negative, close to zero near $`z=z_1`$, but increasing in absolute value as $`z0`$. Since $`A(\rho )G^{}(z)`$ already occurred in Eq. (12), we may have only limited control over this quantity. Fortunately, we can adjust the competing term involving $`A^{}(\rho )G(z)`$: $`A(\rho )`$ can always be chosen with larger $`A^{}(\rho )`$, especially near the throat, but, more importantly, $`\left|G(z)\right|`$ can be made as large as we please by “lowering” $`G(z)`$. \[That $`A^{}(\rho )G(z)`$ is relatively large was already mentioned after Eq. (13).\]
The remaining terms have the form
(17)
$$\begin{array}{c}(1\frac{b(\rho )}{\rho })\{A^{\prime \prime }(\rho )G(z)+[A^{}(\rho )G(z)]^2\hfill \\ \hfill +\frac{1}{\rho }A^{}(\rho )G(z)(1aL(z))\frac{a^2}{\rho ^2}L(z)(1+L(z))\},\end{array}$$
where
$$L(z)=\frac{\frac{F(z)}{\rho ^a}}{1+\frac{F(z)}{\rho ^a}}.$$
At the throat the entire expression is equal to zero. Moving away from the throat, since $`A^{\prime \prime }(\rho )G(z)`$ is negative, care must be taken to keep $`A^{\prime \prime }(\rho )`$ small enough so that the entire sum remains positive.
In the last energy condition (9),
$$\frac{1}{2}(G_{\widehat{t}\widehat{t}}+G_{\widehat{\rho }\widehat{\rho }})+\frac{1}{2}(G_{\widehat{t}\widehat{t}}+G_{\widehat{z}\widehat{z}})+G_{\widehat{\rho }\widehat{z}}0,$$
the first two terms have already been taken care of. The last term is zero at the throat. Away from the throat the heavily dominating $`A^{}(\rho )G(z)`$ contributed by the first two terms carries the day.
## 5. The upper layers
So far we have considered only the region from $`z=0`$ to $`z=z_2`$ and found that the WEC need not be violated on this interval. The next task is to find a convenient value for $`z_2`$ at which to cut off and replace the wormhole material by a transitional layer without introducing any new energy condition violation. The easiest way is to start with Eq. (16) and to assume for the time being that $`b0`$, thereby cutting off the wormhole material. (The transitional layer will be introduced below.) The first term on the right-hand side is zero. Now choose $`z_2`$ so that $`G^{}(z_2)`$ is small enough to keep the right-hand side positive. Again for the sake of simplicity, replace $`G`$ by a straight line with slope $`G^{}(z_2)`$ for $`z>z_2`$ (or by a curve with a very small curvature), thereby retaining continuity. The line will eventually reach zero at some $`z=z_3`$. $`F(z)`$ should also have a small curvature for $`z>z_2`$. Observe that as $`F(z)0`$, $`K(\rho ,z)`$ approaches unity, so that the metric itself approaches that of a flat Minkowski spacetime.
Since $`b0`$, the “negative energy expression” (11), $`\rho db(\rho )/d\rho b(\rho )`$, is also zero, so that the energy conditions are satisfied *a fortiori*. \[The remaining terms, particularly terms involving $`1b(\rho )/\rho `$, do not present any special problems apart from some minor fine tuning: Expressions (13) and (17) suggest that $`F(z)`$ should reach zero before $`G(z)`$. Fortunately, $`|G(z)|`$ is assumed to be large to begin with and to decrease relatively slowly.\]
The final step is to introduce a transitional layer between $`z=z_2`$ and $`z=z_3`$: instead of letting $`b0`$, let $`b=ϵ>0`$, where $`ϵ`$ is so small that the resulting “negative energy expression” is sufficiently close to zero to leave the above comments unaffected. According to line element (2), $`\rho =ϵ`$ now becomes the Schwarzschild radius: since $`A(\rho )`$ has a vertical asymptote at $`\rho =ϵ`$, $`e^{2\mathrm{\Phi }}0`$ as $`\rho ϵ+`$ for all $`z`$, thereby creating an event horizon. So the transitional layer is a Schwarzschild spacetime.
## 6. The throat
To complete the discussion, we need to take a closer look at the throat. So far we considered only the cylindrical surface $`\rho =\rho _0`$ from $`z=0`$ to $`z=z_2`$, the exact analogue of $`r=r_0`$ for the spherically symmetric case. Since the upper layers create a flat top at $`z=z_2`$, there is no violation of the WEC (Visser , chapter 15). But unlike Visser’s cubical wormholes, the edges cannot be rounded off, as this would violate the condition $`b(\rho ,z)/z=0`$. It is precisely this rounding off that causes the violation of the WEC for cubical wormholes.
While keeping the sharp edge does solve one problem, it results in another: by changing the shape function at $`z=z_2`$ (and again at $`z=z_3`$), $`b`$ depends on $`z`$ after all. In fact, there is a jump discontinuity at $`z=z_2`$, so that the partial derivative with respect to $`z`$ does not exist. As a result, the entire solution breaks down at this value. But as noted in the introduction, we are going to make a small change in one of the assumptions, replacing the regular derivative by a one-sided derivative: on any closed interval $`[z_2,z_2^{}]`$ the right-hand derivative is
$$\frac{b(\rho ,z_2+)}{z}=\underset{zz_2+}{lim}\frac{b(\rho ,z)b(\rho ,z_2)}{zz_2}=0.$$
On the open interval $`(0,z_2)`$ the function is everywhere differentiable with respect to $`z`$ and its partial derivative is also equal to zero. With this modification we can retain the requirement
$$\frac{b(\rho ,z)}{z}=0$$
and hence the earlier analysis. Observe that the throat is a piecewise smooth surface.
While the wormholes described here are not likely to occur naturally, a traversable wormhole that does not violate the WEC could in principle be constructed by an advanced cicilization.
## 7. Traversability
Our final topic is a brief look at traversability conditions. Since our variable $`\rho `$ is analogous to $`r`$ in the spherical case, we assume that the traveler approaches the throat along a path perpendicular to the $`z`$ axis. This is best analyzed with the aid of an orthonormal basis in the traveler’s frame:
$$e_{\widehat{0}^{}}=\mu =\gamma e_{\widehat{t}}\gamma \left(\frac{v}{c}\right)e_{\widehat{\rho }},e_{\widehat{1}^{}}=\gamma e_{\widehat{\rho }}+\gamma \left(\frac{v}{c}\right)e_{\widehat{t}},e_{\widehat{2}^{}}=e_{\widehat{\theta }},e_{\widehat{3}^{}}=e_{\widehat{z}}.$$
Here $`\mu `$ is the traveler’s four-velocity, while $`e_{\widehat{1}^{}}`$ points in the direction of travel. (See also Ref. .)
Since $`\rho `$ is analogous to $`r`$, the constraint $`\left|R_{\widehat{1}^{}\widehat{0}^{}\widehat{1}^{}\widehat{0}^{}}\right|`$ is similar to the spherical case discussed in MT .
The other constraints show a different behavior. Thus
(18)
$$\begin{array}{c}\left|R_{\widehat{2}^{}\widehat{0}^{}\widehat{2}^{}\widehat{0}^{}}\right|=\left|\gamma ^2R_{\widehat{\theta }\widehat{t}\widehat{\theta }\widehat{t}}+\gamma ^2\left(\frac{v}{c}\right)^2R_{\widehat{\theta }\widehat{\rho }\widehat{\theta }\widehat{\rho }}\right|=\hfill \\ \hfill |\gamma ^2\frac{\mathrm{\Phi }(\rho ,z)}{\rho }(1\frac{b(\rho )}{\rho })(\frac{1}{K(\rho ,z)}\frac{K(\rho ,z)}{\rho }+\frac{1}{\rho })\\ \hfill +\frac{1}{[K(\rho ,z)]^3}\frac{\mathrm{\Phi }(\rho ,z)}{z}\frac{K(\rho ,z)}{z}|\\ \hfill +\gamma ^2\left(\frac{v}{c}\right)^2|\frac{1}{\rho K(\rho ,z)}\left(\frac{2K(\rho ,z)}{\rho }+\rho \frac{^2K(\rho ,z)}{\rho ^2}\right)\left(1\frac{b(\rho )}{\rho }\right)\\ \hfill \frac{1}{2\rho ^3K(\rho ,z)}(\rho \frac{K(\rho ,z)}{\rho }+K(\rho ,z))(\rho \frac{db(\rho )}{d\rho }b(\rho ))|.\end{array}$$
While the second part is the usual constraint on the velocity of the traveler, the first part shows that the best place to cross the throat is at $`z=z_1`$.
## 8. Further Discussion
We have seen that the price to pay for avoiding the violation of the WEC for a particular type of wormhole is the introduction of a one-sided derivative. While it is tempting to argue that mathematically speaking, the adjustment is rather minor, it does lead to another violation. Fortunately, this violation is physically defendable: we are merely replacing one type of wormhole material by another. This gives our violation the appearance of being less serious than an energy violation.
Similar questions can be asked about the possible violation of the null energy condition (NEC) and the averaged null energy condition (ANEC). The NEC, $`T_{\widehat{\alpha }\widehat{\beta }}k^{\widehat{\alpha }}k^{\widehat{\beta }}0`$ for null vectors, excuses us from considering the timelike vector $`(1,0,0,0)`$, leaving only $`\rho +p_j`$, where $`p_j`$ are the principal pressures. For the functions considered earlier, this condition is met. The ANEC is much more problematical. What is peculiar to our wormhole, however, is that in the critical non-Schwarzschild region, that is, up to the cut-off at $`z=z_2`$, the quantities $`\rho +p_j`$ are bounded away from zero. As a consequence, the ANEC, $`T_{\widehat{\alpha }\widehat{\beta }}k^{\widehat{\alpha }}k^{\widehat{\beta }}𝑑\lambda 0`$ (Ref ), has a good chance of being met (although difficult to quantify). While the price to be paid is still the same, at least the payoff has enjoyed a boost.
## 9. Conclusion
This paper discusses traversable wormholes different from those of the Morris-Thorne type. The wormhole geometry is cylindrically symmetric and the throat a cylindrical surface, a surface that is only piecewise smooth. As a result, the shape function is not differentiable at some $`z=z_2`$ due to a jump discontinuity. It is proposed that the regular derivative be replaced by a one-sided derivative at this value. It is shown that for proper choices of $`b`$, $`\mathrm{\Phi }`$, and $`K`$ in line element (2) the weak energy condition is satisfied for all timelike and null vectors. |
warning/0507/hep-ex0507090.html | ar5iv | text | # Study of the 𝑿(𝟑𝟖𝟕𝟐) and 𝒀(𝟒𝟐𝟔𝟎) in 𝑩^𝟎→𝑱/𝝍𝝅⁺𝝅⁻𝑲^𝟎 and 𝑩⁻→𝑱/𝝍𝝅⁺𝝅⁻𝑲⁻ decays
BABAR-PUB-05/038
SLAC-PUB-11370
thanks: Deceased
The BABAR Collaboration
## Abstract
We present results of a search for the $`X(3872)`$ in $`B^0X(3872)K_S^0,X(3872)J/\psi \pi ^+\pi ^{}`$, improved measurements of $`B^{}X(3872)K^{}`$, and a study of the $`J/\psi \pi ^+\pi ^{}`$ mass region above the $`X(3872)`$. We use 232 million $`B\overline{B}`$ pairs collected at the $`\mathrm{{\rm Y}}(4S)`$ resonance with the BABAR detector at the PEP-II $`e^+e^{}`$ asymmetric-energy storage rings. The results include the 90% confidence interval $`1.34\times 10^6<(B^0X(3872)K^0,XJ/\psi \pi ^+\pi ^{})<10.3\times 10^6`$ and the branching fraction $`(B^{}X(3872)K^{},XJ/\psi \pi ^+\pi ^{})=(10.1\pm 2.5\pm 1.0)\times 10^6`$. We observe a $`(2.7\pm 1.3\pm 0.2)\mathrm{MeV}/c^2`$ mass difference of the $`X(3872)`$ produced in the two decay modes. Furthermore, we search for the $`Y(4260)`$ in $`B`$ decays and set the $`95\%`$ C.L. upper limit $`(B^{}Y(4260)K^{},Y(4260)J/\psi \pi ^+\pi ^{})<2.9\times 10^5`$.
The $`X(3872)`$ was first observed in the charged $`B`$-meson decay chargeConj $`B^{}X(3872)K^{}`$, $`X(3872)J/\psi \pi ^+\pi ^{}`$ by the Belle Collaboration Choi et al. (2003). It has been confirmed by the BABAR Collaboration Aubert et al. (2004a) and observed inclusively in the same final state by the CDF and D0 collaborations Acosta et al. (2004). This narrow-width particle has a mass very near the $`D^0\overline{D}^0`$ threshold and decays into final states containing charmonium ($`J/\psi `$). The most plausible interpretation Eichten et al. (2004) was a $`1^3D_2`$ or $`1^1D_2`$ $`c\overline{c}`$ state which would be narrow since it would be forbidden to decay into open charm $`D\overline{D}`$ states. However, these candidates should have large radiative transitions into $`\chi _c`$ states that have not been observed Choi et al. (2003). Recent studies from Belle that combine angular and kinematic properties of the $`\pi ^+\pi ^{}`$ mass, strongly favor a $`J^{PC}=1^{++}`$ state newbelle . Other explanations include $`2^1P_1`$ $`c\overline{c}`$ ($`1^+`$) or $`2^3P_1`$ $`1^{++}`$ states that should be narrow, but are predicted to be about $`100\mathrm{MeV}/c^2`$ higher than the $`X(3872)`$ mass and are not expected to have a large decay rate into $`J/\psi \pi \pi `$ final states ted . Hence, the $`X(3872)`$ appears not to be a simple quark model $`q\overline{q}`$ meson state.
Many explanations have been proposed for the nature of the $`X(3872)`$. Recent interpretations include the diquark-antidiquark model Maiani et al. (2005) and the S-wave $`D^0\overline{D}^0`$ molecule model Braaten and Kusunoki (2005). The diquark-antidiquark model predicts a spectrum of $`J=0,1,2`$ particles and identifies the $`X(3872)`$ as its $`1^{++}`$ member state with the two quark combinations $`X_u=\left[cu\right]\left[\overline{c}\overline{u}\right]`$ and $`X_d=\left[cd\right]\left[\overline{c}\overline{d}\right]`$ with a mass difference $`m(X_d)m(X_u)(7\pm 2)\mathrm{MeV}/c^2`$. In addition, these two states could form mixed states that are produced in both charged and neutral $`B`$-meson decays with different masses and rates depending on the mixing angle. A search for the predicted charged partner of the $`X(3872)`$ has been addressed in a previous analysis with an upper limit that is still consistent with the model Aubert et al. (2005a). The $`D^0\overline{D}^0`$ molecule model interprets the $`X(3872)`$ as a loosely bound $`D^0\overline{D}^0`$ S-wave state that is produced in weak decays of the $`B`$-meson into $`D^0\overline{D}{}_{}{}^{0}K`$. In this picture, the S-wave molecule must form a $`J^P=1^+`$ state. From factorization, heavy-quark symmetry, and isospin symmetry, the decay $`B^0X(3872)K^0`$ is predicted to be suppressed by an order of magnitude relative to $`B^{}X(3872)K^{}`$ Braaten and Kusunoki (2005). To investigate these predictions, we present in this letter a study of the neutral mode $`B^0X(3872)K^0`$, $`X(3872)J/\psi \pi ^+\pi ^{}`$ and we analyze $`B^{}J/\psi \pi ^+\pi ^{}K^{}`$ decays with increased statistics to obtain improved measurements of $`X(3872)J/\psi \pi ^+\pi ^{}`$. In addition, we examine the higher $`J/\psi \pi ^+\pi ^{}`$ invariant mass region to search for a structure recently observed in initial state radiation (ISR) events Aubert:2005-ISR .
The data were collected with the BABAR detector at the PEP-II asymmetric-energy $`e^+e^{}`$ storage rings on the $`\mathrm{{\rm Y}}(4S)`$ resonance. The integrated luminosity of the data used in this analysis is $`211\text{ fb}^1`$; this corresponds to the production of $`(232\pm 3)\times 10^6`$ $`B\overline{B}`$ pairs.
The BABAR detector is described in detail elsewhere Aubert et al. (2002). Charged-particle trajectories are measured by a combination of a five-layer silicon vertex tracker (SVT) and a 40-layer drift chamber (DCH) in a 1.5-T solenoidal magnetic field. For charged-particle identification, we combine information from a ring-imaging Cherenkov detector (DIRC) and energy-loss measurements provided by the SVT and the DCH. Photons and electrons are detected in a CsI(Tl) electromagnetic calorimeter (EMC). Penetrating muons are identified by resistive-plate chambers in the instrumented magnetic flux return (IFR).
Charged pion candidates are required to be detected in at least 12 DCH layers and have a transverse momentum greater than $`100\mathrm{MeV}/c`$. Kaons, electrons, and muons are separated from pions based on information from the IFR and DIRC, energy loss in the SVT and DCH $`(dE/dx)`$, or the ratio of the candidate EMC energy deposition to its momentum ($`E/p`$). Photon candidates are identified with clusters in the EMC with total energy $`>30\mathrm{MeV}`$ and a shower shape consistent with that expected from a photon.
The $`B^0J/\psi \pi ^+\pi ^{}K_S^0`$ and $`B^{}J/\psi \pi ^+\pi ^{}K^{}`$ decays are reconstructed in the following way. Electron candidates and bremsstrahlung photons satisfying $`2.95<m(e^+e^{}(\gamma ))<3.14\mathrm{GeV}/c^2`$ are used to form $`J/\psi e^+e^{}`$ candidates. A pair of muon candidates within the mass interval $`3.06<m(\mu ^+\mu ^{})<3.14\mathrm{GeV}/c^2`$ is required for a $`J/\psi \mu ^+\mu ^{}`$ candidate. A mass constraint to the nominal $`J/\psi `$ mass Eidelman et al. (2004) is imposed in the fit of the lepton pairs. We reconstruct $`K_S^0\pi ^+\pi ^{}`$ candidates from pairs of oppositely charged tracks forming a vertex with a $`\chi ^2`$ probability larger than $`0.1\%`$, a flight-length significance $`l/\sigma (l)>3`$ and an invariant mass within $`15\mathrm{MeV}/c^2`$ of the nominal $`K_S^0`$ mass Eidelman et al. (2004). $`X(3872)`$ candidates are formed by combining $`J/\psi `$ candidates with two oppositely charged pion candidates fitted to a common vertex. Finally, we form $`B^0(B^{})`$ candidates by combining $`X(3872)`$ candidates with $`K_S^0(K^{})`$ candidates. To suppress continuum background, we select only events with a ratio of the second to the zeroth Fox-Wolfram moment G.C. Fox (179) less than $`0.5`$.
We use two kinematic variables to identify signal events from $`B`$ decays: the difference between the energy of the $`B`$ candidate and the beam energy, $`\mathrm{\Delta }EE_B^{}\sqrt{s}/2`$, and the energy-substituted mass $`m_{\mathrm{ES}}\sqrt{(s/2+𝐩_i𝐩_B)^2/E_i^2𝐩_B^2}`$. Here $`(E_i,𝐩_i)`$ is the four-vector (in the laboratory frame) and $`\sqrt{s}`$ is the center-of-mass (CM) energy of the $`e^+e^{}`$ system. $`E_B^{}`$ is the energy of the $`B`$ candidate in the CM system and $`𝐩_B`$ the momentum in the laboratory frame. The signature of signal events is $`\mathrm{\Delta }E0`$, and $`m_{\mathrm{ES}}m_B`$ where $`m_B`$ is the mass of the $`B`$-meson Eidelman et al. (2004).
We optimize the signal selection criteria by maximizing the ratio $`n_s^{mc}/(3/2+\sqrt{n_b^{mc}})`$ Punzi (2003) where $`n_s^{mc}`$ ($`n_b^{mc}`$) are the number of reconstructed Monte Carlo (MC) signal (background) events. The optimization was performed by varying the selection criteria on $`\mathrm{\Delta }E`$, $`m_{\mathrm{ES}}`$, the candidate masses of the $`X(3872)`$ and $`K_S^0`$, and the particle identification (PID) requirements of leptons, pions, and charged kaons. The criteria $`|\mathrm{\Delta }E|<15\mathrm{MeV}`$, $`|m_{\mathrm{ES}}m_B|<6\mathrm{MeV}/c^2`$ and $`|m(J/\psi \pi ^+\pi ^{})3872\mathrm{MeV}/c^2|<6\mathrm{MeV}/c^2`$ (signal region) were found to be optimal for selecting signal events. In case of multiple candidates in an event, we select the candidate with the smallest value of $`|\mathrm{\Delta }E|`$. Applying our optimized selection criteria, we compute the $`J/\psi \pi ^+\pi ^{}`$ invariant mass in the range $`3.83.95\mathrm{GeV}/c^2`$ shown in Figs. 1(a) and 1(b) for the $`B^{}`$ and $`B^0`$ mode, respectively. The shaded area shows events in the sideband region $`|m_{\mathrm{ES}}5260|<6\mathrm{MeV}/c^2`$.
We extract the number of signal events with an extended unbinned maximum-likelihood fit to the two-dimensional distribution $`y(m_{\mathrm{ES}},m_X)`$ where $`m_X`$ is the $`J/\psi \pi ^+\pi ^{}`$ invariant mass. The probability density function (PDF) (normalized to the total number of events) is $`𝒫(y)=_tn_t𝒫_t(y)`$ where $`n_t`$ is the number of events of category $`t`$. We consider three different event categories: signal, $`B`$ decays with the same final-state particles as the signal that accumulate near $`m_{\mathrm{ES}}m_B`$ (peaking background), and combinatorial background. The individual PDFs $`𝒫_t`$ are assumed to be uncorrelated in $`m_{\mathrm{ES}}`$ and $`m_X`$ and can therefore be factorized as $`𝒫_t(y)=g_t(m_{\mathrm{ES}})h_t(m_X)`$, where $`g_t`$ and $`h_t`$ represent the $`m_{\mathrm{ES}}`$ and $`m_X`$ probability distributions, respectively. The $`BX(3872)K`$ signal events are modeled by a Gaussian distribution in $`m_{\mathrm{ES}}`$. The resolution function in $`m_X`$ for those events is best described by a Cauchy function Cauchy due to the mass constraint of the $`J/\psi `$ candidate. The PDF for peaking background events is parameterized by a Gaussian distribution in $`m_{\mathrm{ES}}`$ and a linear function in $`m_X`$. We model combinatorial background events by an ARGUS function Albrecht et al. (1990) in $`m_{\mathrm{ES}}`$ and a linear function in $`m_X`$. The fit performance was validated with MC experiments. The mean and width of the $`m_{\mathrm{ES}}`$ Gaussian distribution for signal and peaking background and the width of the $`m_X`$ Cauchy distribution for the $`B^0`$ mode were fixed to values obtained from MC samples. Other parameters are allowed to vary in the fit.
The fit is performed in the region $`5.2<m_{\mathrm{ES}}<5.3\mathrm{GeV}/c^2`$ and $`3.80<m_X<3.95\mathrm{GeV}/c^2`$ without applying the optimized selection criteria on those two variables. The signal region projections of the two-dimensional fit are shown in Fig. 1 for the $`B^{}`$ (a,c) and $`B^0`$ (b,d) modes. We obtain $`61.2\pm 15.3`$ signal events for the $`B^{}`$ mode ($`n_s^{}`$) and $`8.3\pm 4.5`$ signal events for the $`B^0`$ mode ($`n_s^0`$), respectively. In the following we interpret the observed events in the $`B^0`$ mode as the $`X(3872)`$.
The efficiency is determined from MC samples with an $`X(3872)`$ signal of zero width at $`3.872\mathrm{GeV}/c^2`$. The decay model consists of the sequential isotropic decays $`BX(3872)K`$, $`X(3872)J/\psi \rho ^0`$, and $`\rho ^0\pi ^+\pi ^{}`$. Compared to a three-body decay, this gives a more accurate description of the observed $`\pi ^+\pi ^{}`$ invariant mass distribution Aubert et al. (2004a). Efficiencies are corrected for the small differences in PID and tracking efficiencies that are found by comparing data and MC control samples. The final efficiencies are $`(17.4\pm 0.2)\%`$ for the $`B^0`$/$`K_S^0`$ mode and $`(22.2\pm 0.2)\%`$ for the $`B^{}`$/$`K^{}`$ mode.
The branching fraction systematic errors ($`B^{}`$, $`B^0`$ mode in %) include uncertainties in the number of $`B\overline{B}`$ events (1.1, 1.1), secondary branching fractions (5.0, 5.0) Eidelman et al. (2004), efficiency calculation due to limited MC statistics (0.7, 1.9), MC decay model of the $`X(3872)`$ (1.0, 1.6), differences between data and MC (1.8, 8.9), PID (5.0, 5.0), charged particle tracking (6.0, 4.8), and $`K_S^0`$ reconstruction (-, 1.6). The production ratio of $`B^0`$ and $`B^{}`$ mesons in $`\mathrm{{\rm Y}}(4S)`$ decays is $`1.006\pm 0.048`$ Aubert et al. (2004b). The total fractional error obtained by adding the uncertainties in quadrature is $`9.6\%`$ and $`12.8\%`$ for the $`B^{}`$ and $`B^0`$ mode, respectively.
Assuming Gaussian systematic errors with a PDF $`P_{sys}(n)\mathrm{exp}[(nn_S)^2/2\sigma _{sys}^2]`$, the negative log-likelihood (NLL) function including systematic errors is $`L_{sys}=([1/L(n)][1/\mathrm{ln}P_{sys}(n)])^1`$ where $`L(n)=\mathrm{ln}((n)/_{max})`$ is the NLL projection of the parameter estimate $`n`$ of the number of signal events and $`\sigma _{sys}`$ is the systematic error on the number of signal events. The significance including systematic errors obtained from $`\sqrt{2L_{sys}(n=0)}`$ is $`2.5\sigma `$ for the $`B^0`$ mode and $`6.1\sigma `$ for the $`B^{}`$ mode. The statistical significance ($`\sigma _{sys}=0`$) of the signal is $`2.6\sigma `$ and $`7.5\sigma `$, respectively.
Using $`n_S^0`$ and $`n_S^{}`$, the efficiencies, the secondary branching fractions and the number of $`B\overline{B}`$ events, we obtain the branching fractions $`^0(B^0X(3872)K^0,XJ/\psi \pi ^+\pi ^{})=(5.1\pm 2.8\pm 0.7)\times 10^6`$ and $`^{}(B^{}X(3872)K^{},XJ/\psi \pi ^+\pi ^{})=(10.1\pm 2.5\pm 1.0)\times 10^6`$. For the ratio of branching fractions, $`R^0/^{}`$, where most of the systematic errors cancel, we obtain $`R=0.50\pm 0.30\pm 0.05`$. We calculate a $`90\%`$ confidence level (C.L.) likelihood interval Cowan (1998) $`[n_l,n_h]`$ for the number of signal events in the $`B^0`$ mode by solving the equation $`2L_{sys}(n_{l,h})=[\text{erf}^1(0.95)]^2`$. With $`n_l=2.2`$ and $`n_h=16.9`$ the $`90\%`$ C.L. interval on $`^0`$ is $`1.34\times 10^6<^0<10.3\times 10^6`$. Using the same strategy, the confidence interval on the ratio of branching fractions becomes $`0.13<R<1.10`$ at $`90\%`$ C.L.
We measure the mass of the $`X(3872)`$ in both modes in reference to the precisely measured $`\psi (2S)`$ mass Eidelman et al. (2004). We fit the $`J/\psi \pi ^+\pi ^{}`$ invariant mass in the $`\psi (2S)`$ and $`X(3872)`$ region and calculate $`m_X=m_{X,\text{fit}}m_{\psi (2S),\text{fit}}+m_{\psi (2S)}`$. The result for the $`B^0`$ mode is $`(3868.6\pm 1.2\pm 0.2)\mathrm{MeV}/c^2`$ and $`(3871.3\pm 0.6\pm 0.1)\mathrm{MeV}/c^2`$ for the $`B^{}`$ mode, where the first error is the statistical uncertainty on $`m_{X,\text{fit}}`$ and the second is the uncertainty on $`m_{\psi (2S),\text{fit}}`$ and $`m_{\psi (2S)}`$ Eidelman et al. (2004). The mass difference of the $`X(3872)`$ produced in $`B^0`$ and $`B^{}`$ decays is $`\mathrm{\Delta }m=(2.7\pm 1.3\pm 0.2)\mathrm{MeV}/c^2`$. The full width at half maximum of the $`X`$-mass distribution from the fit on data is $`(6.7\pm 2.7)\mathrm{MeV}/c^2`$, which is consistent with the MC-determined value of $`(5.4\pm 0.1)\mathrm{MeV}/c^2`$. From this we calculate the $`90\%`$ C.L. upper limit on the natural width as $`\mathrm{\Gamma }<4.1\mathrm{MeV}/c^2`$.
Recent observations by BABAR Aubert:2005-ISR in ISR events provide evidence for at least one broad resonance in the invariant mass spectrum of $`J/\psi \pi ^+\pi ^{}`$ at $`4.259\mathrm{GeV}/c^2`$ that can be characterized by a single resonance with a full width of $`88\mathrm{MeV}/c^2`$. This structure is referred to as $`Y(4260)`$. We search in $`B^{}`$ decays for states decaying into $`J/\psi \pi ^+\pi ^{}`$ above $`4\mathrm{GeV}/c^2`$ and impose the additional selection criterion $`|m(K^{}\pi ^+\pi ^{})1273\mathrm{MeV}/c^2|>250\mathrm{MeV}/c^2`$, which removes backgrounds from $`K_1(1270)`$ decays. In the resulting mass distribution, Fig. 2, we observe large combinatoric backgrounds and cannot reliably determine the parameters of one or more resonances. We use a two-dimensional PDF identical to the previous model, but fix the central value and width of the signal component to the ISR results Aubert:2005-ISR . The natural width of $`88\mathrm{MeV}/c^2`$ has been enlarged by the detector resolution, which is found to be the same as for the mass region around $`3.87\mathrm{GeV}/c^2`$. The $`m_X`$ projection of the two-dimensional fit is overlaid in Fig. 2 and yields $`128\pm 42`$ signal events. The statistical significance calculated from $`\sqrt{2\mathrm{ln}_0/}`$ is $`3.1\sigma `$ where $``$ and $`_0`$ are the maximum likelihood of the fit and the null hypothesis fit, respectively. Using a phase-space MC simulation of a state at $`4.26\mathrm{GeV}/c^2`$ decaying into $`J/\psi \pi ^+\pi ^{}`$ and assuming the same systematic uncertainties and efficiency corrections as for the $`X(3872)`$, we obtain the $`95\%`$ C.L. upper limit on the branching fraction $`_Y=(B^{}Y(4260)K^{},Y(4260)J/\psi \pi ^+\pi ^{})<2.9\times 10^5`$. The $`90\%`$ C.L. likelihood interval on the branching fraction $`_Y=(2.0\pm 0.7\pm 0.2)\times 10^5`$ is $`1.2\times 10^5<_Y<2.9\times 10^5`$.
In conclusion, our studies of the $`J/\psi \pi ^+\pi ^{}`$ invariant mass below $`4\mathrm{GeV}/c^2`$ yield a signal of $`2.5\sigma `$ and $`6.1\sigma `$ significance in the $`B^0`$ and $`B^{}`$ mode, respectively, with a ratio of branching fractions $`R=0.50\pm 0.30\pm 0.05`$. We observe an excess of events above background in the $`J/\psi \pi ^+\pi ^{}`$ invariant mass between $`4.2`$ and $`4.4\mathrm{GeV}/c^2`$. These events are consistent with the broad structure observed in ISR events Aubert:2005-ISR .
If one narrow state is observed in the mode $`B^{}X(3872)K^{}`$, $`XJ/\psi \pi ^+\pi ^{}`$, the diquark-antidiquark model Maiani et al. (2005) predicts one amplitude (from $`X_d`$ or $`X_u`$) to be dominant in the charged mode and the other amplitude to be dominant in the neutral mode. In this case, the model predicts the relative rates to be equal ($`R=1`$) and the mass difference to be $`(7\pm 2)\mathrm{MeV}/c^2`$. The ratio of branching fractions is consistent with our measurement, $`0.13<R<1.10`$ at $`90\%`$ C.L., and the observed mass difference of $`(2.7\pm 1.3\pm 0.2)\mathrm{MeV}/c^2`$ is both consistent with zero and the model prediction within two standard deviations. In the S-wave molecule model Braaten and Kusunoki (2005), the neutral mode branching fraction is predicted to be at least 10 times smaller ($`R<0.1`$) than the charged mode. However, we obtain a ratio of neutral to charged branching fractions which is slightly more consistent with isospin-conserving decays.
We are grateful for the excellent luminosity and machine conditions provided by our PEP-II colleagues, and for the substantial dedicated effort from the computing organizations that support BABAR. The collaborating institutions wish to thank SLAC for its support and kind hospitality. This work is supported by DOE and NSF (USA), NSERC (Canada), IHEP (China), CEA and CNRS-IN2P3 (France), BMBF and DFG (Germany), INFN (Italy), FOM (The Netherlands), NFR (Norway), MIST (Russia), and PPARC (United Kingdom). Individuals have received support from CONACyT (Mexico), A. P. Sloan Foundation, Research Corporation, and Alexander von Humboldt Foundation. |
warning/0507/cs0507072.html | ar5iv | text | # Reliable Data Storage in Distributed Hash Tables
## 1 Introduction
Grid Computing and Content Distribution applications place demanding scalability and fault tolerance requirements on underlying services. Client-Server based solutions have often encountered problems meeting these requirements .
Peer to peer technology has demonstrated excellent scalability in file sharing applications. It seems likely this potential will be invaluable in solving scalability problems for more legitimate grid service applications.
Much recent research in the area of peer to peer computing has focused around various kinds of distributed hash table (DHT). These algorithms provide scalable fault tolerant key based routing. Key lookups can be routed reliably to the host responsible for them, which return any associated data.
In this paper we explore the potential of Distributed Hash Tables to provide reliable, scalable and consistent storage of mutable data. We show how the choice of replication algorithm can affect the performance, reliability and bandwidth costs of storing data in a Hash Table, a topic that has previously received little attention.
We will first give a brief overview of distributed hash tables, and describe the aims and design space of replication algorithms.
We then describe both an existing replication algorithm and a new replication algorithm based on dynamic replication, but adapted to provide reliability in an unstable environment, investigating the performance and reliability of both.
We then provide an analysis of the impact of various replica placement strategies possible with our dynamic replication algorithm. We show that replica placement can have a significant impact on the reliability and performance of the system.
Finally, we use a simulation to give a quantitative comparison of the performance and bandwidth costs of the algorithms we have described.
## 2 Distributed Hash Tables
Distributed hash tables provide a solution to the lookup problem in distributed systems. Given the name of a data item stored in the system, we can locate the node on which that data item should be stored. Most DHTs aim to find the responsible node with delay logarithmic in the number of nodes in the network.
There are many different DHT implementations, including PAST, Tapestry, CAN, Kademlia and Chord/DHash . A node in a DHT is typically responsible for the data close to it according to some distance function. Each node maintains knowledge about a small proportion of other nodes in the network, and uses these to forward requests for keys it does not own to nodes which are closer to the requested key. Usually, the distance is reduced by a constant fraction at each routing hop, leading to lookup time logarithmic in the number of nodes.
In any large scale distributed system, nodes are likely to be joining and leaving the system constantly. These changes in the set of participating nodes are called *churn*. DHT routing is often tolerant of node churn, but storage is not. When a node fails, the key-value pairs it was responsible for become inaccessible, and must be recovered from elsewhere. This means that to provide reliability a *replication algorithm* must store and maintain backup copies of data. It must do this in a manner scalable both in the number of nodes and the quantity of data stored in the DHT.
In this paper, we will discuss replication in the context of Chord. Briefly, Chord nodes and data items are given IDs between 0 and some maximum $`K`$, which map to positions on a ring. The distance function between IDs is the clockwise distance around the ring between them. A node *owns*, or is responsible for, data that it is the first clockwise successor of. In order to route key requests, each node maintains knowledge of its immediate clockwise neighbors, called its *successors*, and several other nodes at fractional distances around the ring from it, called its *fingers*. Space concerns prohibit giving a full description of Chord, and for a full description readers are encouraged to consult the original paper
## 3 The Aims of Replication
A Replication algorithm aims to achieve some combination of the following design goals.
The replication algorithm must not rely on any single node, and must recover from churn without user intervention.
The replication algorithm should scale to storing large quantities of data on a large number of nodes, $`N`$. So as not to limit the scalability of Chord, per node replication algorithm state and bandwidth usage should be $`O(log(N))`$.
The replication algorithm may reduce the time taken to look up information by placing replicas of data in a manner that allows network locality to be exploited.
Updating data involves enumerating all replicas. Once this is done, a distributed commit protocol can be used to update the data consistently at all locations. For this reason, enumerating all replicas should be as fast as possible.
The replication algorithm may provide caches of more popular items in order that the load is evenly balanced among all the nodes in the network.
If the replication algorithm is to work with mutable data, it should seek to provide clear guarantees about the consistency of replicas.
Different algorithms achieve these aims to varying extents. The choice of replication strategy may depend on which goals are more important to the task being considered.
Other aims may include resilience to malicious nodes, anonymity or privacy. These are important goals, but we consider them orthogonal to the replication problem, and best treated separately.
## 4 Replication Algorithms
A replication algorithms can be characterized by its approach to four key problems:
Node churn will cause replicas to be lost. The replication algorithm must detect and repair these lost replicas without using excessive bandwidth.
In order to update an item, we need to locate all replicas of that item. Ways of doing this include limiting replica placement to a fixed number of nodes, searching for replicas, and periodically deleting old replicas.
The replica placement strategy determines which nodes replicas should be stored on. This can have an impact on both performance and reliability.
The number of replicas we keep of a given key may either be fixed in advance, or allowed to vary according to the key’s popularity. Variable cardinality often provides superior load balancing, but makes addressability more difficult to achieve.
In this paper, we consider replication by storing a complete copy of the data associated with each key on another node. We believe that the algorithms we will describe could, with some adaptation, also be applied to erasure coded fragments of the original data, with possible performance benefits in some circumstances.
We will now give details of how two replication algorithms approach these tasks.
## 5 DHash Replication
This replication algorithm is a combination of the replica placement strategy first proposed in the original Chord paper, and a maintenance algorithm proposed in . The techniques are combined in the DHash storage system.
The placement strategy is simple, replicas of a data item are placed on the $`r`$ successors of the node responsible for that data item and nowhere else.
Newly joined nodes will inherit a keyspace from the node they precede, and are sent the data they become responsible for when the maintenance algorithm next runs.
The DHash maintenance algorithm runs two maintenance protocols in order to prevent the number of replicas from either dropping too low or rising too high:
A node sends a SYNCHRONIZE message containing the key range it is responsible for to its $`r`$ successors. These nodes then synchronize the contents of their databases so that all keys in this range are stored on both the root node and its successors.
Efficient methods for database synchronization, such as Merkle Tree hashing, are discussed in .
A node periodically checks its database of keys to see if it has any keys it is no longer responsible for. To do this, it looks up the owner of each key it stores, and checks the successor list of that owner.
If it is within $`r`$ hops of the node, then it will be within the first $`r`$ items of the successor list.
If its ID is not in this list, the node is no longer responsible for keeping a replica of this item, and the node offers the data item to the current owner. It can then safely delete the key.
If all replicas are to be repaired by a single maintenance call, the local maintenance algorithm must run two passes, the first gathering all replicas in the root nodes key range onto the root node, the second distributing these replicas onto all successors.
### 5.1 Fetch Algorithm
If we adopt this replication algorithm, we can use the following fetch algorithm to retrieve the data associated with a key. This algorithm will also share the load of providing a data item between all the replica holders.
### 5.2 Maintenance and Reliability
In order to keep the system reliable, we must both store replicas and repair them regularly. The more replicas we keep, the less frequently maintenance is required. This means that, for a given level of reliability, there is a trade-off between the bandwidth used by maintenance algorithm runs and the disk space used for storing replicas.
Here, we attempt to give some new insight into this trade-off, and show the level of replication and maintenance necessary to provide a given level of reliability.
For a given data item to be lost, all $`r`$ of the nodes holding replicas of it must fail. If a replica is missing from the system with probability $`p`$, the probability of a given data item being lost is simply $`p^r`$.
For the purposes of providing reliable storage, however, we are concerned not with the probability of a given item being lost, but with the probability of data loss anywhere in the system. For this to happen, a node and its $`r`$ successors must all have failed at some point in the ring.
To determine this probability, we model the chord ring as a sequence of $`N`$ nodes, each of which is missing replicas with probability $`p`$. The probability of data loss in this model is equivalent to that of obtaining a sequence of $`r`$ successful outcomes in $`N`$ Bernoulli trials with probability of success $`p`$. This is known as the *Run Problem*, and the general solution, $`RUN(p,r,N)`$, can be given in terms of a generating function.
$`F(p,r,s)`$ $`=`$ $`{\displaystyle \frac{p^rs^r(1ps)}{1s+(1p)p^rs^{r+1}}}{\displaystyle \underset{i=r}{\overset{\mathrm{}}{}}}c_i^ps^i`$
$`RUN(p,r,N)`$ $`=`$ $`{\displaystyle \underset{i=r}{\overset{N}{}}}c_i^p`$
In order to relate this to maintenance intervals, we need to determine the proportion of replicas missing from the ring. We will do this in terms of the number of maintenance calls in one *half life*, $`S`$. A network’s half life is the minimum of the times taken for $`N/2`$ of the nodes to leave, and the time $`N/2`$ nodes take to join. We will consider only stable state systems, where nodes join and leave the network at the same rate.
After one half-life, half of the nodes are newly joined, and contain no data<sup>1</sup><sup>1</sup>1This is also a simplification. Although most missing data is caused by empty new nodes, a node also stores data on its $`r`$ successors, and so a failure causes one fewer replica to be stored. Taking this into account gives the fraction of missing replicas as $`\frac{2r+1}{4rS}`$. Our approximation is fair for large $`r`$., so we take $`p=\frac{1}{2}`$. We assume churn occurs at a constant rate, so that a fraction $`\frac{1}{S}`$ of the way through a half life, $`p=\frac{1}{2S}`$.
We make the simplifying assumption that data transfer time is a negligible proportion of the half life, so that a maintenance call will instantly return the system to its ideal state, provided no data has been lost completely before it runs. Each half life can then be divided into $`S`$ independent maintenance intervals, during each of which data is lost with probability $`RUN(\frac{1}{2S},r,N)`$. The overall data loss probability is the probability that any of these $`S`$ maintenance intervals loses data, and so is given by:
$`FAIL(N,r,S)=1(1RUN({\displaystyle \frac{1}{2S}},r,N))^S`$
We can use this to determine how often the maintenance algorithm needs to run in order to maintain a given failure probability. Figure 1 shows the minimum maintenance frequency necessary to maintain $`FAIL(500,r,S)=10^6`$, where $`r`$ varies from 4 to 20. We can see that there is a clear trade-off between bandwidth use and storage space. The number of maintenance calls necessary drops rapidly as we increase the number of replicas.
Figure 2 shows how network size affects the level of maintenance necessary. $`N`$ is varied from 50 to 500 for several values of $`r`$. For small $`r`$, the network size has a significant impact on the required maintenance frequency. As we increase $`r`$, network size becomes less important in determining maintenance frequency.
It should be noted that in systems where data size is large or half-life is short, so that data transfer times become a significant fraction of a half life, our assumptions are not valid and more frequent maintenance or a larger numbers of replicas will be necessary.
### 5.3 Maintenance and Fetch Latency
As we allow the number of replicas to drop between maintenance intervals, we increase the likelihood that we will need to contact more than one node to find the data. If we again approximate the probability of a node not having data as $`\frac{1}{2S}`$ then the number of nodes we need to contact, or probe, before data is found will be geometrically distributed, with expected value:
$`E(probes)={\displaystyle \frac{2S}{2S1}}`$
## 6 Dynamic Replication
Replica Enumeration, as proposed in , aims to remove some of the placement and cardinality restrictions imposed by successor replication, whilst preserving addressability and the ability to make consistent updates.
The placement strategy for Replica enumeration is based around an *allocation function*, $`h(m,d)`$. For each document with ID $`d`$, the replicas are placed at replica locations determined by $`h(m,d)`$ where $`m1`$ is the index of that instance. The allocation function is intended to be pseudo-random, so that the replicas are evenly distributed about the address space.
The replication cardinality is variable in a fixed range $`1rR_{MAX}`$, allowing greater replication for items in greater demand. The mechanism used to decide the exact number of replicas is not specified, but could be designed to adapt to both the reliability of the network, and the load on those nodes providing each document.
To provide addressability, the following invariants are maintained:
1. Replicas of an item $`d`$ are only placed at addresses given by $`h(m,d),`$ where $`1mR_{MAX}`$.
2. For any document $`d`$ in the system, there always exists an initial replica at $`h(1,d)`$
3. Any further replica with $`(m>1)`$ can only exist if a replica currently exists for $`m1`$
Various strategies for finding data are possible in this scheme. One that is generally efficient is to do a binary search over the range \[$`1..R_{MAX}`$\] starting from a randomly selected index in that range. If the data is not replicated at a given location, we use invariant three to reduce the search range accordingly.
Dynamic replication can help alleviate the lookup bottleneck that affects Successor replication. Successor replication requires that all lookups for a popular key are directed to that key’s owner. With an appropriate allocation function, dynamic replication can place replicas at evenly spread, well known, locations around the ring. Lookup queries for the owner of a popular item are then distributed more evenly.
### 6.1 Maintenance algorithm
A major difficulty with the dynamic replication algorithm as originally provided is that maintaining these invariants in a system with a high churn rate is very difficult. Node departures and arrivals could cause any of the invariants to be violated.
It can be shown that lookups will proceed correctly as long as at least invariant two holds. However, for replica addressability, invariants one and four must also hold.
We modify the algorithm given in with a maintenance algorithm to allow it to operate correctly and reliably in a system with high churn rates.
We will use the following definitions to refer to the various roles nodes play in holding replicas:
1. The node responsible for $`h(1,d)`$ is the *owner* of $`d`$.
2. The *replica group* for an item $`d`$ are those nodes whose keyspace includes a replica location from the set $`\{h(m,d):1mR_{MAX}\}`$.
3. The *core group* for an item $`d`$ is the set of replica holders for which $`mR_{MIN}`$
4. The *peripheral group* are those replica holders for which $`R_{MIN}<mR_{MAX}`$.
Since we can not rely on any single node being available in an unreliable environment, we must modify our invariants.
1. Replicas of an item $`d`$ can only be retrieved from addresses given by $`h(m,d)`$ where $`\{1mR_{MAX}\}`$.
2. For any document $`d`$ in the system, there exists with high probability a replica in the core group.
3. Any peripheral replica with $`(m>R_{MIN})`$ can only be retrieved for a single maintenance interval if no replica currently exists for $`m1`$.
We can now give three maintenance protocols which maintain these invariants under churn.
The owner of a data item calculates and looks up the nodes in the core group for the data range it is responsible for. For each core replica holder, the owner and the replica holder synchronize databases over the part of owner’s keyspace which is mapped into that replica holder’s keyspace.
Core maintenance must also deal with *allocation collisions*, as described in the next section.
In order to maintain Invariant 3, a node which stores a replica with index $`m>R_{MIN}`$ must check that a replica of that item is also held on the replica predecessor, the owner of the location with index $`m1`$.
For each peripheral replica a node holds, it must obtain a summary of the items with the previous index on the replica predecessor. Bloom filters can be used to reduce the cost of these exchanges.
These summaries can be used to remove orphaned peripheral replicas from the system. Orphaned peripheral replicas should not be used to answer fetch requests, but should still be stored for at least one maintenance interval, as simulation shows maintenance often replaces the missing replica.
Each node calculates the replica group for each item it holds. Any items it is no longer a replica holder for are offered to their owner, then deleted.
We cache lookups made during maintenance to reduce bandwidth costs. Cache validity is checked at regular intervals during maintenance.
This algorithm attempts to restore the system to its ideal state each time it is run. However, between runs, the system is rarely in its ideal state. Thus we must ensure the system operates correctly where the invariants do not hold. Invariant two is sufficient to ensure lookups proceed correctly, though less efficiently, as nodes fail.
In order to update information, we must discover how many peripheral replicas are currently in the system. To be completely certain of consistency, we must offer all updates to all owners of peripheral replica locations.
If a lower probability of a temporary inconsistency is acceptable, we can improve performance by offering updates only until we encounter a certain number of empty peripheral replica locations since, by invariant three, it becomes increasingly unlikely that any further locations are occupied. This could dramatically improve performance if $`R_{MAX}R_{MIN}`$.
### 6.2 Allocation Collisions
It is important that the images of a node’s key-range under $`h(m,d)`$ are owned by different nodes for each $`d`$. In some cases however, the allocation function will map two replicas into the keyspace of the same node. We call this an *allocation collision*. Each Allocation collision reduces the number of nodes in the core replica group, reducing reliability.
The core maintenance algorithm must keep track of which nodes have been allocated replicas from which key ranges. If an allocation collision occurs, the core maintenance algorithm must instead place the collided keyspace must in the peripheral group, effectively extending the core group. This means we must choose $`R_{MAX}`$ so that sufficient peripheral locations are available to recover from all allocation collisions with high probability. To do this, we must understand how nodes are distributed around the ring.
Since we have $`N`$ nodes uniformly distributed throughout a keyspace of size $`K`$, the probability of a node being at a given ID is $`\frac{K}{N}`$. Therefore, the number of keys between nodes is geometrically distributed with $`p=\frac{K}{N}`$. Using standard properties of a geometric distibution, we can find the mean and variance of this distribution.
$`\mu ={\displaystyle \frac{KN}{N}}{\displaystyle \frac{K}{N}}(sinceKN)`$
$`\sigma ^2={\displaystyle \frac{K(K1)}{N^2}}{\displaystyle \frac{K^2}{N^2}}(sinceK1)`$
To recover from allocation collisions, the range of available replica locations should be at least this size. Thus, by the central limit theorem, the space occupied by $`r`$ nodes will be normally distributed with
$`\mu {\displaystyle \frac{rK}{N}}`$
$`\sigma ^2{\displaystyle \frac{rK^2}{N^2}}`$
And so, by standard properties of the normal distribution, 95% of the time, the keyspace between $`r`$ nodes will be less than $`(r+1.645\sqrt{r})\frac{K}{N}`$ keys in length. To allow $`r`$ replicas can be stored in 95% of cases, the range of available replica locations should be at least this size.
### 6.3 Dynamic fetch algorithm
The dynamic fetch algorithm needs to choose which indexes to lookup and in which order. In many cases, algorithm 2 will give good performance.
In situations where load balancing is not critical, shorter fetch times can be attained by searching the core replica locations before trying peripheral ones.
If maintenance is infrequent, eliminating the entire range of peripheral replicas with higher indexes if one peripheral replica is found empty may result in poor performance, and simply removing the replica known to be empty may be preferable.
### 6.4 Recursive data lookup
In order to increase performance when looking up data, we use algorithm 3 to perform recursive gets. This combines the Chord lookup and get messages, which allows any node on the lookup path of a request to return a replica of the requested data, if it holds one. This avoids further lookup hops, reducing fetch latency.
Recursive data lookup can interfere with load balancing, since some replicas are passed queries more often than others. To prevent this, an overloaded node may choose to forward a recursive get request rather than answer using its own replica.
### 6.5 Allocation Functions
The choice of allocation function is critical to maintenance performance. For each item $`d`$ that a node owns, it must lookup and contact every node in the core replica group in order to run the core maintenance protocol.
In order that this is scalable, we must ensure that as many of these lookups as possible can be satisfied with little network communication. This requires that the allocation function maps one nodes data onto a limited number of replica holders.
We suggest that for a given $`m`$, $`h(m,d)`$ is a *translation* in d. This means that the image of one node’s key-range is another continuous linear range of the same size. Since Chord nodes are equally distributed throughout the key space, an image of one node’s key-range is owned by $`O(1)`$ other nodes.
We will now give four allocation functions and explore how they impact reliability and performance. All these functions make use of $`N`$, the number of nodes in the system. This value may either be supplied by the user, or estimated at run time.
#### 6.5.1 Successor Allocation
$`h(m,d)=(d+(m\frac{K}{N}))\text{mod}K`$
Attempting to map replicas onto successors is efficient, as the Chord Protocol maintains a list of each nodes successors on that node, so lookups can often be performed without consulting another node.
Because replica locations with different indices are relatively close together under this allocation function, we expect some allocation collisions, which must result in the creation of peripheral replicas. Using the result from section 6.2, we recommend that $`R_{MAX}R_{MIN}`$ is at least $`1.645\sqrt{R_{MIN}}`$, so that $`R_{MIN}`$ distinct replicas can be stored in 95% of cases. This consideration also applied to the predecessor and block allocation functions.
#### 6.5.2 Predecessor Allocation
$`h(m,d)=(d(m\frac{K}{N}))\text{mod}K`$
Because queries are routed around the ring clockwise towards the node responsible for them, a lookup for one node is frequently routed through one of its predecessors.
Predecessor allocation aims to exploit this fact to reduce lookup latency. When a request for a document is routed through one of the replica holders, the recursive get algorithm allows them to satisfy the request before it ever reaches the node responsible, reducing the fetch latency by one or more hops.
#### 6.5.3 Block Allocation
$`h(m,d)=(d(d\text{mod}{\displaystyle \frac{KR_{MAX}}{N}})`$
$`+(d\text{mod}{\displaystyle \frac{K}{N}})`$
$`+(m{\displaystyle \frac{K}{N}}))\text{mod}K`$
This allocation function attempts to make replica groups overlap entirely with core replica groups of other nearby keyranges. As will be seen in the next section, this policy provides a lower probability of data loss than other placement functions.
It also provides most of the benefits of both successor and predecessor replication, since most nodes will have replicas placed on both successors and predecessors.
This allocation function is discontinuous in $`d`$, and the maintenance algorithm must be able to deal with this when mapping ranges of keys onto other nodes.
#### 6.5.4 Finger Allocation
$`\delta =log_2({\displaystyle \frac{K}{N}})`$
$`h(m,d)=(d+2^{(m+\delta )})\text{mod}K`$
This allocation function again takes advantage of the information already maintained by the chord algorithm. Chord maintains routing information about nodes at fractional distances around the ring, called fingers. Placing replicas on these finger nodes reduces the number of lookups that must be made, and distributes replicas evenly around the ring.
Because of the distance between replica locations, allocation collisions are rare, and $`R_{MAX}`$ may be set more conservatively than for other allocation functions.
### 6.6 Allocation functions and Reliability
The allocation function chosen has a significant impact on reliability. Block allocation, in which core replica groups for one data range overlap entirely with core groups for nearby data ranges, produces only a very few core replica groups. Finger allocation produces core replica groups which overlap very little with other replica groups, producing a large number of distinct groups.
This large number of distinct groups leads to a higher probability that any one of them will fail. We produced a simple model of a 500 node network, in which 250 nodes are marked as failed. We produced $`10^5`$ sample networks with this model, and used them to estimate the probability of any data loss occurring in the network with varying numbers of replicas. Figure 3 shows the probability of data loss for finger, block and successor allocation functions. Block allocation is able to achieve a more reliable system with a smaller number of replicas.
We also simulated random replica placement, where replicas were placed entirely randomly using a pseudo random number generator. The results for random placement are almost identical to those for finger placement.
In figure 4, we compare the quantity of data lost, given that a failure occurs. More data is lost with block allocation functions when a rare failure does occur than other allocation functions. In many applications however, any quantity of data loss would be considered catastrophic.
## 7 Simulation
We now attempt order to quantitatively compare the performance and bandwidth usage of these replication algorithms. Due to the difficulty of managing large numbers of physical nodes, we chose to test the algorithms through simulation rather than through deployment.
Our simulator is based around the SimPy discrete-event simulation framework, which uses generator functions rather than full threads to achieve scalability. The simulator implements a message level model of a Chord network running each of the replication algorithms described.
### 7.1 Simulation Parameters
In our simulation, we chose parameters that might resemble a data center built from cheap commodity components. While the simulator is capable of running thousands of nodes, we have limited it to two hundred in most scenarios. This was in order to keep runtime reasonable for the large number of scenarios and algorithms we wish to test, and the need to repeat simulations to estimate errors.
We simulate a steady state system in which a node fails every 24 hours, shortly after which a new node replaces it. Latency between nodes is assumed to be uniform, and bandwidth is assumed to be unlimited - messages always take the same time to deliver regardless of size.
We chose fixed parameters for the chord algorithm in all simulations, with a successor list length of 10 and finger table size of 12. Chord repair is carried out at thirty minute intervals.
We also configure the GET algorithm to search the core replica group before trying the peripheral replicas. Local and Core Maintenance algorithms run two passes at each maintenance interval.
Failures are detected by timeouts, which are set to 3 hops for round trip communications. Recursive lookup timeouts are based on network size, and are set to 15 hops in a 200 node network. A shorter timeout could have been chosen, leading to shorter average lookup times, as failed lookups are detected more quickly. Short timeouts also increase bandwidth usage, however, as long-running lookups are reissued before they complete.
The system is simulated for one complete half life, during which 50,000 sample data fetches are made for randomly chosen data items, and fetch times are logged. Bandwidth usage is also logged by type, allowing separation of maintenance messages from chord repair messages. Each simulation is repeated 4 times to obtain a good estimate of the average latencies and bandwidth usage.
## 8 Simulation Results
### 8.1 Fetch Latency
The fetch latency each maintenance algorithm can achieve depends on the network size, the frequency with which it is run, and the number of replicas in the system.
Figure 6 shows how fetch times scale with maintenance frequency. The Finger and DHash algorithms both scale as predicted in section 5.3.
Maintenance frequency has less effect on Successor, Predecessor and Block fetch times. The proximity of different replica indexes means that other replica holders often preemptively return data, so that it is less important that specified replica is present.
The predecessor algorithm achieves the shortest fetch times. With predecessor allocation, queries for core replicas are more often routed through peripheral replicas, which return the data preemptively. DHash fetches are never returned preemptively, and are routed through the requesting node, so that they take at least one hop longer than dynamic allocation.
Figure 5 shows how fetch times scale with network size. All algorithms show logarithmic behavior, though finger allocation compares increasingly badly to other dynamic algorithms since as networks size grows, the probability of a preemptive return drops.
Increasing the number of replicas reduces the lookup times slightly, whereas increasing the number of distinct data items in the system has no impact on lookup times.
### 8.2 Maintenance Bandwidth Costs
Maintenance bandwidth can be divided into two separate costs: The cost of identifying which data should be stored where, and the cost of moving the data to that location. We refer to the former as the maintenance *overhead*.
In figure 7 we can see how overhead varies with maintenance frequency. DHash maintenance has lower overhead than the dynamic algorithms. Dynamic maintenance typically involves $`O(r)`$ lookups, where DHash requires $`O(1)`$. The dynamic algorithms all have similar overhead, which increases linearly with maintenance frequency. Notably, all algorithms overhead bandwidth is so small as to be negligible in most network environments.
Data movement bandwidth is likely to be the bottleneck in distributed storage systems. Figure 8 shows all dynamic algorithms move very similar quantities of data. At high maintenance levels, significantly more data is moved with DHash than the dynamic algorithms, since a single node joining produces changes in the membership of $`r`$ nearby replica groups, with one node leaving each group. This causes the node expelled from each of these replica groups to send any replicas it no longer owns to their owner.
Figure 9 shows the linear relationship between the number of key value pairs and maintenance overhead. Again, maintenance overhead is low for all algorithms, so that it should be feasible to store a very large numbers of items per node without maintenance bandwidth becoming a bottleneck.
Per node maintenance overhead bandwidth and total data movement remain constant as we vary the number of nodes in the system, for all algorithms. This means all algorithms will scale well to very large numbers of nodes.
### 8.3 Fault Tolerance
We have so far investigated the performance of these data replication algorithms under a steady state of churn, in which new nodes join at the same rate as other nodes fail. The DHT can also recover from far higher failure rates, although there is a substantial performance impact.
We simulate a simultaneous failure of a varying proportion of the nodes in the chord network, and then launch 50,000 data fetches immediately afterward. Unlimited retries are allowed, and the average fetch time, including retries is shown in Figure 10.
The DHash algorithm is particularly affected in this scenario, due to its reliance on reaching a single node. The dynamic algorithms are more resilient to faulty routing, as they may select multiple indexes to look up, and because preemptive returns are possible even when the specific node requested is unreachable.
## 9 Related Work
Metadata based algorithms, such the Version ID system used by OceanStore are another possible replica management option. These algorithms do not have placement restrictions, but instead use a metadata item to locate replicas. These solutions may perform well in many scenarios, but require an underlying reliable storage system for metadata. The algorithms we have described could be used to provide such a storage system.
Soft-state storage is another common replica management system. With soft-state storage, no attempt is made to maintain replicas. Instead, data expires after some timeout and the system relies on data periodically being refreshed and reinserted by some external system. Though this can be useful for frequently refreshed data, failure of the external storage system can cause unwanted data loss.
## 10 Conclusions
We have used a combination of analysis and simulation to assess the ability of various replication algorithms to meet the goals we set out in section 3.
We can see that dynamic replication can achieve faster lookups, greater reliability and may require less replica movement than the DHash algorithm, with only a slightly higher maintenance overhead. We have also shown how the allocation function choice can have a dramatic impact on performance. Of the allocation functions we considered, block allocation provides the best reliability and represents a good compromise for most systems, though predecessor placement might be preferable if performance is critical.
Possible drawbacks of dynamic replication are its slightly higher maintenance bandwidth usage and its reliance on an even distribution of node IDs, which may make it unsuitable for small Chord Rings.
System size scalability is good for all maintenance algorithms. Lookup times scale with the logarithm of network size and total system bandwidth consumption scales linearly with the number of nodes.
On an internet wide scale, bandwidth and uptime are likely to be more limited than in the data center scenarios we have considered. In such a system, the number of nodes is likely to vary throughout the day rather than remain constant. Although our work may provide insight into performance in such scenarios, further work needs to be done to assess the reliability of a system which incorporates user desktop systems. |
warning/0507/physics0507145.html | ar5iv | text | # On Statistical Significance of Signal1footnote 11footnote 1Published in High Ener. Phys. Nucl. Phys. 30 (2006) 331.
## I Introduction
The statistical significance of a signal in an experiment of particle physics is to quantify the degree of confidence that the observation in the experiment either confirm or disprove a null hypothesis $`H_0`$, in favor of an alternative hypothesis $`H_1`$. Usually the $`H_0`$ stands for a known or background processes, while the alternative hypothesis $`H_1`$ stands for a new or a signal process plus background processes with respective production cross section. This concept is very useful for usual measurements that one can have an intuitive estimation, to what extent one can believe the observed phenomena are due to backgrounds or a signal. It becomes crucial for the measurements which claim a new discovery or a new signal. As a convention in particle physics experiment, the ”$`5\sigma `$” standard, namely the statistical significance $`S5`$ is required to define the sensitivity for discovery; while in the cases $`S3`$ ($`S2`$), one may claim that the observed signal has strong (weak) evidence.
However, as pointed out in Ref. Sinervo , the concept of the statistical significance has not been employed consistently in the most important discoveries made over the last quarter century. Also, the definitions of the statistical significance in different measurements differ from each other. Listed below are various definitions for the statistical significance in counting experiment (see, for example, refs. Bity1 Bity2 Narsky ):
$$S_1=(nb)/\sqrt{b},$$
(1)
$$S_2=(nb)/\sqrt{n},$$
(2)
$$S_{12}=\sqrt{n}\sqrt{b},$$
(3)
$$S_{B1}=S_1k(\alpha )\sqrt{n/b},$$
(4)
$$S_{B12}=2S_{12}k(\alpha ),$$
(5)
$$_{\mathrm{}}^{S_N}N(0,1)𝑑x=\underset{i=0}{\overset{n1}{}}e^b\frac{b^i}{i!},$$
(6)
where $`n`$ is the total number of the observed events, which is the Poisson variable with the expectation $`s+b`$, $`s`$ is the expected number of signal events to be searched, while $`b`$ is the known expected number of Poisson distributed background events. All numbers are counted in the ”signal region” where the searched signal events are supposed to appear. In equation (4) and (5), the $`k(\alpha )`$ is a factor related to $`\alpha `$ that the corresponding statistical significance assumes $`1\alpha `$ acceptance for positive decision about signal observation, and $`k(0.5)=0,k(0.25)=0.66,k(0.1)=1.28,k(0.05)=1.64`$, etc Bity2 . In equation (6), $`N(0,1)`$ is a notation for the normal function with the expectation and variance equal to 0 and 1, respectively. On the other hand, the measurements in particle physics often examine statistical variables that are continuous in nature. Actually, to identify a sample of events enriched in the signal process, it is often important to take into account the entire distribution of a given variable for a set of events , rather than just to count the events within a given signal region of values. In this situation, I. Nasky Narsky gives a definition of the statistical significance via likelihood function
$$S_L=\sqrt{2\mathrm{ln}L(b)/L(s+b)}$$
(7)
under the assumption that $`2\mathrm{ln}L(b)/L(s+b)`$ distributes as $`\chi ^2`$ function with degree of freedom of 1.
Upon above situation, it is clear that we desire to have a self-consistent definition for statistical significance, which can avoid the danger that the same $`S`$ value in different measurements may imply virtually different statistical significance, and can be suitable for both counting experiment and continuous test statistics. In this letter we propose a definition of the statistical significance, which could be more close to the desired property stated above.
## II Definition of the statistical significance
The $`p`$value is defined to quantify the level of agreement between the experimental data and a hypothesis Ref. Sinervo ; PDG . Assume an experiment makes a measurement for test statistic $`t`$ being equal to $`t_{obs}`$, and $`t`$ has a probability density function $`g(t|H_0)`$ if a null hypothesis $`H_0`$ is true. We futher assume that large $`t`$ values correspond to poor agreement between the data and the null hypothesis $`H_0`$, then the $`p`$value of an experiment would be
$$p(t_{obs})=P(t>t_{obs}|H_0)=_{t_{obs}}^{\mathrm{}}g(t|H_0)𝑑t.$$
(8)
A very small $`p`$value tends to reject the null hypothesis $`H_0`$.
Since the $`p`$value of an experiment provides a measure of the consistency between the $`H_0`$ hypothesis and the measurement, our definition for statistical significance $`S`$ relates with the $`p`$value in the form of
$$_S^SN(0,1)𝑑x=1p(t_{obs})$$
(9)
under the assumption that the null hypothesis $`H_0`$ represents that the observed events can be described merely by background processes. Because a small $`p`$value means a small probability of $`H_0`$ being true, corresponds to a large probability of $`H_1`$ being true, one would get a large signal significance $`S`$ for a small $`p`$value, and vice versa. The left side of equation (9) represents the integral probability of the normal distribution in the region within $`\pm S`$ standard deviation ($`\pm S\sigma `$), therefore, this definition conforms itself to the meaning of that the statistical significance should have. In such a definition, some correlated $`S`$ and $`p`$values are listed in Table 1.
## III Statistical significance in counting experiment
A group of particle physics experiment involves the search for new phenomena or signal by observing a unique class of events that can not be described by background processes. One can address this problem to that of a ”counting experiment”, where one identifies a class of events using well-defined criteria, counts up the number of observed events, and estimates the average rate of events contributed by various backgrounds in the signal region, where the signal events (if exist) will be clustered. Assume in an experiment, the number of signal events in the signal region is a Poisson variable with the expectation $`s`$, while the number of events from backgrounds is a Poisson variable with a known expectation $`b`$ without error, then the observed number of events distributes as the Poisson variable with the expectation $`s+b`$. If the experiment observed $`n_{obs}`$ events in the signal region, then the $`p`$value is
$`p(n_{obs})`$ $`=`$ $`P(n>n_{obs}|H_0)={\displaystyle \underset{n=n_{obs}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{b^n}{n!}}e^b`$
$`=`$ $`1{\displaystyle \underset{n=0}{\overset{n_{obs}1}{}}}{\displaystyle \frac{b^n}{n!}}e^b.`$
Substituting this relation to equation (9), one immediately has
$$_S^SN(0,1)𝑑x=\underset{n=0}{\overset{n_{obs}1}{}}\frac{b^n}{n!}e^b.$$
(11)
Then, the signal significance $`S`$ can be easily determined. Comparing this equation with equation (6) given by Ref. Narsky , we notice the lower limit of the integral is different.
## IV Statistical significance in continuous test statistics
The general problem in this situation can be addressed as follows. Suppose we identify a class of events using well-defined criteria, which are characterized by a set of $`N`$ observations $`X_1,X_2,,X_N`$ for a random variable $`X`$. In addition, one has a hypothesis to test that predicts the probability density function of $`X`$, say $`f(X|\stackrel{}{\theta })`$, where $`\stackrel{}{\theta }=(\theta _1,\theta _2,\mathrm{},\theta _k)`$ is a set of parameters which need to be estimated from the data. Then the problem is to define a statistic that gives a measure of the consistency between the distribution of data and the distribution given by the hypothesis.
To be concrete, we consider the random variable $`X`$ is, say, an invariant mass, and the $`N`$ observations $`X_1,X_2,\mathrm{},X_N`$ give an experimental distribution of $`X`$. Assuming parameters $`\stackrel{}{\theta }=(\theta _1,\theta _2,\mathrm{},\theta _k)(\stackrel{}{\theta _s};\stackrel{}{\theta _b)}`$, where $`\stackrel{}{\theta _s}`$ and $`\stackrel{}{\theta _b}`$ represent the parameters related to signal (say, a resonance) and backgrounds contribution, respectively. We assume the null hypothesis $`H_0`$ stands for that the experimental distribution of $`X`$ can be described merely by the background processes, while the alternative hypothesis $`H_1`$ stands for that the experimental distribution of $`X`$ should be described by the backgrounds plus signal; namely, the null hypothesis $`H_0`$ specifies fixed value(s) for a subset of parameters $`\stackrel{}{\theta _s}`$ (the number of fixed parameter(s) is denoted as r), while the alternative hypothesis $`H_1`$ leaves the r parameter(s) free to take any value(s) other than those specified in $`H_0`$. Therefore, the parameters $`\stackrel{}{\theta }`$ are restricted to lie in a subspace $`\omega `$ of its total space $`\mathrm{\Omega }`$. On the basis of a data sample of size $`N`$ from $`f(X|\stackrel{}{\theta })`$ we want to test the hypothesis $`H_0:\stackrel{}{\theta }`$ belongs to $`\omega `$. Given the observations $`X_1,X_2,,X_N`$, the likelihood function is $`L=_{i=1}^Nf(X_i|\stackrel{}{\theta })`$. The maximum of this function over the total space $`\mathrm{\Omega }`$ is denoted by $`L(\widehat{\mathrm{\Omega }})`$; while within the subspace $`\omega `$ the maximum of the likelihood function is denoted by $`L(\widehat{\omega })`$, then we define the likelihood-ratio $`\lambda L(\widehat{\omega })/L(\widehat{\mathrm{\Omega }})`$. It can be shown that for $`H_0`$ true, the statistic
$$t2\mathrm{ln}\lambda 2(\mathrm{ln}L_{max}(s+b)\mathrm{ln}L_{max}(b))$$
(12)
is distributed as $`\chi ^2(r)`$ when the sample size $`N`$ is large Eadie . In equation (12) we use $`L_{max}(s+b)`$ and $`L_{max}(b)`$ denoting $`L(\widehat{\mathrm{\Omega }})`$ and $`L(\widehat{\omega })`$, respectively. If $`\lambda `$ turns out to be in the neighborhood of 1, the null hypothesis $`H_0`$ is such that it renders $`L(\widehat{\omega })`$ close to the maximum $`L(\widehat{\mathrm{\Omega }})`$, and hence $`H_0`$ will have a large probability of being true. On the other hand, a small value of $`\lambda `$ will indicates that $`H_0`$ is unlikely. Therefore, the critical region of $`\lambda `$ is in the neighborhood of 0, corresponding to large value of statistic $`t`$. If the measured value of $`t`$ in an experiment is $`t_{obs}`$, from equation (8) we have $`p`$value
$$p(t_{obs})=_{t_{obs}}^{\mathrm{}}\chi ^2(t;r)𝑑t.$$
(13)
Therefore, in terms of equation (9), one can calculate the signal significance according to following expression:
$$_S^SN(0,1)𝑑x=1p(t_{obs})=_0^{t_{obs}}\chi ^2(t;r)𝑑t.$$
(14)
For the case of $`r=1`$, we have
$`{\displaystyle _S^S}N(0,1)𝑑x`$ $`=`$ $`{\displaystyle _0^{t_{obs}}}\chi ^2(t;1)𝑑t`$
$`=`$ $`2{\displaystyle _0^{\sqrt{t_{obs}}}}N(0,1)𝑑x,`$
and immediately obtain
$`S`$ $`=`$ $`\sqrt{t_{obs}}`$
$`=`$ $`[2(\mathrm{ln}L_{max}(s+b)\mathrm{ln}L_{max}(b))]^{1/2},`$
which is identical to the equation (7) given by Ref. Narsky .
## V Discussion and Summary
In section 2, the $`p`$value defined by equation (8) is based on the assumption that large $`t`$ values correspond to poor agreement between the null hypothesis $`H_0`$ and the observed data, namely, the critical region of statistic $`t`$ for $`H_0`$ lies on the upper side of its distribution. If the situation is such that the critical region of statistic $`t`$ lies on the lower side of its distribution, then equation (8) should be replaced by
$$p(t_{obs})=P(t<t_{obs}|H_0)=_{\mathrm{}}^{t_{obs}}g(t|H_0)𝑑t,$$
(16)
and the definition of statistical significance $`S`$ expressed by equation (9) is still applicable. For the case that the critical region of statistic $`t`$ for $`H_0`$ lies on both lower and upper tails of its distribution, and one determined from an experiment the observed $`t`$ values in both sides: $`t_{obs}^U`$ and $`t_{obs}^L`$, then equation (8) should be replaced by
$`p(t_{obs})`$ $`=`$ $`P(t<t_{obs}^L|H_0)+P(t>t_{obs}^U|H_0)`$
$`=`$ $`{\displaystyle _{\mathrm{}}^{t_{obs}^L}}g(t|H_0)𝑑t+{\displaystyle _{t_{obs}^U}^{\mathrm{}}}g(t|H_0)𝑑t.`$
In summary, we proposed a definition for the statistical significance by establishing a correlation between the normal distribution integral probability and the $`p`$value observed in an experiment, which is suitable for both counting experiment and continuous test statistics. The explicit expressions to calculate the statistical significance for counting experiment and continuous test statistics in terms of the Poisson probability and likelihood-ratio are given. |
warning/0507/cond-mat0507052.html | ar5iv | text | # Tight-Binding model for semiconductor nanostructures
## I Introduction
Semiconductor quantum dots Woggon (1997); Michler (2000) (QDs) are of particular interest, both concerning basic research and possible applications. QDs are considered to be zero dimensional objects, i.e. systems confined in all three directions of space with a typical size of the magnitude of several nanometers. Therefore, these systems are realizations of “artificial atoms” whose form and size can be manipulated. Concerning basic research these nanostructures (QDs) are interesting, as the methods of quantum theory can be applied to systems on new scales and with new symmetries in between that of atoms or molecules and of macroscopic crystals. On the other side light emission and absorption just from the localized states in such devices may be important for optoelectronic applications Passaseo et al. (2001); Grundmann (2002), quantum cryptography Michler et al. (2000) and quantum computing Li et al. (2003).
Semiconductor QDs can be realized by means of metallic gates providing external (electrostatic) confinement potentials van der Wiel et al. (2003), by means of selforganized clustering of certain atoms in the Stranski-Krastanow (SK) growth mode Merz et al. (1998); Yang et al. (2001); Zhang et al. (2001) or chemically by stopping the crystallographic growth using suitable surfactant materials Kim et al. (2001); Guzelian et al. (1996); Manna et al. (2000); Murray et al. (2001). Here we deal only with the latter two types of QDs. The QDs created in the SK growth mode emerge self-assembled or self-organized in the epitaxial growth process because of the preferential deposition of material in regions of intrinsic strain or along certain crystallographic directions. In epitaxial growth of a semiconductor material A on top of a semiconductor material B only one or a few monolayers of A material may be deposited homogeneously as a quasi two dimensional (2d) A-layer on top of the B-surface forming the so called wetting layer (WL). Under certain conditions and for certain materials further deposited A-atoms will not form a further homogeneous layer but they will cluster and form islands of A-material because this may lower the elastic energy due to the lattice mismatch of the A- and B-material. If one then stops the growth process, one has free A-QDs on top of an A-WL on the B-material. If one continues the epitaxial growth process with B-material, one obtains embedded quantum dots (EQDs), i.e. QDs of A-material on top of an A-WL embedded within B-material.
The chemically realized QDs emerge by means of colloidal chemical synthesis Kim et al. (2001); Guzelian et al. (1996). Thereby the crystal growth of semiconductor material in the surrounding of soap-like films called surfactants is stopped when the surface is covered by a monolayer of surfactant material. Thus one obtains tiny crystallites of the nanometer size in all three directions of space, why these QDs are also called “nanocrystals” (NCs). The size and the shape of the grown NCs can be controlled by external parameters like growth time, temperature, concentration and surfactant material Manna et al. (2000); Murray et al. (2001). Certain physical properties like the band gap (and thus the color) depend crucially on the size of the NCs. Typical diameters for both, EQDs and NCs, are between 3 and 30 nm, i.e. they contain between $`10^3`$ up to $`10^5`$ atoms. Therefore, EQDs and NCs can be considered to be a new, artficial kind of condensed matter in between molecules and solids. For the in SK-modus grown EQDs lens-shaped dots Wojs et al. (1996), dome shaped and pyramidal dots Fry et al. (2000); Ruvimov et al. (1995); Merz et al. (1998), and also truncated cones Shumway et al. (2001) have been found and considered.
Of course the fundamental task is the calculation of the electronic properties of EQDs and NCs. But here one encounters the difficulty that these systems are much larger than conventional molecules and that the fundamental symmetry of solid state physics, namely translational invariance, is not fulfilled. Therefore, neither the standard methods of theoretical chemistry nor the ones of solid state theory can immediately be applied to systems with up to $`10^5`$ atoms. Conventional ab-initio methods of solid state theory based on density functional theory (DFT) and local density approximation (LDA) would require supercell calculations. But the size of a supercell must be larger than the EQD or NC, and such large supercells are still beyond the possibility of present day computational equipment. Therefore, only systems with up to a few hundred atoms can be investigated in the framework of the standard ab-initio DFT methods Puzder et al. (2004); Sarkar and Springborg (2003); Deglmann et al. (2002). Simple model studies based on the effective mass approximation Wojs et al. (1996); Grundmann et al. (1995) or a multi-band $`\stackrel{}{k}\stackrel{}{p}`$-model Stier et al. (1998); Fonoberov and Baladin (2003); Pryor (1998) describe the QD by a confinement potential caused by the band offsets, for instance; they give qualitative insights into the formation of bound (hole and electron) states, but they are too crude for quantitative, material specific results or predictions. More suitable for a microscopic description are empirical pseudopotential methods Wang et al. (1999, 2000); Williamson et al. (2000); Kim et al. (1998) as well as empirical tight-binding models Saito and Arakawa (2002); Lee et al. (2001, 2002); Allan et al. (2000); Niquet et al. (2001); Santoprete et al. (2003); Hill and Whaley (1994); Leung et al. (1998); Leung and Whaley (1999); Pokrant and Whaley (1999); Schrier and Whaley (2003); Lee et al. (2004). The empirical pseudopotential methods allow for a detailed variation of the wave functions on the atomic scale. This is certainly the most accurate description from a microscopic, atomistic viewpoint, but it requires a large set of basis states. Within a TB-model some kind of coarse graining is made and one studies spatial variations only on inter-atomic scales and no longer within one unit cell. The advantage is that usually a small basis set is sufficient, which allows for the possibility to study larger systems. Furthermore the TB-model provides a simple physical picture in terms of the atomic orbitals and on-site and inter-site matrix elements between these orbitals. A cutoff after a few neighbor shells is usually justified for orbitals localized at the atomic sites.
Semiempirical TB-models have been used already to describe ”nearly” spherical InAs and CdSe NCs for which the dangling bonds at the surfaces are saturated by hydrogen Lee et al. (2001, 2002); Niquet et al. (2001); Allan et al. (2000) or organic ligands Leung et al. (1998); Leung and Whaley (1999); Pokrant and Whaley (1999). Also uncapped Saito et al. (1998) and capped Santoprete et al. (2003) pyramidal InAs QDs were ivestigated by use of an empirical TB-model. In the latter work an $`sp^3s^{}`$-basis was used leading to a $`10N\times 10N`$ Hamiltonian matrix, where $`N`$ is the number of atoms, with 33 independent parameters. In the present paper we apply a similar TB-model to II-VI nanostructures, namely CdSe EQDs embedded within ZnSe and spherical CdSe NCs. We show that a smaller TB-basis is sufficient, namely an $`s_cp_a^3`$-basis, i.e. 4 states per unit cell and spin direction. This requires only 8 independent parameters and, in principle, allows for the investigation of larger nanostructures than were accessible in Ref. Santoprete et al., 2003. Strictly speaking, the $`s_cp_a^3`$-basis-set leads to a smaller matrix-dimension and also to a smaller number of nonzero matrix elements compared to $`sp^3s^{}`$ TB-model. So the $`s_cp_a^3`$ TB-model is numerically less demanding regarding both memory requirements and computational time. For the bulk system the valence p-bands are excellently reproduced and the conduction s-band is well reproduced close to the $`\mathrm{\Gamma }`$-point. Therefore, we expect that also for the QDs all the hole states and at least the lowest lying electron states (close to the gap) are well reproduced. We investigate, in particular, the influence of strain effects on the electronic structure. To examine the accuracy of our model we compare the results to other microscopic and macroscopic models. Furthermore TB-results obtained for CdSe-NCs are compared to experimental results, and very good agreement, for instance for the dependence of the energy gap on the NC-diameter, is obtained. This demonstrates that our TB-model with a reduced basis set is reliable and sufficient for the reproduction of the most essential electronic properties of the nanostructures.
This work is organized as follows. In Sec. II our TB-model is presented. The formalism how to obtain the TB-parameters and how to apply them to the description of EQDs and NCs is described. In Sec. III the inclusion of strain effects in our model is introduced. Results for the pyramidal CdSe EQDs are presented. For the spherical CdSe NCs the results and the comparison with the experimental data are presented in Sec. IV. Section V contains a summary and a conclusion.
## II Theory
### II.1 TB-Model for bulk materials
In this work we use a TB-Model with 8 basis states per unit cell. Such a model has been succesfully used for the investigation of optical properties in ZnSe-quantum wells Dierks and Czycholl (1998). For the description of the bulk semiconductor compounds CdSe and ZnSe we choose an $`s_cp_a^3`$ basis set. That implies that the set of basis states $`|\nu ,\alpha ,\sigma ,\stackrel{}{R}`$ is given by four orbitals $`\alpha =s,p_x,p_y,p_z`$ with spin $`\sigma =\pm \frac{1}{2}`$. One $`s`$-orbital at the cation ($`\nu =c`$) and three $`p`$-orbitals at the anion ($`\nu =a`$) site in each unit cell $`\stackrel{}{R}`$ are chosen. The TB matrix elements are given by
$$E_{\alpha ,\alpha ^{}}(\stackrel{}{R}^{}\stackrel{}{R})_{^{\nu ,\nu ^{}}}=\nu ^{},\alpha ^{},\sigma ^{}\stackrel{}{R}^{}|H^{\text{bulk}}|\nu ,\alpha ,\sigma ,\stackrel{}{R}.$$
(1)
The coupling of the basis orbitals is limited to nearest and next nearest neighbors. Following Ref. Chadi, 1977, the spin-orbit component of the bulk Hamiltonian $`H^{\text{bulk}}`$ couples only $`p`$ orbitals at the same atom. With the two center approximation of Slater and Koster Slater and Koster (1954) we are left with only 8 independent matrix elements $`E_{\alpha ,\alpha ^{}}(\stackrel{}{R}^{}\stackrel{}{R})_{^{\nu ,\nu ^{}}}`$.
In $`\stackrel{}{k}`$ space, with the basis states $`|\stackrel{}{k},\nu ,\alpha ,\sigma `$, the electronic properties of the pure bulk material are modelled by an $`8\times 8`$ matrix $`𝐇^{\text{bulk}}(\stackrel{}{k})`$ (for each $`\stackrel{}{k}`$ from the first Brillouin zone). This matrix depends on the different TB-parameters $`E_{\alpha ,\alpha ^{}}(\stackrel{}{R}^{}\stackrel{}{R})_{^{\nu ,\nu ^{}}}`$. By analytical diagonalization for special $`\stackrel{}{k}`$ directions, the electronic dispersion $`E_n(\stackrel{}{k})`$ is obtained as a function of the TB-parameters; here $`n`$ is the band index. Equations for the different TB-parameters $`E_{\alpha ,\alpha ^{}}(\stackrel{}{R}^{}\stackrel{}{R})_{^{\nu ,\nu ^{}}}`$ can now be deduced as a function of the Kohn-Luttinger-parameters ($`\gamma _1`$,$`\gamma _2`$,$`\gamma _3`$), the energy gap $`E_{\text{gap}}`$, the effective electron mass $`m_e`$ and the spin-orbit-splitting $`\mathrm{\Delta }_{\text{so}}`$. The zero level of the energy scale is fixed to the valence-band maximum. The used material parameters for CdSe and ZnSe are given in Table 1. The resulting numerical values for the different TB-parameters (obtained by optimizing them so that the resulting TB band-structure reproduces the parameters given in Table 1) are summarized in the Table 2 (with and without taking into account a site-diagonal parameter for the spin-orbit coupling). Within this approach, the characteristic properties of the band structure in the region of the $`\mathrm{\Gamma }`$ point are well reproduced, as can be seen from Fig. 1, which shows the TB-bands of bulk CdSe and ZnSe (using the TB-parameters with spin-orbit coupling). When comparing with band structure results from the literatureCohen and Chelikowski (1989), one sees that the three valence ($`p`$-) bands are excellently reproduced whereas the $`s`$-like conduction band is well reproduced only close to the $`\mathrm{\Gamma }`$-point. This is understandable, because higher (unoccupied) conduction bands are neglected, and can be improved by taking into account more basis states per unit cell. But for a reproduction of the electronic properties in the region near the $`\mathrm{\Gamma }`$-point, which is important for a proper description of the optical properties of the semiconductor materials, the $`s_cp_a^3`$-TB-model is certainly sufficient and satisfactory.
### II.2 TB-Model for Embedded Quantum Dots and Nanocrystals
Having determined suitable TB-parameters for the bulk materials (here CdSe and ZnSe) a EQD or NC can be modelled simply by using the TB-parameters of the bulk materials for those sites (or unit cells) occupied by atoms (or molecules) of this material. Concerning the on-site matrix elements this condition is unambiguous. Concerning the intersite matrix elements one also uses the bulk matrix elements, if the two sites are occupied by the same kind of material, but one has to use suitable averages of the bulk intersite matrix elements for matrix elements over interfaces between different material, i.e. if the two sites (or unit cells) are occupied by different atoms (or molecules). Concerning the surfaces or boundaries of the nanostructure there are different possibilities. One can use fixed boundary conditions, i.e. effectively use zero for the hopping matrix elements from a surface atom to its fictitious nearest neighbors, or (for the embedded QDs) one can use periodic boundary conditions to avoid any surface effects, which artificially arise from the finite cell size used for the EQD-modelling. For the NCs the best thing to do is a realistic, atomistic modelling of the organic ligands covering the NC-surface, as described in Refs.Pokrant and Whaley (1999); Rabani (2003); Puzder et al. (2003). Within the restricted basis set thus selected the ansatz for an electronic eigenstate of the EQD or NC is, of course, a linear combination of the atomic orbitals $`|\nu ,\alpha ,\sigma ,\stackrel{}{R}`$:
$$|\mathrm{\Phi }=\underset{\alpha ,\nu ,\sigma ,\stackrel{}{R}}{}u_{\nu ,\alpha ,\sigma ,\stackrel{}{R}}|\nu ,\alpha ,\sigma ,\stackrel{}{R}.$$
(2)
Here $`\stackrel{}{R}`$ denotes the unit cell, $`\alpha `$ the orbital type, $`\sigma `$ the spin and $`\nu `$ an anion or cation. Then the Schrödinger equation leads to the following finite matrix eigenvalue problem:
$$\underset{\alpha ,\nu ,\sigma ,\stackrel{}{R}}{}\nu ^{},\alpha ^{},\sigma ^{},\stackrel{}{R}^{}|H|\nu ,\alpha ,\sigma ,\stackrel{}{R}u_{\nu ,\alpha ,\sigma ,\stackrel{}{R}}Eu_{\nu ^{},\alpha ^{},\sigma ^{},\stackrel{}{R}^{}}=0,$$
(3)
where $`E`$ is the energy eigenvalue. The shortcut notation $`\nu ^{},\alpha ^{},\sigma ^{},\stackrel{}{R}^{}|H|\nu ,\alpha ,\sigma ,\stackrel{}{R}=H_{l\stackrel{}{R}^{},m\stackrel{}{R}}`$ is used in the following for the matrix elements with $`l=\nu ^{},\alpha ^{},\sigma ^{}`$ and $`m=\nu ,\alpha ,\sigma `$.
The matrix elements for CdSe and ZnSe without strain are denoted by $`H_{l\stackrel{}{R}^{},m\stackrel{}{R}}^0`$. For these matrix elements the TB-parameters $`E_{\alpha ,\alpha ^{}}(\stackrel{}{R}^{}\stackrel{}{R})_{^{\nu ,\nu ^{}}}`$ of the bulk materials, determined in Sec. II.1, are used. For the off-diagonal matrix elements over interfaces and the diagonal matrix elements of the selen atoms at the interface between dot and barrier, which can not unambiguously be referred to belong to the ZnSe or CdSe, respectively, we choose the mean value of the parameters for the two materials. Furthermore, a parameter for the valence-band offset $`\mathrm{\Delta }E_V`$ has to be included in the model. This means that for CdSe in a heterostructure, i.e. surrounded by a barrier ZnSe material, all diagonal matrix elements are shifted just by $`\mathrm{\Delta }E_V`$ compared to the bulk CdSe diagonal matrix elements. In the literature different values for $`\mathrm{\Delta }E_V`$ can be found, they vary in the range of 10 %-30 % of the band gap difference between CdSe and ZnSe Kurtz et al. (2001); Yang et al. (2001); Lakes et al. (1996). We have performed calculations with valence-band offsets of $`\mathrm{\Delta }E_V=0.108`$ eV, $`\mathrm{\Delta }E_V=0.22`$ eV and $`\mathrm{\Delta }E_V=0.324`$ eV, which corresponds to 10 %, 20 % and 30 % of the difference of the band gaps. We find that these different choices for $`\mathrm{\Delta }E_V`$ shift the EQD energy gap $`E_{\text{gap}}^{\text{QD}}`$ by less than 2 %. This shows, that the results are not much affected by the specific choice of the valence-band offset $`\mathrm{\Delta }E_V`$. Therefore, in the following, an intermediate value of $`\mathrm{\Delta }E_V=0.22`$ eV is chosen.
Furthermore, in a heterostructure of two materials with different lattice constants, strain effects have to be included for a realistic description of the electronic states, because the distance between two CdSe unit cells and the bond angles are not the same as the corresponding equilibrium values in bulk CdSe. This means that the intersite TB matrix elements $`H_{l\stackrel{}{R}^{},m\stackrel{}{R}}`$ in the EQD differ from the $`H_{l\stackrel{}{R}^{},m\stackrel{}{R}}^0`$ matrix elements in the bulk material. In general, a relation
$$H_{l\stackrel{}{R}^{},m\stackrel{}{R}}=H_{l\stackrel{}{R}^{},m\stackrel{}{R}}^0f(\stackrel{}{d}_{\stackrel{}{R}^{}\stackrel{}{R}}^{\mathrm{\hspace{0.17em}0}},\stackrel{}{d}_{\stackrel{}{R}^{}\stackrel{}{R}})$$
(4)
has to be expected, where $`\stackrel{}{d}_{\stackrel{}{R}^{}\stackrel{}{R}}^{\mathrm{\hspace{0.17em}0}}`$ and $`\stackrel{}{d}_{\stackrel{}{R}^{}\stackrel{}{R}}`$ are the bond vectors between the atomic positions of the unstrained and strained material, respectively. The function $`f(\stackrel{}{d}^{\mathrm{\hspace{0.17em}0}},\stackrel{}{d})`$ describes, in general, the influence of the bond length and the bond angle on the intersite (hopping) matrix elements. For lack of a microscopic theory for the functional form we use as a simplified model assumption $`f(\stackrel{}{d}_{\stackrel{}{R}^{}\stackrel{}{R}}^0,\stackrel{}{d}_{\stackrel{}{R}^{}\stackrel{}{R}})=\left(d_{\stackrel{}{R}^{}\stackrel{}{R}}^0/d_{\stackrel{}{R}^{}\stackrel{}{R}}\right)^2`$. With this $`d^2`$ ansatz, the interatomic matrix elements $`H_{l\stackrel{}{R}^{},m\stackrel{}{R}}`$, with $`\stackrel{}{R}^{}\stackrel{}{R}`$, are given by
$$H_{l\stackrel{}{R}^{},m\stackrel{}{R}}=H_{l\stackrel{}{R}^{},m\stackrel{}{R}}^0\left(\frac{d_{\stackrel{}{R}^{}\stackrel{}{R}}^0}{d_{\stackrel{}{R}^{}\stackrel{}{R}}}\right)^2.$$
(5)
This corresponds to Harrison’s Froyen and Harrison (1979) $`d^2`$ rule, the validity of which has been demonstrated for II-VI-materials by Sapra et al. Sapra et al. (2002). More sophisticated ways to treat the scaling of the interatomic matrix elements, e.g. by calculating the dependence of energy bands on volume effects and different exponents for different orbitals, can be found in the literature Jancu et al. (1998); Santoprete et al. (2003); Lee et al. (2004). Furthermore the results of Bertho et al. Bertho et al. (1991) for the calculations of hydrostatic and uniaxial deformation potentials in case of ZnSe show that the $`d^2`$ rule should be a reasonable approximation. Our model assumption for the function $`f(\stackrel{}{d}^0,\stackrel{}{d})`$ means that we neglect the influence of bond angle distortion. Though energy shifts due to bond angle distortions have been found for InAs EQDsSantoprete et al. (2003), here the negligence of bond angle distortion can be justified when exclusively taking into account the coupling between s- and p-orbitals at nearest neighbor sites. Piezoelectric fields, which are usually considered to be less important for the zinc blende structures realized in CdSe and ZnSe Fonoberov and Baladin (2003), are also not taken into account in our model.
The problem is now reduced to the diagonalization of a finite but very large matrix. To calculate the eigenvalues of this matrix, in particular the bound electronic states in the QD, the folded spectrum method Wang and Zunger (1994) is applied to the eigenvalue problem of Eq. 3.
## III Results for a pyramidal $`\mathrm{𝐂𝐝𝐒𝐞}`$ Embedded Quantum Dot
### III.1 Geometry and Strain
To model a CdSe QD embedded into a ZnSe barrier material we choose a finite (zinc blende) lattice within a box with fixed boundary conditions. Within this box we consider a CdSe WL of thickness $`1a`$ (lattice constant of the conventional unit cell, i.e. about two anion and two cation layers), and on top of this wetting layer there is a pyramidal QD with base length $`b`$ and height $`h=b/2`$. For the matrix elements corresponding to sites within the WL or the QD we choose the TB-values appropriate for CdSe, for all other sites within the box the ones for ZnSe. Figure 2 shows a schematic picture of this geometry we use to model the EQD. We investigate EQDs with a base lengths $`b`$ of $`6a`$, $`8a`$ and $`10a`$, where $`a=5.668\text{Å}`$ is the lattice constant of the bulk ZnSe material. Cells with the dimensions of $`18a\times 18a\times 15a`$ (38 880 atoms), $`20a\times 20a\times 16a`$ (51 200 atoms) and $`22a\times 22a\times 17a`$ (65 824 atoms) are used for the calculations.
Figure 2 shows the EQD with a base length $`b`$ of $`10a`$. Fixed boundary conditions are applied to avoid a dot-dot coupling in contrast to periodic boundary conditions Santoprete et al. (2003). The total size of the cells is chosen so that the boundary conditions affect the energy gap of the EQD by less than 2 %.
To consider strain effects in our model the knowledge of the strain tensor $`\mathit{ϵ}`$ is necessary. The strain tensor $`\mathit{ϵ}`$ is related to the strain dependend relative atomic positions $`\stackrel{}{d}_{\stackrel{}{R}^{}\stackrel{}{R}}`$ by
$$\stackrel{}{d}_{\stackrel{}{R}^{}\stackrel{}{R}}=(\mathrm{𝟙}+\mathit{ϵ})\stackrel{}{d}_{\stackrel{}{R}^{}\stackrel{}{R}}^{\mathrm{\hspace{0.17em}\hspace{0.17em}0}}.$$
(6)
To appoint the strain tensor outside the EQD, the WL is treated as a quantum film. In the absence of a shear strain ($`ϵ_{i,j}\delta _{i,j}`$) for a coherently grown film, the strain components are given by Chuang (1995)
$`ϵ_{||}=ϵ_{xx}=ϵ_{yy}`$ $`=`$ $`{\displaystyle \frac{a_Sa_D}{a_D}}`$ (7)
$`ϵ_{}=ϵ_{zz}`$ $`=`$ $`{\displaystyle \frac{C_{12}}{C_{11}}}ϵ_{||}.`$ (8)
Here $`a_D`$ is the lattice constant of the unstrained film material and $`a_S`$ denotes the parallel lattice constant of the substrate. In Table 1 the cubic elastic constants $`C_{ij}`$ of the bulk materials are given. The resulting strain profile for a line scan in z-direction outside the dot is shown in Fig. 3 (a). In Ref. Stier et al., 1998 Stier et al. considered a similar strain profile for an InAs/GaAs EQD. The lattice mismatch of approximately 7 % in the InAs/GaAs system is nearly the same as for the CdSe/ZnSe system. So our calculated strain profile shows the same behavior as the profile in Ref. Stier et al., 1998 for a line scan in z-direction outside the EQD.
To obtain the strain profile inside the EQD we use a model strain profile, which shows a similar behavior as the strain profiles which are given in Refs. Pryor et al., 1998; Stier et al., 1998 for a line scan in z-direction through the tip of the pyramid. This model strain profile is displayed in Fig. 3 (b). The shear components, $`ϵ_{xy}`$, $`ϵ_{xz}`$ and $`ϵ_{yz}`$, can be neglected, at least away from the boundaries of the dot Grundmann et al. (1995).
### III.2 Bound single particle states
We have calculated the first five states for electrons and holes for three different EQD sizes. These calculations are done with and without including strain effects. For the evaluations without strain we have chosen the exponent in Eq. (5) to be zero. The energy spectrum obtained from these calculations is shown in Fig. 4 (a) without strain and in Fig. 4 (b) including strain effects. The states are labeled by $`e_1`$ and $`h_1`$ for electron and hole ground states, $`e_2`$ and $`h_2`$ for the first excited states, and so on. All energies are measured relative to the valence-band maximum of ZnSe. Figure 4 also shows the size dependence of the electron and hole energy levels. The energies are compared to the ground state energies for electrons and holes in the $`1a`$ thick CdSe WL ($`\text{WL}_{e_1}`$ and $`\text{WL}_{h_1}`$, respectively), which is calculated separately for a coherently strained quantum film (i.e. the WL without the QD). As expected from a naive particle in a box picture, the binding of electrons and holes becomes stronger in the EQD when the dot size is increased. The quantum confinement causes the number of bound states to decrease when the dot size is reduced. For the EQDs with a base length $`b=8a`$ and $`b=10a`$ the calucated hole states are well above the WL energy ($`\text{WL}_{h_1}`$). This is valid for the strain-unaffected and strained EQD. For the system with $`b=6a`$ we obtain at least four bound hole states in both models. The energy splitting between the different states is only slightly influenced by the strain. Furthermore we see from Fig. 4 that the number of bound electron states is influenced by the strain. For the system with a base length of $`b=10a`$ we get at least three bound-electron states when we take strain effects into account (Fig. 4 (b)). Without strain effects at least 5 bound states are found. So the confinement potential for the electrons is effectively reduced by the strain.
The bound electron states $`e_2`$ and $`e_3`$ are energetically not degenerate even without strain. This arises from the $`C_{2v}`$ symmetry of the system. Already from the geometry of the EQD-system it is clear that there is no $`(001)`$ mirror plane. Furthermore, if one considers a $`(001)`$-plane with sites occupied by Se anions, the nearest neighbor (cation) planes in $`\pm z`$-direction are not equivalent, as in the zinc blende structure the nearest neighbors above the plane are found in $`[111]`$-direction and below the plane in $`[1\overline{1}\overline{1}]`$-direction. So also for crystallographic reasons a $`(001)`$-plane is not a mirror plane. Finally, if one considers the base plane of the EQD (or the WL) to be this anion $`(001)`$-plane, there are different cations, namely Cd above and Zn below this plane. Therefore, the QD-system has reduced $`C_{2v}`$-symmetry. In theories based on continuum models, e.g. effective mass approximations Wojs et al. (1996); Grundmann et al. (1995), the discussed effects cannot be accounted for.
These interfacial effects also affect the one particle wave functions in the system. In Fig. 5 the isosurfaces for the squared electron wavefunctions $`|\mathrm{\Phi }_i(\stackrel{}{r})|^2`$ are displayed with and without strain, respectively. The light and dark isosurface levels are selected as 0.1 and 0.5 of the maximum probability density, respectively. For both calculations, the lowest electron state $`e_1`$ is an $`s`$-like state according to its nodal structure. The next two states $`e_2`$ and $`e_3`$ are $`p`$-like states. These states are oriented along the $`[1\overline{1}0]`$ and the $`[110]`$ direction, respectively. Due to the different atomic structure along these directions we find a $`p`$-state splitting $`\mathrm{\Delta }_{e_2,e_3}^0=E_{e_3}E_{e_2}`$ for the unstrained EQD of about $`0.43`$meV. In conventional $`\stackrel{}{k}\stackrel{}{p}`$ models Grundmann et al. (1995); Pryor (1998) an unstrained, square-based pyramidal EQD is modelled with a $`C_{4v}`$ symmetry. In our microscopic model the resulting degeneracy is lifted and a splitting occurs as a consequence of the reduction of $`C_{4v}`$ symmetry to a $`C_{2v}`$ zinc blende symmetry.
The strain splits the states $`e_2`$ and $`e_3`$ further. Due to the different atomic structure, the strain profile within each plane (perpendicular to the growth z-direction) along the $`[110]`$ and $`[1\overline{1}0]`$ direction is different Pryor et al. (1998). This effect contributes also to the anisotropy. Due to the fact, that the base is larger than the top, there is a gradient in the strain tensor between the top and the bottom of the pyramid. In the EQD, the cation neighbors above each anion are found in $`[111]`$ direction while the cation neighbors below are found in $`[1\overline{1}\overline{1}]`$-direction. Therefore, the cations along the $`[1\overline{1}0]`$ direction are systematically more stressed than the cations along the $`[110]`$ direction. In case of strain we find a $`p`$-state splitting of $`\mathrm{\Delta }_{e_2,e_3}^{\text{strain}}=7.1`$ meV. Compared to the states $`e_2`$ and $`e_3`$ of the unstrained EQD, the two lumps of the light isosurfaces are well separated. The states $`e_2`$ and $`e_3`$ reveal nodal planes along the $`[110]`$ and $`[1\overline{1}0]`$ direction, respectively.
The state $`e_4`$ for the strained dot is resonant with a WL-state, so the wave function is leaking into the WL. Also the wave function of the state $`e_5`$ is localized at the base of the pyramid but clearly shows already a finite probability density inside the WL. The states $`e_4`$ and $`e_5`$ of the unstrained EQD are still mainly localized inside the dot. The classification of the state $`e_4`$ by its nodal structure is difficult. $`e_4`$ is similar to a $`p`$-state which is oriented along the $`[001]`$ direction. The electron state $`e_5`$ is $`d`$-like.
Figure 6 shows the isosurface plots of the squared wavefunction $`|\mathrm{\Phi }_i(\stackrel{}{r})|^2`$ for the lowest five hole states $`h_1`$-$`h_5`$ with and without strain. The light and dark isosurface levels are again selected as 0.1 and 0.5 of the maximum probability density, respectively. Our atomistic calculation shows that the hole states cannot be classified by $`s`$-like ($`h_1`$), $`p`$-like ($`h_2`$ and $`h_3`$) or $`d`$-like ($`h_5`$) shape according to their nodal structures. With and without strain the hole states underly a strong band mixing. So the calculated hole states show no nodal structures. Therefore the assumption of a single heavy-hole valence-band for the description of the bound hole states in a EQD even qualitatively yields incorrect results. In contrast to quantum well systems, the light-hole and heavy-hole bands are strongly mixed in a EQD. This result is in good agreement with other multiband approaches Stier et al. (1998); Fonoberov and Baladin (2003); Pryor (1998); Wang et al. (1999, 2000); Williamson et al. (2000); Kim et al. (1998).
From Fig. 6 we can also estimate the influence of strain on the different hole states. Without strain the states $`h_1`$ and $`h_2`$ are only slightly elongated along the $`[1\overline{1}0]`$ and $`[110]`$ direction, respectively. Due to strain these states are clearly elongated along these directions. The states $`h_3`$-$`h_5`$ are only slightly affected by strain.
Another interesting result is that strain effects shift the electron states to lower energies and the hole states to higher energies as displayed in Fig. 4. Figure 4 also reveals that the WL ground-state for electrons and holes is shifted in a similar way due to strain. We observe that strain decreases the EQD gap $`E_{\text{gap}}^{\text{QD}}=E_{e_1}E_{h_1}`$ by about 1.4%, lowering it from the strain-unaffected value 2.12 eV to the value 2.09 eV. For a biaxial compressive strain in a zinc blende structure, the conduction-band minimum of a bulk material is shifted to higher energies while the energy shift of the valence-band maximum depends on the magnitude of the hydrostatic and shear deformation energies Chuang (1995). So one would expect that the electron states are shifted to higher energies due to the fact that CdSe is compressively strained in the ZnSe-Matrix. This is in contradiction to the behavior we observe here. To investigate the influence of the WL states on the one-particle spectrum we use the same model geometry as shown in Fig. 2 but with a considerably smaller WL thickness of only one monolayer (ML). A 1 ML thick WL was also used before by Santoprete et al. Santoprete et al. (2003), Stier et al. Stier et al. (1998) and Wang et al. Wang et al. (1999) for an InAs/GaAs EQD. Figure 7 shows the comparison of the results for a strain-unaffected and a strained pyramidal CdSe EQD with a 1 ML thick WL and a base length of $`b=10a`$.
On the left-hand side of Fig. 7 the first five electron and hole-state energies for an unstrained EQD are displayed while the right-hand side shows the energies for the strained EQD. For a 1 ML thick WL the lowest electron state is, by strain effects, shifted to higher energies. This is what one would expect for biaxial compression of the bulk material. Furthermore the splitting of the $`p`$-like states $`e_2`$ and $`e_3`$ is larger compared to the results for a $`1a`$ thick WL. The splitting $`\mathrm{\Delta }_{e_2,e_3}^0`$ of the unstrained EQD with a $`1a`$ thick WL is $`\mathrm{\Delta }_{e_2,e_3}^0=0.43`$ meV whereas for the system with a 1 ML thick WL one has $`\mathrm{\Delta }_{e_2,e_3}^0=0.5`$ meV. So the spltting $`\mathrm{\Delta }_{e_2,e_3}^0`$ is increased by about 16 %. With strain-effects, the splitting for the system with 1 ML WL thickness $`\mathrm{\Delta }_{e_2,e_3}^{\text{strain}}=10.9`$ meV is about 54 % larger than the splitting in the system with $`1a`$ WL thickness $`\mathrm{\Delta }_{e_2,e_3}^{\text{strain}}=7.1`$ meV. Also the energy splitting $`\mathrm{\Delta }_{e_1,e_2}`$ between the ground state $`e_1`$ and the first excited state $`e_2`$ is strongly influenced by the WL thickness, namely $`\mathrm{\Delta }_{e_1,e_2}=162.8`$ meV for the unstrained sytem with $`1a`$ WL but $`\mathrm{\Delta }_{e_1,e_2}=204.1`$ meV for a 1 ML WL; with strain effects the splitting $`\mathrm{\Delta }_{e_1,e_2}`$ is increased by about 27 % if the WL thickness is decreased from $`1a`$ to 1 ML. The results are summerized in Table 3. This effect mainly arises from the fact, that the bound states inside the dot are also coupled to the WL-states. For a $`1a`$ WL the wave functions of the bound states show also a probability density inside the WL. For a thinner WL the leaking of the states into the region of the WL is much less pronounced. In this case, the microscopic structure inside the EQD and also the strain-affect are much more important. This explains the larger energy splittings in case of the 1 ML WL compared to the results for a $`1a`$ WL. The hole states are influenced in a similar manner. In case of a 1 ML WL the energy spectrum of the hole states is shifted to higher energies due to the strain-effects. This behavior is similar to the behavior obtained from the calculations for a $`1a`$ WL (Fig. 7). In the 1 ML WL system the energy splittings $`\mathrm{\Delta }_{h_1,h_2}`$ and $`\mathrm{\Delta }_{h_2,h_3}`$ for the first three hole-states are larger than the values we obtain for the system with $`1a`$ WL. These splittings are also summarized in Table 3. The WL thickness also influences the EQD energy gap $`E_{\text{gap}}^{\text{QD}}`$. For a 1 ML WL the electron-states are shifted to higher energies in contrast to the behavior of the hole states (compare Figs. 4 and 7). In case of the 1 ML WL the gap energy $`E_{\text{gap}}^{\text{QD}}`$ is only slightly affected by the strain. We observe here that the strain has opposite effect for electrons and hole states: electron states become shallower, approaching the conduction-band edge, while the hole states become deeper, moving away from the valence-band edge.
The knowledge of the single-particle wave functions makes the examination of many-particle effects in EQDs possible. The single particle wave functions can be used for the calculation of Coulomb- and dipole- matrix elements as input parameters. For example the investigation of multi-exciton emission spectra Baer et al. (2004), carrier capture and relaxation in semiconductor quantum dot lasers Nielsen et al. (2004) or a quantum kinetic description of carrier-phonon interactions Seebeck et al. (2005) is possible.
## IV Results for $`\mathrm{𝐂𝐝𝐒𝐞}`$ Nanocrystals
### IV.1 Geometry and Strain
In this section we investigate the single particle states of CdSe nanocrystals within our TB-model. These nanostructures are chemically synthesized Kim et al. (2001); Guzelian et al. (1996). The nanocrystals are nearly spherical in shape Banin et al. (1999); Guzelian et al. (1996); Alivisatos (1996) and the surface is passivated by organic ligands. Due to the flexible surrounding matrix, these nanostructures are nearly unstrained Alivisatos (1996). The size of these nanostructrues is in between 10 and 40 Å in radius Banin et al. (1998, 1999); Alperson et al. (1999); Kim et al. (2001).
We model such a chemically synthesized NC as an unstrained, spherical crystallite with perfect surface passivation. The zincblende structure is assumed for the CdSe nanocrystal. We neglect surface reconstructions Pokrant and Whaley (1999); Puzder et al. (2003); Deglmann et al. (2002) and that the surface coverage with ligands is often not perfect Taylor et al. (2002), though these effects can be important especially for very small NCs. However, we concentrate on considerably larger NCs than in the before mentioned referenceses. Therefore, unlike previous TB work we concentrate here on size and the size dependence of the results obtained for the electronic structure of the NCs. The TB-parameters, which describe the coupling between the dot material and the ligand molecules, are chosen to be zero. This corresponds to an infinite potential barrier at the surface and is commonly used because of the larger band gap of the surrounding material Lippens and Lannoo (1989). An alternative approach to treat the ligand molecules is discussed by Sapra et al. in Ref. Sapra and Sarma, 2004. The influence of the organic ligands on the electronic structure can also be investigated more realistically in the framework of microscopic descriptions Pokrant and Whaley (1999); Rabani (2003); Puzder et al. (2003).
### IV.2 Single particle states and comparsion with experimental results
We have performed TB-calculations for finite, spherical, unstrained NCs of diameter between 1.82 nm and 4.85 nm (corresponding to 3-8 $`a`$, when $`a6.07`$ Å is the CdSe lattice constant of the conventional unit cell). The finite matrix diagonalizations yield both, the discrete eigenenergies and the eigenstates. For the largest NCs (of diameter 4.85 nm) results for the five lowest lying electron and hole eigenstates are shown in Fig. 8 again in the form of an isosurface plot. The lowest lying electronic state $`e_1`$ obviously has spherical symmetry and can be classified as a $`1s`$-state. Correspondingly the second state $`e_2`$ has the form of a $`2s`$-state and the states $`e_{3,4}`$ are $`p`$-states and $`e_5`$ is a $`d`$-like state. Despite the spherical symmetry of the system this simple classification is no longer possible for the hole states, however. Even the lowest lying hole state $`h_1`$ has no full rotational invariance, i.e. strictly speaking it cannot be classified as being an $`s`$-state. This is due to the intermixing of different atomic TB-valence electron states in the NC. Similarly the higher hole states $`h_2h_4`$ cannot clearly be classified as an $`s`$\- or $`p`$-like state. This is an effect, which simple effective mass models cannot account for, but which will have implications in the calculation of matrix elements between these states, which enter selection rules for optical transitions etc.
In the case of an ideal zinc blende structure as considered here we do not obtain any indications of quasi-metallic behavior, i.e. of a non-vanishing (quasi continuous) spectrum of states at the Fermi energy in contrast to previous work (assuming an ideal wurtzite structure for CdSe nanocrystals) Puzder et al. (2004); Leung and Whaley (1999); Pokrant and Whaley (1999). This is probably due to the fact that this quasi-metallic behavior is due to surface states in the case when no passivation and surface reconstruction is taken into account. These surface states are formed by the dangling bonds of unsaturated Se at the NC surface, which cause $`s`$-states in the band gap region Leung and Whaley (1999). In our simplified and restricted TB $`s_cp_a^3`$-basis set these $`s`$-orbitals at the anions (Se) are not taken into account. Therefore, these surface states, which in reality and in more realistic models are removed (i.e. energetically drawn down and filled) due to passivation and surface reconstruction, do not occur.
The discrete electronic states of semiconductor NCs are experimentally accessible by scanning tunneling microscopy (STM) Banin et al. (1999); Alperson et al. (1999); Banin et al. (1998). The tunnel current $`I`$ between the metallic tip of the STM and the CdSe nanocrystal, which is e.g. epitaxially electrodeposited onto a template-stripped gold film, is measured as a function of the bias voltage $`V`$. The conductance ($`dI/dV`$) is related to the local tunneling density of states. In the $`dI/dV`$ versus $`V`$ diagram, several discrete peaks can be observed. These peaks correspond to the addition energies (charging energies) of holes and electrons. The spacing between the various peaks can be attributed to the Coulomb charging (addition spectrum) and/or charge transfer into higher energy levels (excitation spectrum). From these measurements the energy gap $`E_{\text{gap}}^{\text{nano}}`$ as well as the splitting $`\mathrm{\Delta }_{e_1,e_2}`$ between electron ground state $`e_1`$ and the first excited state $`e_2`$ can be determined.
Alperson et al. Alperson et al. (1999) investigated CdSe nanocrystals with an STM. Here we compare our calculated energy gap $`E_{\text{gap}}^{\text{nano}}`$, which is given by the difference between the electron, $`e_1`$, and hole, $`h_1`$, groundstate, with measured data from Ref. Alperson et al., 1999. Figure 9 displays the results for CdSe NCs with diameters in between $`1.82`$ nm and $`4.85`$ nm. Alperson et al. Alperson et al. (1999) compare the STM results (dashed dotted line) with optical spectroscopy measurements (dotted line) from Ekimov et al. Ekimov et al. (1993). The overall agreement with the TB results is very good, especially for the larger NCs. Deviations in the case of the small 2 nm NC arise from surface reconstructions Puzder et al. (2004); Deglmann et al. (2002); Pokrant and Whaley (1999) which are neglected here. When the same calculation is done without spin-orbit coupling (TB-NO SO), the energy gap $`E_{\text{gap}}^{\text{nano}}`$ is always strongly overestimated by the TB-model, in particular for smaller nanocrystals. So the spin-orbit coupling is important for a satisfactory reproduction of the experimental results. For the calculations without spin-orbit coupling, the TB-parameters are re-optimized to reproduce the characteristic properties (band gap, effective masses) of the bulk material. The re-optimized parameters are given in Table 2.
Taking into account the electron spin, the lowest electron state $`e_1`$ is twofold degenerated and $`s`$-like. This is consistent with the experimentally observed doublet Alperson et al. (1999) in the $`dI/dV`$ characteristic. The next excited level is (quasi) sixfold degenerated. The spin-orbit coupling splits this into one twofold and one fourfold degenerate state Niquet et al. (2001). In the STM measurement Alperson et al. Alperson et al. (1999) observed such a higher multiplicity of the second group of peaks. This behavior has also experimentally Banin et al. (1999) and theoretically Niquet et al. (2001) been observed for InAs nanocrystals. The electron energy spectrum for NCs of different diameter is shown in Fig. 10 (a). Here the first five electron states $`e_1e_5`$ are displayed. Note that every state is twofold degenerated due to the spin.
For the hole states the situation is more complicated. Alperson et al. Alperson et al. (1999) observed a high density of states at negative bias. The distinction between addition and excitation peaks is difficult, due to the large number of possibilities and the close proximity between the charging energy and the level spacing. For the holes we obtain that the first two states $`(h_1,h_2)`$ and $`(h_3,h_4)`$ are fourfold degenerated. The energy splitting of these states is also very small. These results are consistent to the observations of Alperson et al. Alperson et al. (1999). Figure 10 (b) shows the hole energy versus diamater $`d`$ for the spherical CdSe NCs. Obviously, for all diameters displayed the states $`h_1h_4`$ are almost degenerate, i.e. including spin there is almost an 8-fold degeneracy of these states.
Furthermore the calculated splitting $`\mathrm{\Delta }_{e_1,e_2}=E_{e_2}E_{e_1}`$ between the first two electron states $`e_1`$ and $`e_2`$ is compared with experimentally observed results for this quantity. Figure 11 shows $`\mathrm{\Delta }_{e_1,e_2}`$ as a function of the nanocrystal diameter $`d`$. The influence of the spin-orbit coupling on our results is also investigated. We have done the calculations without (TB-NO SO) and with spin orbit-coupling (TB). The results of our TB-model for the splitting $`\mathrm{\Delta }_{e_1,e_2}`$ are compared with results obtained by STM Alperson et al. (1999) and by optical methods (optical) Ekimov et al. (1993). This splitting $`\mathrm{\Delta }_{e_1,e_2}`$ was independently determined experimentally by Guyot-Sionnest and Hines Guyot-Sionnest and Hines (1998) using infrared spectroscopy (IR). Without spin-orbit coupling the TB-model always overestimates the splitting $`\mathrm{\Delta }_{e_1,e_2}`$. Especially for smaller nanocrystals the spin-orbit coupling is very important to describe the electronic structure. With spin-orbit coupling the results of the TB-model show good agreement with the experimentally observed results.
## V Conclusion
We have applied an empirical $`s_cp_a^3`$ TB-model to the calculation of the electronic properties of II-VI semiconductor EQDs and NCs. Assuming a zinc blende lattice and (per spin direction) one $`s`$-like orbital at the cation sites and three $`p`$-orbitals at the anion sites, the TB-parameters for different materials (here ZnSe and CdSe) are determined so that the most essential properties (band gap, effective masses etc.) of the known band structure of the (3 dimensional) bulk materials are well reproduced by the TB band structure. Then a CdSe QD (on top of a two-dimensional, a few atomic layers thick WL) embedded within a ZnSe matrix is modelled by using the TB-parameters of the dot material for those sites occupied by CdSe and the ZnSe TB matrix elements for the remaining sites; suitable averages have to be chosen for intersite matrix elements over and for on-site matrix elements on anion (Se) sites at interfaces between QD and barrier material. Spherical CdSe NCs can be modelled similarly by setting the intersite matrix elements between surface atoms and atoms in the monolayer of surfactant material to zero. The effects of the spin-orbit interaction, the band offsets and for the EQDs strain effects are taken into account.
For the EQD systems the numerical diagonalization yields a discrete spectrum of bound electron and hole states localized in the region of the EQD. Energetically these discrete states are below the continuum of the WL states. We have investigated the dependence on the EQD size and find that the number of the bound states and their binding energy increases with increasing dot size, therefore the effective band gap decreases. We have also investigated the dependence of the bound eigenenergies and their degeneracy on strain and on the thickness of the WL. Looking at the states themselves one sees that conduction band (electron) states can be roughly classified as $`s`$-like, $`p`$-like, etc. states but the valence band (hole) states cannot be classified according to such simple ($`s`$,$`p`$,$`d`$) symmetries because they are determined by a mixing between the different (anion) $`p`$-states. This cannot be accounted for by simple effective mass models but it will be important for instance for the calculation of dipole matrix elements between electron and hole states which determine the selection rules for optical transitions. For the NCs the whole spectrum is discrete, but in spite of the spherical symmetry the hole states do not have the simple $`s`$,$`p`$,$`d`$-symmetry but are intermixtures of atomic $`p`$-orbitals. Even the lowest hole state has no spherical $`s`$-symmetry but it is 4-fold (8-fold including spin) degenerate. The spin-orbit interaction is very important. Including the spin-orbit interaction we obtained nearly perfect agreement with experimental results obtained by STM for the dependence of the band gap and of the splitting of the lowest electronic states on the diameter of the NC.
Compared to (two-band) effective mass Wojs et al. (1996); Grundmann et al. (1995) and multi-band $`\stackrel{}{k}\stackrel{}{p}`$-models Stier et al. (1998); Fonoberov and Baladin (2003); Pryor (1998) for EQDs our TB model clearly has the advantage of a microscopic, atomistic description. Different atoms and constituents of the nanostructure and their actual positions are considered, and this may lead to a reduction of symmetries (for instance the $`C_{2v}`$ symmetry instead of a $`C_{4v}`$-symmetry). This may automatically lift certain degeneracies and lead automatically to a splitting, for instance between $`e_2`$ and $`e_3`$ states, whereas an 8-Band-$`\stackrel{}{k}\stackrel{}{p}`$ model still yields degenerate $`e_2`$ and $`e_3`$ states Pryor (1998). The effects of inhomogeneous strain can be easily incorporated into a TB model by considering the deviations of the actual atomic positions from the ideal position in the bulk crystal. Only the (empirical) pseudopotential treatmentWang et al. (1999, 2000); Williamson et al. (2000); Kim et al. (1998) may be still superior and more accurate than the TB approach, but in a pseudopotential descripion a variation of the wave functions within the individual atoms is accounted for and a large number of basis states is required. Therefore, a TB description is simpler and quicker and allows for the investigation of larger nanostructures without loosing information on the essential, microscopic details of the structure. Compared to other TB models of QD structures, we do not consider free standing, isolated QDs (as in Ref. Saito et al., 1998) but we can describe realistic QDs (with a WL) embedded into another barrier material. We show here that a reduced $`s_cp_a^3`$-basis is already sufficient for a satisfying reproduction of properties like optical gaps, energy splittings, etc., and their size dependence. Much larger basis sets, namely a $`sp^3s^{}`$-basisSantoprete et al. (2003) or even a $`sp^3d^5s^{}`$-basisJancu et al. (1998) were used in previous TB-models of EQDs. Our reduced, smaller basis set, of course, leads to computational simplifications and allows for the treatment of larger QDs. Furthermore, we apply our TB-model to different materials than investigated previously, namely II-VI CdSe nanostructures, and we investigate also NCs, for which excellent agreement with experimental STM-results could be demonstrated.
In the future further applications of our TB-model for embedded semiconductor QDs and NCs are planned. Of course applications to QDs of other materials, for instance nitride systems, and other (e.g. wurtzite) crystal structures are possible. Furthermore, EQDs of other shape and size (dome-shaped, lens-shaped, truncated cones, etc.) or two coupled QDs or freestanding (capped and uncapped) QDs can be investigated. A combination with ab-initio calculations is also possible by determining the TB-parameters from a first-principles band structure calculation of the bulk material. Furthermore the influence of surface reconstructions and the surfactant material on the results for NCs should be investigated. Especially for small NCs these effects are important. Finally matrix elements of certain observables like dipole matrix elements between the calculated QD electron and hole states can be determined, which are important for selection rules and the optical properties of these systems.
###### Acknowledgements.
This work has been supported by a grant from the Deutsche Forschungsgemeinschaft No. Cz-31/14-1,2 and by a grant for CPU time from the John von Neumann Institute for Computing at the Forschungszentrum Jülich. |
warning/0507/hep-ex0507031.html | ar5iv | text | # Measurement of branching fractions and mass spectra of 𝐵→𝐾𝜋𝜋𝛾
BABAR-PUB-05/31
SLAC-PUB-11288
thanks: Deceased
The BABAR Collaboration
## Abstract
We present a measurement of the partial branching fractions and mass spectra of the exclusive radiative penguin processes $`BK\pi \pi \gamma `$ in the range $`m_{K\pi \pi }<1.8\mathrm{GeV}/c^2`$. We reconstruct four final states: $`K^+\pi ^{}\pi ^+\gamma `$, $`K^+\pi ^{}\pi ^0\gamma `$, $`K_S^0\pi ^{}\pi ^+\gamma `$, and $`K_S^0\pi ^+\pi ^0\gamma `$, where $`K_S^0\pi ^+\pi ^{}`$. Using 232 million $`e^+e^{}B\overline{B}`$ events recorded by the BABAR experiment at the PEP-II asymmetric-energy storage ring, we measure the branching fractions $`(B^+K^+\pi ^{}\pi ^+\gamma )=(2.95\pm 0.13(\mathrm{stat}.)\pm 0.20(\mathrm{syst}.))\times 10^5`$, $`(B^0K^+\pi ^{}\pi ^0\gamma )=(4.07\pm 0.22(\mathrm{stat}.)\pm 0.31(\mathrm{syst}.))\times 10^5`$, $`(B^0K^0\pi ^+\pi ^{}\gamma )=(1.85\pm 0.21(\mathrm{stat}.)\pm 0.12(\mathrm{syst}.))\times 10^5`$, and $`(B^+K^0\pi ^+\pi ^0\gamma )=(4.56\pm 0.42(\mathrm{stat}.)\pm 0.31(\mathrm{syst}.))\times 10^5`$.
preprint: BABAR-PUB-05/31preprint: SLAC-PUB-11288
In the standard model (SM) the radiative penguin decay $`BX_s\gamma `$, where $`X_s`$ is a hadronic system with unit strangeness, proceeds via weak-interaction loop diagrams. New physics, beyond the SM, may also contribute to the loop amplitude, and lead to differences from the SM. This possibility has been pursued in inclusive measurements, which are theoretically clean but experimentally challenging, and in exclusive measurements, such as $`BK\pi \gamma `$. We report measurements of the branching fractions and mass spectra for the decays $`BK\pi \pi \gamma `$ in four channels. SM predictions of the rates and resonance structure of these decays have large uncertainties SMpred . The $`K^+\pi ^+\pi ^{}\gamma `$ and $`K^0\pi ^+\pi ^{}\gamma `$ decay channels have previously been observed Belle . Throughout this Letter, stated decays include charge conjugate modes.
The decays $`BK\pi ^+\pi ^0\gamma `$, which have not previously been observed, are of particular interest because these three-body hadronic states permit the measurement, given sufficient statistics, of the photon polarization Gronau02 . The polarization measurement depends on the interference between processes such as $`(K\pi ^+)\pi ^0\gamma `$ and $`(K\pi ^0)\pi ^+\gamma `$, where $`()`$ indicates resonant substructure. This measurement may be compared with the SM prediction of nearly complete left-handed polarization.
We use a sample of $`(232\pm 1.5)\times 10^6`$ $`B\overline{B}`$ pairs in a 210.9 fb<sup>-1</sup> dataset collected at the $`\mathrm{{\rm Y}}(4S)`$ resonance with the BABAR detector at the PEP-II asymmetric-energy $`e^+e^{}`$ collider. For background studies, we also use a $`21.7\text{fb}^1`$ sample collected below the $`B\overline{B}`$ threshold. The measurement procedure was designed using simulated signal and background events, data in sideband kinematic regions, and reconstructed $`BD\pi ^+`$, $`DK\pi \pi `$ decays. Only after we established the selection and fit procedures did we examine signal candidates in the data sample.
A description of the detector exists elsewhere detector . For this measurement, the most important detector elements are the five-layer silicon microstrip tracking detector (SVT) and the forty-layer drift chamber (DCH), situated in a 1.5 T solenoidal magnetic field, which measure charged particle momenta; the CsI(Tl) electromagnetic calorimeter (EMC), which measures the energies and directions of the photons; and the detector of internally reflected Cherenkov light (DIRC). The DIRC response and energy loss ($`\mathrm{d}E/\mathrm{d}x`$) measured in the SVT and DCH are used to identify charged kaons and pions.
We reconstruct the photon candidate in the $`K\pi \pi \gamma `$ decay from an EMC shower not associated with a charged track. The photon must be in the fiducial region of the EMC, have a shower-profile consistent with a single photon, and be well-separated from other showers. To remove photons from $`\pi ^0`$ ($`\eta `$) decays, we combine the candidate with other photons having energies of at least 50 (250)$`\mathrm{MeV}/c^2`$, and reject it if the invariant mass of any combination is within 25 (40)$`\mathrm{MeV}/c^2`$ of the $`\pi ^0`$ ($`\eta `$) mass.
We select $`K^\pm `$ and $`\pi ^\pm `$ candidates from charged tracks consistent with a kaon or pion mass hypothesis in the DIRC and in the $`\mathrm{d}E/\mathrm{d}x`$ in the SVT and DCH. We reconstruct $`K_S^0`$ candidates from pairs of oppositely-charged tracks, and determine the decay vertex with a fit. We require that the invariant mass falls within 11$`\mathrm{MeV}/c^2`$ of the $`K_S^0`$ mass; that the distance between the $`B`$ decay vertex and the $`K_S^0`$ vertex exceeds 5 times the uncertainty on the distance; and that the angle between the $`K_S^0`$ trajectory and its momentum is less than 100 mrad. We reconstruct $`\pi ^0`$ candidates from pairs of EMC showers each with energy $`>`$50$`\mathrm{MeV}`$. We require the invariant mass to be within 16$`\mathrm{MeV}/c^2`$ of the $`\pi ^0`$ mass, and that the energy of each pair in the $`\mathrm{{\rm Y}}(4S)`$ center of mass (CM) frame exceeds 450$`\mathrm{MeV}`$; this last selection is about 83% efficient.
The dominant source of background is continuum production of light quark-antiquark pairs, in which a high-energy photon typically is produced either by initial state radiation, or from the decay of a $`\pi ^0`$ or $`\eta `$ in which one photon is not detected. To reject these backgrounds, we construct a Fisher discriminant Fisher from the polar angle of the $`B`$ candidate in the CM frame, the angle between the thrust axis of the $`B`$ and the thrust axis of the remaining charged and neutral particles, and the ratio of the second to zeroth angular moments of the remaining charged and neutral particles around the thrust axis of the $`B`$. We optimize the coefficients independently in each channel to discriminate between simulated signal and continuum.
We perform a geometric fit to the reconstructed $`B`$ candidate, with production vertex constrained to the nominal beam spot, rejecting the candidate if the final state is inconsistent with decay from a single vertex. We define $`\mathrm{\Delta }E^{}E_B^{}E_{\text{beam}}^{}`$ and $`m_{\mathrm{ES}}\sqrt{E_{\text{beam}}^2𝐩_B^2}`$, where $`E_B^{}`$ and $`𝐩_B^{}`$ are the CM energy and momentum of the $`B`$ candidate, and $`E_{\text{beam}}^{}`$ is the CM energy of each beam. We require $`m_{\mathrm{ES}}>5.2\mathrm{GeV}/c^2`$ and $`|\mathrm{\Delta }E^{}|<0.15\mathrm{GeV}`$. We also require that the invariant mass of the $`K\pi \pi `$ system, $`m_{K\pi \pi }`$, fall below 1.8$`\mathrm{GeV}/c^2`$; this eliminates much of a rising continuum background with very little expected signal loss. It also removes $`K\pi \pi `$ combinations from $`D`$ decays, in $`BD\pi ^0`$ and $`BD\eta `$ where the $`\pi ^0`$ or $`\eta `$ are mis-reconstructed as a photon. In an event in which we reconstruct multiple candidates in one channel that pass the selection requirements (occurring in 11-27% of selected signal events, depending on the channel), we keep the candidate with the largest vertex probability (with the best $`\pi ^0`$ mass in the $`K_S^0\pi ^+\pi ^0\gamma `$ channel; with the $`\pi ^0`$ mass as a tie breaker in the $`K^+\pi ^{}\pi ^0\gamma `$ channel) and reject the others. Candidates reconstructed in different channels are allowed in the same event. The dependence of the efficiency of our selection requirements on intermediate resonance and on $`m_{K\pi \pi }`$ has been checked and found to be small; systematic uncertainties are discussed later.
The dominant backgrounds from $`B\overline{B}`$ events after the selection criteria have been applied are $`bs\gamma `$ processes. We categorize these backgrounds: (i) “crossfeed” from mis-reconstructed $`K\pi \pi \gamma `$ decays, such as by choosing incorrectly a particle from the other $`B`$; (ii) $`BK\pi \gamma `$ decays that combine with a track from the other $`B`$ to form a $`K\pi \pi \gamma `$ candidate; and (iii) backgrounds from all other $`bs\gamma `$ decays. A crossfeed candidate may be reconstructed in the same decay channel in which it is produced, or in a different channel, and can also be produced in a $`BK\pi \pi \gamma `$ decay that is not used in this analysis (such as $`B^+K^+\pi ^0\pi ^0\gamma `$). We model our signal as well as crossfeed backgrounds with simulated $`BK_X\gamma `$ decays, where $`K_X`$ is any of the five lowest-lying $`J>0`$ kaon resonances above the $`K^{}(892)`$. We study backgrounds from $`K\pi \gamma `$ using simulated $`BK^{}(892)\gamma `$ and $`BK_2^{}(1430)\gamma `$ decays. We study backgrounds from other $`bs\gamma `$ decays using an inclusive simulation according to the model of Kagan and Neubert KN99 with $`m_b=4.8\mathrm{GeV}/c^2`$, tuned to match multiplicity distributions measured in inclusive $`bs\gamma `$ decays semiinclusive . The largest final background contributions from $`bs\gamma `$ processes are crossfeed backgrounds, for which we obtain yields ranging from 55% to 95% of the signal yields.
We estimate other sources of background candidates from $`B`$ decays other than $`bs\gamma `$ processes by simulating generic $`B`$ decays. We pay special attention to $`B`$ decays with $`K\pi \pi \pi ^0`$ and $`K\pi \pi \eta `$ final states; if the $`\pi ^0`$ or $`\eta `$ decays asymmetrically and we don’t detect the lower-energy photon, the kinematic properties of the resulting $`B`$ candidate may resemble a signal candidate. We study these decays using high-statistics simulated samples, and look for signal candidates that are reconstructed from a single $`BK\pi \pi \pi ^0`$ or $`BK\pi \pi \eta `$ decay. We expect to reconstruct fewer than two such candidates per channel.
We perform a maximum likelihood fit to the joint $`m_{\mathrm{ES}}`$$`\mathrm{\Delta }E^{}`$ distribution of our selected candidates. We fit all four channels simultaneously to account for crossfeed backgrounds between channels. The likelihood function contains terms for correctly reconstructed signal candidates, crossfeed background candidates between all 16 combinations of the production and reconstruction channels, backgrounds from $`BK\pi \gamma `$ and from other $`bs\gamma `$ decays, and backgrounds from continuum events. We have determined from simulations that the dominant continuum background component adequately accounts for combinatoric backgrounds from other $`B\overline{B}`$ decays, which do not show strong peaks in $`m_{\mathrm{ES}}`$ and $`\mathrm{\Delta }E^{}`$.
The likelihood function for a candidate reconstructed in decay channel $`i`$ with kinematic variables $`y(m_{\mathrm{ES}},\mathrm{\Delta }E^{})`$ is given by,
$$\begin{array}{c}^i(y)=N_{B\overline{B}}\left(^iϵ_s^if_s^i(y)+\underset{j}{}^jϵ_x^{ji}f_x^{ji}(y)\right)\hfill \\ \hfill +n_c^if_c^i(y)+n_b^if_b^i(y),\end{array}$$
where $`N_{B\overline{B}}`$ is the number of $`B\overline{B}`$ pairs in our dataset; $`^i`$ is the branching fraction for decay channel $`i`$; $`ϵ_s^i`$ and $`f_s^i`$ are the efficiency and probability density function (PDF) for correctly reconstructed signal candidates in decay channel $`i`$; $`ϵ_x^{ji}`$ and $`f_x^{ji}`$ are the efficiency and PDF for crossfeed background candidates produced in channel $`j`$ and reconstructed in channel $`i`$; $`n_c^i`$ and $`f_c^i`$ are the yield and PDF for backgrounds from continuum and generic $`B\overline{B}`$ decay events in channel $`i`$; and $`n_b^i`$ and $`f_b^i`$ are the yield and PDF for backgrounds from other $`bs\gamma `$ processes in channel $`i`$. We further parameterize the likelihood function by the four data-taking runs during which data were collected, accounting for slight changes in experimental conditions.
The branching fractions $`^i`$, yields $`n_c^i`$, and shape parameters of $`f_c^i`$ are varied in the fit; other efficiencies, yields, and PDF shapes are fixed from simulation studies. We parameterize $`f_s^i`$ as the product of Crystal Ball functions crystalball of $`m_{\mathrm{ES}}`$ and of $`\mathrm{\Delta }E^{}`$, $`f_x^i`$ as the product of a Crystal Ball function of $`m_{\mathrm{ES}}`$ and a linear function of $`\mathrm{\Delta }E^{}`$, and $`f_c^i`$ as the product of an Argus function argus of $`m_{\mathrm{ES}}`$ and an exponential function of $`\mathrm{\Delta }E^{}`$. We use a binned parameterization for $`f_b^i`$. As the signal and crossfeed terms are both scaled by the parameters $`^i`$, the crossfeed background yields vary with the signal branching fractions, and we measure the branching fractions from yields of both signal and crossfeed candidates.
Table 1 shows the fit results. Projections in $`m_{\mathrm{ES}}`$, along with the fit results, are displayed in Fig. 1. The fit probability ($`P`$-value) is evaluated with a likelihood ratio statistic , assuming Poisson-distributed bin contents, to be 10%. The distribution of the test statistic under the null hypothesis is evaluated by simulation.
Figure 2 shows background-subtracted $`m_{K\pi \pi }`$ mass spectra. Background subtraction is achieved using the results of the fits to calculate event-by-event weights to extract the signal component sPlots . We present branching fractions in bins of $`m_{K\pi \pi }`$, which are largely model-independent, instead of extracting $`BK_X\gamma `$ branching fractions for specific $`K_X`$ resonances. Disentangling the resonance structure requires careful modeling of amplitudes and relative phases of interfering processes, including in the decays of the $`K_X`$ resonances, not all of which are well measured. A partial wave analysis to extract the resonance structure and measure the photon polarization should be possible with future datasets.
We validate the procedure for extracting branching fractions and $`m_{K\pi \pi }`$ distributions using fits to simulated samples. We verify that the branching fractions and mass spectra obtained from these toy fits reproduce on average the simulation inputs. We use the same procedure to extract the $`m_{K\pi \pi }`$ distributions for continuum and generic $`B\overline{B}`$ backgrounds and for backgrounds from $`bs\gamma `$ decays; these are consistent with the expected distributions.
Systematic uncertainties arise from various sources, shown in Table 2. The largest sources are: (i) The $`\mathrm{{\rm Y}}(4S)`$ branching fractions to $`B^+B^{}`$ and $`B^0\overline{B}^0`$ are each assumed to be 0.5. We assign a 2.6% systematic uncertainty to this, based on current information upsBF . (ii) The uncertainty on the photon selection efficiency determined from simulated events is estimated to be 2.7%. (iii) From studies of $`BD\pi ^\pm `$, $`DK\pi \pi `$ events, we assign an uncertainty of 4.2% to the charged kaon identification efficiency. (iv) The uncertainty of the $`\pi ^0`$ selection efficiency is estimated at 3.0%. (v) There is considerable uncertainty in the models we use to estimate backgrounds, including cross-feed dependence, from $`bs\gamma `$ processes. We estimate the effect of this uncertainty on both the branching fractions and mass spectra by simulating these backgrounds with substantially different models. The largest effect is in the $`K_S^0\pi ^{}\pi ^+\gamma `$ channel, where the uncertainty is 4.0%. (vi) We measure a shift in the beam energy in $`BD\pi ^\pm `$ decays, on average 0.6$`\mathrm{MeV}`$; we estimate the effect of this on our fits. (vii) We estimate bias in the fit due to uncertain parameterization of the signal and background PDFs. The largest effect is in the $`K_S^0\pi ^+\pi ^0\gamma `$ channel, where the uncertainty is 3.5%.
We have measured branching fractions for $`BK\pi \pi \gamma `$ in four decay channels for $`m_{K\pi \pi }<1.8\mathrm{GeV}/c^2`$. The $`K\pi ^+\pi ^{}`$ channels are consistent with the previous measurement Belle . We present first observations of decays in the $`K\pi ^+\pi ^0\gamma `$ channels that are important to measuring the photon polarization. The branching fractions are relatively large in the context of $`BX_s\gamma `$ decays, providing encouragement that a polarization measurement may be possible with future datasets. Mass spectra for the $`K\pi \pi `$ system are also presented. We observe an enhancement near 1.3$`\mathrm{GeV}/c^2`$ and substantial branching fractions at higher masses. Untangling the resonant contributions presents a challenge for the polarization measurement.
We are grateful for the excellent luminosity and machine conditions provided by our PEP-II colleagues, and for the substantial dedicated effort from the computing organizations that support BABAR. The collaborating institutions wish to thank SLAC for its support and kind hospitality. This work is supported by DOE and NSF (USA), NSERC (Canada), CEA and CNRS-IN2P3 (France), BMBF and DFG (Germany), INFN (Italy), FOM (The Netherlands), NFR (Norway), MES (Russia), MEC (Spain), and STFC (United Kingdom). Individuals have received support from the Marie Curie EIF (European Union) and the A. P. Sloan Foundation. |
warning/0507/quant-ph0507207.html | ar5iv | text | # One-Dimensional Three-State Quantum Walk
## I Introduction
Quantum walks Aharonov1993 ; Mayer1996 ; Ambainis2001 ; Aharonov2001 ; Tregenna2002 are very simple quantum processes, however they are connected with a wide variety of field. Studies of the quantum walks are initially motivated by developing techniques for quantum algorithms. For example, the Grover’s search algorithm Grover1997 , which is one of the most famous quantum algorithms, is related to a discrete quantum walk Shenvi2002 ; Ambainis2004 . Oka et al. Oka2005 studied breakdown of an electric-field driven system by mapping the Landau-Zener transition dynamics to a quantum walk.
The Hadamard walk plays a key role in studies of the quantum walk and it has been analyzed in detail. Thus the generalization of the Hadamard walk is one of fascinating challenges. The simplest classical random walker on a one-dimensional lattice moves to the left or to the right with probability $`1/2`$. On the other hand, the quantum walker according to the Hadamard walk on a line moves both to the left and to the right. It is well known that the spatial distribution of probability of finding a particle governed by the Hadamard walk after long time is quite different from that of the classical random walk, see Kon2002 ; Kon2005 ; GJS ; KFK , for examples. However there are some commonalities. In both walks the probability of finding probability at a fixed lattice site converges to zero after infinite long time.
Let us consider a classical random walker in which it can stay at the same position with non-zero probability in a single time-evolution. In the limit of long time, dose the probability of finding a walker at a fixed position converges to a positive value ? If the probability of staying at the same position in a single step is very close to 1, the walker diffuses very slowly. However the probability of existence at a fixed position surely converges to zero if the jumping probability is not zero. The classical random walker who can remain the same position is essentially regarded as the same process with the random walk right or left with probability with 1/2 by scaling the time. Is the same conclusion valid for a quantum walk? The answer is “No”. We show that the profile of the Hadamard walk changes drastically by appending only one degree of freedom to the inner states.
If the quantum particle in three-state quantum walk exists at only one site initially, the particle is trapped with high probability near the initial position. Similar localization has been already seen in other quantum walks. The first simulation showing the localization was presented by Mackay et al. Mackay2002 in studying the two-dimensional Grover walk. After that, more refined simulations were performed by Tregenna et al. Tregenna2003 and an exact proof on the localization was given by Inui et al. Inui2004 . The second is found in the four-state quantum walk Brun2003B ; Inui2005Physica . In this walk, a particle moves not only the nearest sites but also the second nearest sites according to the four inner states. The four-state quantum walk is also a generalized Hadamard walk, and it is similar with the three-state quantum walk. The significant difference between them is that the wavefunction of the four-state quantum walk dose not converge, but that of the three-state quantum walk converges in the limit as time tends to infinity. We focus on localized stationary distribution of the three-state quantum particle and calculate it rigorously. Moreover Konno Kon2005b proved that a weak limit distribution of a continuous-time rescaled two-state quantum walk has a similar form to that of the discrete-time one. In fact, both limit density functions for two-state quantum walks have two peaks at the two end points of the supports. As a corollary, it is easily shown that a weak limit distribution of the corresponding continuous-time three-state quantum walk has the same shape as that of the two-state walk. So the localization does not occur in the continuous-time case in contrast with the discrete-time case. In this situation the main aim of this paper is to show that the localization occurs for a discrete-time three-state quantum walk rigorously.
The rest of the paper is organized as follows. After defining the three-state quantum walk, eigenvalues and eigenvectors of the time evolution operator are calculated, and the wave function is shown in Sec. II. A time-averaged probability of finding the particle is introduced in Sec. III, and it is shown that the time-averaged probability of the three-state quantum walk converges to non-zero values. In Sec. IV, we prove the probability of finding the particle at a fixed itself converges to a non-zero value after infinite long time.
## II Definition of the three-state quantum walk
The three-state quantum walk (3QW) considered here is a kind of the generalized Hadamard walk on a line. The particle ruled by 3QW is characterized in the Hilbert space which is defined by a direct product of a chirality state space $`|s\{|L,|0,|R\}`$ and a position space $`|n\{\mathrm{},|2,|1,|0,|1,|2,\mathrm{}\}`$. The chirality states are transformed at each time step by the next unitary transformation:
$`|L={\displaystyle \frac{1}{3}}(|L+2|0+2|R),|0={\displaystyle \frac{1}{3}}(2|L|0+2|R),|R={\displaystyle \frac{1}{3}}(2|L+2|0|R).`$
Let $`\mathrm{\Psi }(n,t)[\psi _L(n,t),\psi _0(n,t),\psi _R(n,t)]^T`$ be the amplitude of the wave function of the particle corresponding to the chiralities “L”, “0” and “R” at the position $`n`$ and the time $`t\{0,1,2,\mathrm{}\}`$, where $`T`$ denotes the transpose operator and $``$ is the set of integers. We assume that a particle exists initially at the origin. Then the initial quantum states are determined by $`[\psi _L(0,0),\psi _0(0,0),\psi _R(0,0)][\alpha ,\beta ,\gamma ],`$ where $`\alpha ,\beta ,\gamma `$ with $`|\alpha |^2+|\beta |^2+|\gamma |^2=1`$. Here $``$ is the set of complex numbers.
Before we define the time-evolution of the wavefunction, we introduce the following three operators:
$`U_L={\displaystyle \frac{1}{3}}\left[\begin{array}{ccc}1& 2& 2\\ 0& 0& 0\\ 0& 0& 0\end{array}\right],U_0={\displaystyle \frac{1}{3}}\left[\begin{array}{ccc}0& 0& 0\\ 2& 1& 2\\ 0& 0& 0\end{array}\right],U_R={\displaystyle \frac{1}{3}}\left[\begin{array}{ccc}0& 0& 0\\ 0& 0& 0\\ 2& 2& 1\end{array}\right].`$
If the matrix $`U_L`$ is applied to the function $`\mathrm{\Psi }(n,t)`$, the only “L”-component is selected after carrying out the superimposition between $`\psi _L(n,t)`$, $`\psi _0(n,t)`$ and $`\psi _R(n,t)`$. Similarly the “0”-componet and “R”-componet are selected in $`U_0\mathrm{\Psi }(n,t)`$ and $`U_R\mathrm{\Psi }(n,t)`$.
We now define the time evolution of the wavefunction by
$`\mathrm{\Psi }(n,t)`$ $`=`$ $`U_L\psi (n+1,t)+U_0\psi (n,t)+U_R\psi (n1,t).`$
One finds clearly that the chiralities “L” and “R” are corresponding to the left and the right, and the chirality “0” is corresponding to the neutral state for the motion.
Using the Fourier analysis, which is often used in the calculations of quantum walks, we obtain the wavefunction. The spatial Fourier transformation of $`\mathrm{\Psi }(k,t)`$ is defined by
$`\stackrel{~}{\mathrm{\Psi }}(k,t)={\displaystyle \underset{n}{}}\mathrm{\Psi }(n,t)e^{ikn}.`$
The dynamics of wavefunction in the Fourier domain is given by
$`\stackrel{~}{\mathrm{\Psi }}(k,t+1)`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left[\begin{array}{ccc}e^{ik}& 0& 0\\ 0& 1& 0\\ 0& 0& e^{ik}\end{array}\right]\left[\begin{array}{ccc}1& 2& 2\\ 2& 1& 2\\ 2& 2& 1\end{array}\right]\stackrel{~}{\mathrm{\Psi }}(k,t)`$ (8)
$``$ $`\stackrel{~}{U}\stackrel{~}{\mathrm{\Psi }}(k,t).`$ (9)
Thus the solution of Eq. (9) is formally given by $`\stackrel{~}{\mathrm{\Psi }}(k,t)=\stackrel{~}{U}^t\stackrel{~}{\mathrm{\Psi }}(k,0)`$. Let $`e^{i\theta _{j,k}}`$ and $`|\mathrm{\Phi }_k^j`$ be the eigenvalues of $`\stackrel{~}{U}`$ and the orthonormal eigenvector corresponding to $`e^{i\theta _{j,k}}`$ $`(j=1,2,3)`$. Since the matrix $`\stackrel{~}{U}`$ is a unitary matrix, it is diagonalizable. Therefore the wavefunction $`\stackrel{~}{\mathrm{\Psi }}(k,t)`$ is expressed by
$`\stackrel{~}{\mathrm{\Psi }}(k,t)=\left({\displaystyle \underset{j=1}{\overset{3}{}}}e^{i\theta _{j,k}t}|\mathrm{\Phi }_k^j\mathrm{\Phi }_k^j|\right)\stackrel{~}{\mathrm{\Psi }}(k,0),`$
where $`\stackrel{~}{\mathrm{\Psi }}(k,0)=[\alpha ,\beta ,\gamma ]^T^3`$ with $`|\alpha |^2+|\beta |^2+|\gamma |^2=1.`$ The eigenvalues are given by
$`\theta _{j,k}`$ $`=`$ $`\{\begin{array}{cc}0,& j=1,\\ \theta _k,& j=2,\\ \theta _k,& j=3,\end{array}`$ (13)
$`\mathrm{cos}\theta _k`$ $`=`$ $`{\displaystyle \frac{1}{3}}(2+\mathrm{cos}k),`$
$`\mathrm{sin}\theta _k`$ $`=`$ $`{\displaystyle \frac{1}{3}}\sqrt{(5+\mathrm{cos}k)(1\mathrm{cos}k)},`$ (14)
for $`k[\pi ,\pi ).`$ The eigenvectors of $`\stackrel{~}{U}`$ which are orthonormal basis are obtained after some algebra:
$`|\mathrm{\Phi }_k^j`$ $`=`$ $`\sqrt{c_k(\theta _{j,k})}\left[\begin{array}{cc}\frac{1}{1+e^{i(\theta _{j,k}k)}}& \\ \frac{1}{1+e^{i\theta _{j,k}}}& \\ \frac{1}{1+e^{i(\theta _{j,k}+k)}}& \end{array}\right],`$ (18)
where
$`c_k(\theta )=2\left\{{\displaystyle \frac{1}{1+\mathrm{cos}(\theta k)}}+{\displaystyle \frac{1}{1+\mathrm{cos}\theta }}+{\displaystyle \frac{1}{1+\mathrm{cos}(\theta +k)}}\right\}^1.`$
In the above calculation, it is most important to emphasize that there is the eigenvalue “1” independently on the value of $`k`$. The eigenvalues of the Hadamard walk are given by $`e^{i\theta _k}`$ and $`e^{i(\pi \theta _k)}`$ with the arguments satisfying $`\mathrm{sin}\theta _k=\mathrm{sin}k/\sqrt{2}`$. Therefore the eigenvalues does not take the value “1” except the special case $`k=0`$. We have shown the existence of strongly degenerated eigenvalue such as “1” in Eq. (14) is necessary condition in the quantum walks showing the localization Inui2004 . The significant difference between 3QW and the previous quantum walks showing the localization such as the Grover walk and the four-state quantum walk is the number of the degenerate eigenvalues of the time evolution matrix. There are two degenerate eigenvalues with values “1” and “-1” independently of the value of $`k`$ in both the Grover walk and the four-state quantum walk. In contrast there is only one degenerate eigenvalue independently of the value $`k`$ for 3QW. This distinctive property causes the particular time evolution in 3QW.
Let us number each chirality “L”,“0”, and “R” using $`l=1`$, 2, and 3, respectively. Then the wavefunction in real space is obtained by the inverse Fourier transform: for $`\alpha ,\beta ,\gamma `$ with $`|\alpha |^2+|\beta |^2+|\gamma |^2=1`$,
$`\mathrm{\Psi }(n,t;\alpha ,\beta ,\gamma )`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _\pi ^\pi }\stackrel{~}{\mathrm{\Psi }}(k,t)e^{ikn}𝑑k`$ (19)
$`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _\pi ^\pi }\left({\displaystyle \underset{j=1}{\overset{3}{}}}e^{i\theta _{j,k}t}|\mathrm{\Phi }_k^j\mathrm{\Phi }_k^j|\stackrel{~}{\mathrm{\Psi }}(k,0)\right)e^{ikn}𝑑k`$
$`=`$ $`[\mathrm{\Psi }(n,t;1;\alpha ,\beta ,\gamma ),\mathrm{\Psi }(n,t;2;\alpha ,\beta ,\gamma ),\mathrm{\Psi }(n,t;3;\alpha ,\beta ,\gamma )]^T`$
$`=`$ $`{\displaystyle \underset{j=1}{\overset{3}{}}}[\mathrm{\Psi }_j(n,t;1;\alpha ,\beta ,\gamma ),\mathrm{\Psi }_j(n,t;2;\alpha ,\beta ,\gamma ),\mathrm{\Psi }_j(n,t;3;\alpha ,\beta ,\gamma )]^T,`$
where
$`\mathrm{\Psi }_j(n,t;l;\alpha ,\beta ,\gamma )={\displaystyle \frac{1}{2\pi }}{\displaystyle _\pi ^\pi }c_k(\theta _{j,k})\phi _k(\theta _{j,k},l)e^{i(\theta _{j,k}t+kn)}𝑑k,(l=1,2,3)`$ (20)
with
$`\phi _k(\theta ,l)=\zeta _{l,k}(\theta )[\alpha \overline{\zeta _{1,k}(\theta )}+\beta \overline{\zeta _{2,k}(\theta )}+\gamma \overline{\zeta _{3,k}(\theta )}],(l=1,2,3)`$
$`\zeta _{1,k}(\theta )=(1+e^{i(\theta k)})^1,\zeta _{2,k}(\theta )=(1+e^{i\theta })^1,\zeta _{3,k}(\theta )=(1+e^{i(\theta +k)})^1,`$ (21)
where $`\overline{z}`$ is conjugate of $`z.`$ Sometimes we omit the initial qubit state $`[\alpha ,\beta ,\gamma ]`$ such as $`\mathrm{\Psi }_j(n,t;l)=\mathrm{\Psi }_j(n,t;l;\alpha ,\beta ,\gamma )`$. The probability of finding the particle at the position $`n`$ and time $`t`$ with the chirality $`l`$ is given by $`P(n,t;l)=|\mathrm{\Psi }(n,t;l)|^2.`$ Thus the probability of finding the particle at the position $`n`$ and time $`t`$ is $`P(n,t)=_{l=1}^3P(n,t;l).`$
## III time-averaged probability
We focus our attention on the spatial distribution of the probability of finding the particle after long time. Equation (19) is rather complicated, but the value of $`lim_t\mathrm{}P(n,t)`$ can be exactly calculated in following sections. We start showing a numerical result for the probability of finding the particle at the origin before carrying analytical calculation. Figure 1 shows the time dependence of $`P(0,t)`$ with the initial state $`\alpha =i/\sqrt{2},\beta =0,\gamma =1/\sqrt{2}`$. The probability decreases quickly near $`t=0`$ and fluctuates near 0.2. The inserted figure in Fig. 1 is plotted as the inversed time, and the amplitude of fluctuation decreases as the time. The horizontal dashed line in Fig. 1 is the value at the origin of the time-averaged probability defined by
$`\overline{P}_{\mathrm{}}(0;\alpha ,\beta ,\gamma )`$ $`=`$ $`\underset{N\mathrm{}}{lim}\left(\underset{T\mathrm{}}{lim}{\displaystyle \frac{1}{T}}{\displaystyle \underset{l=1}{\overset{3}{}}}{\displaystyle \underset{t=0}{\overset{T1}{}}}P_N(0,t;l;\alpha ,\beta ,\gamma )\right),`$
where $`P_N(n,t;l;\alpha ,\beta ,\gamma )`$ is the probability of finding a particle with the chirality $`l`$ at the position $`n`$ and time $`t`$ on a cyclic lattice containing $`N`$ sites. The time-averaged probability introduced here has already used to study the quantum walk. In the case of the Hadamard walk on a cycle containing odd sites, the time-averaged probability takes a value $`1/N`$ independently of the initial states Aharonov2001 ; Inui2005 . Therefore the time-averaged probability converges to zero in the limit $`N\mathrm{}`$. On the other hand the time-averaged probability of quantum walks which exhibits the localization converges to a non-zero value. For this reason we firstly calculate time-averaged probability of 3QW, and we will show in the next section that the probability $`P(n,t)`$ itself converges as well.
We calculate the time-averaged probability at the origin in 3QW on a cycle with $`N`$ sites. We assume that the number $`N`$ is odd. The argument of eigenvalues of 3QW with a finite $`N`$ is $`\theta _{j,2m\pi /N}`$ for $`m[(N1)/2,(N1)/2]`$. Since the eigenvalues corresponding to $`m`$ is the same with the eigenvalues corresponding to $`m`$, the wavefunction is formally expressed by
$`\mathrm{\Psi }_N(0,t;l;\alpha ,\beta ,\gamma )`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{3}{}}}{\displaystyle \underset{m=0}{\overset{(N1)/2}{}}}c_{j,m,l}(N)e^{i\theta _{N,j,m}t},`$
where $`\theta _{N,j,m}`$ is $`\theta _{j,2m\pi /N}`$. We note here that the coefficients $`c_{j,m,l}(N)`$ depend on the initial state $`[\alpha ,\beta ,\gamma ]`$, but we omit to describe it. Then the probability $`P_N(0,t;\alpha ,\beta ,\gamma )`$ at the origin is given by
$`P_N(0,t;\alpha ,\beta ,\gamma )`$ $`=`$ $`{\displaystyle \underset{l_1,l_2,j_1,j_2=1}{\overset{3}{}}}{\displaystyle \underset{m_1,m_2=0}{\overset{(N1)/2}{}}}c_{j_1,m_1,l_1}^{}(N)c_{j_2,m_2,l_2}(N)e^{i(\theta _{N,j_2,m_2}\theta _{N,j_1,m_1})t}.`$
The coefficients $`c_{j,m,l}(N)`$ are determined from the product of eigenvectors. Although lengthy calculations are required to express the coefficients $`c_{j,m,l}(N)`$ explicitly, it is shown below that coefficients $`c_{j,m,l}(N)`$ except $`j=1`$ dose not contribute $`P_N(t)`$. Noting that the following equation
$`\underset{T\mathrm{}}{lim}{\displaystyle \frac{1}{T}}{\displaystyle \underset{t=0}{\overset{T1}{}}}e^{i\theta t}`$ $`=`$ $`\{\begin{array}{cc}1,& \theta =0,\\ 0,& \theta 0,\end{array}`$
we have
$`\overline{P}_N(0;\alpha ,\beta ,\gamma )`$ $`=`$ $`{\displaystyle \underset{l=1}{\overset{3}{}}}\left(\right|{\displaystyle \underset{m=0}{\overset{(N1)/2}{}}}c_{1,m,l}(N)|^2`$ (23)
$`+|{\displaystyle \underset{j=2}{\overset{3}{}}}c_{j,0,l}(N)|^2+{\displaystyle \underset{j=2}{\overset{3}{}}}{\displaystyle \underset{m=1}{\overset{(N1)/2}{}}}|c_{j,m,l}(N)|^2).`$
The difference between the first term and the other terms is caused by the difference of the degree of degenerate eigenvalues.
Let $`\varphi _{j,k}(N)`$ be the eigenvectors corresponding to the eigenvalue $`e^{i\theta _{j,k}t}`$ of the time-evolution matrix of the 3QW with a matrix size $`3N\times 3N`$. They are easily obtained from $`|\mathrm{\Phi }_k^j`$ by
$`|\varphi _m^j(N)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3N}}}[|\mathrm{\Phi }_m^j,\omega |\mathrm{\Phi }_m^j,\omega ^2|\mathrm{\Phi }_m^j,\mathrm{},\omega ^{N1}|\mathrm{\Phi }_m^j],`$
where $`\omega =e^{2\pi i/N}`$. Since the coefficients $`c_{j,m,l}(N)`$ are proportion to the product of eigenvectors, the orders with respect to $`N`$ of the first term and the second term in Eq. (23) are $`𝒪`$$`(1)`$ and $`𝒪(N^\mathit{1})`$, respectively. Thus we can neglect the second term in the limit of $`N\mathrm{}`$. Using the eigenvectors in Eq. (18), we have
$`\overline{P}_{\mathrm{}}(0;\alpha ,\beta ,\gamma )`$ $`=`$ $`(52\sqrt{6})(1+|\alpha +\beta |^2+|\beta +\gamma |^22|\beta |^2).`$
The time-averaged probability takes the maximum value $`2(52\sqrt{6})`$ at $`\beta =0`$, which is the value indicated by the horizontal dashed line in Fig. 1. The component of $`\overline{P}_{\mathrm{}}(0;\alpha ,\beta ,\gamma )`$ corresponding to $`l=1,2,3`$ are respectively given by
$`\overline{P}_{\mathrm{}}(0;1;\alpha ,\beta ,\gamma )`$ $`=`$ $`{\displaystyle \frac{|\sqrt{6}\alpha 2(\sqrt{6}3)\beta +(125\sqrt{6})\gamma |^2}{36}},`$
$`\overline{P}_{\mathrm{}}(0;2;\alpha ,\beta ,\gamma )`$ $`=`$ $`{\displaystyle \frac{(\sqrt{6}3)^2|\alpha +\beta +\gamma |^2}{9}},`$
$`\overline{P}_{\mathrm{}}(0;3;\alpha ,\beta ,\gamma )`$ $`=`$ $`{\displaystyle \frac{|\sqrt{6}\gamma 2(\sqrt{6}3)\beta +(125\sqrt{6})\alpha |^2}{36}}.`$ (24)
We stress here that the time-averaged probability is not always positive. If $`\alpha =1/\sqrt{6}`$, $`\beta =2/\sqrt{6}`$ and $`\gamma =1/\sqrt{6},`$ then the time-averaged probability becomes zero, that is, $`\overline{P}_{\mathrm{}}(0;1/\sqrt{6},2/\sqrt{6},1/\sqrt{6})=0.`$
## IV Stationary distribution of the particle
We showed that the time-averaged probability of 3QW converges to non-zero values except special initial states. This result, however, does not mean that a particle is observed with almost the same probability at the origin after long time evolution. Indeed the probabilities of finding a particle in Grover walk and the four-state quantum walk, whose time-averaged probability converge to non-zero values, do not converge. Because the time-averaged probabilities in both cases depend on the parity of the time. Figure 1 suggests that the probability $`P(n,t)=P(n,t;\alpha ,\beta ,\gamma )`$ itself converges in a limit of $`t\mathrm{}`$. In this section we calculate the limit $`P_{}(n)lim_t\mathrm{}P(n,t)`$ rigorously and consider the dependence of $`P_{}(n)`$ on the position $`n`$.
The wavefuntion given by Eq. (19) is infinite superimposition of the wavefunction $`e^{i(\theta _{j,k}t+kn)}`$. If the argument $`\theta _{j,k}`$ given in Eq. (14) for $`j=2,3`$ is not absolute zero, then the $`e^{i\theta _{j,k}t}`$ oscillates with high frequency in $`k`$-space for large $`t`$. On the other hand, $`c_k(\theta _{j,k})\phi _k(\theta _{j,k},l)`$ in Eq. (20) changes smoothly with respect to $`k`$. Therefore we expect that the integration for $`j=2,3`$ in Eq. (20) becomes small as the time increases due to cancellation and converges to zero in the limit of $`t\mathrm{}`$. That is,
$`\underset{t\mathrm{}}{lim}{\displaystyle \underset{j=2}{\overset{3}{}}}\mathrm{\Psi }_j(n,t;l;\alpha ,\beta ,\gamma )=0,`$ (25)
for $`l=1,2,3`$ and $`\alpha ,\beta ,\gamma `$ with $`|\alpha |^2+|\beta |^2+|\gamma |^2=1`$. This conjecture can be rigorously proved using the Riemann-Lebesgue lemma (see Appendix A). Consequently the probability $`P_{}(n)=P_{}(n;\alpha ,\beta ,\gamma )`$ is determined from the eigienvectors corresponding to the eigenvalue “1”, that is, $`\theta _{1,k}=0`$ (see Eq. (14)), and $`l`$-th component of $`P_{}(n;l)=P_{}(n;l;\alpha ,\beta ,\gamma )`$ is given by
$`P_{}(n;l;\alpha ,\beta ,\gamma )=\left|\mathrm{\Psi }_1(n,t;l;\alpha ,\beta ,\gamma )\right|^2.`$ (26)
Note that $`\mathrm{\Psi }_1(n,t;l;\alpha ,\beta ,\gamma )`$ does not depend on time $`t`$, since $`\theta _{1,k}=0`$. By transforming integration in the left-hand side in Eq. (26) into complex integral, we have
$`P_{}(n;1;\alpha ,\beta ,\gamma )`$ $`=`$ $`\left|2\alpha I(n)+\beta J_+(n)+2\gamma K_+(n)\right|^2,`$
$`P_{}(n;2;\alpha ,\beta ,\gamma )`$ $`=`$ $`\left|\alpha J_{}(n)+{\displaystyle \frac{\beta }{2}}L(n)+\gamma J_+(n)\right|^2,`$
$`P_{}(n;3;\alpha ,\beta ,\gamma )`$ $`=`$ $`\left|2\alpha K_{}(n)+\beta J_{}(n)+2\gamma I(n)\right|^2,`$ (27)
where $`c=5+2\sqrt{6}((1,0))`$ and
$`I(n)={\displaystyle \frac{2c^{|n|+1}}{c^21}},L(n)=I(n1)+2I(n)+I(n+1),`$
$`J_+(n)=I(n)+I(n+1),J_{}(n)=I(n1)+I(n),`$
$`K_+(n)=I(n+1),K_{}(n)=I(n1),`$ (28)
for any $`n.`$ We should remark that it is confirmed $`\overline{P}_{\mathrm{}}(0;l;\alpha ,\beta ,\gamma )=P_{}(0;l;\alpha ,\beta ,\gamma )`$ for $`l=1,2,3`$ by using Eqs. (24), (27) and (28).
Here we give an example. From Eqs. (27) and (28), we obtain
$`P_{}(0;i/\sqrt{2},0,1/\sqrt{2})={\displaystyle \frac{4c^2(5c^2+2c+5)}{(1c^2)^2}}=104\sqrt{6}=0.202\mathrm{},`$ (29)
$`P_{}(|n|;i/\sqrt{2},0,1/\sqrt{2})={\displaystyle \frac{2(5c^4+2c^3+10c^2+2c+5)}{(1c^2)^2}}c^{2|n|},`$ (30)
for any $`n1`$. Furthermore,
$`0<{\displaystyle \underset{n}{}}P_{}(n;i/\sqrt{2},0,1/\sqrt{2})=1/\sqrt{6}=0.408\mathrm{}<1.`$ (31)
That is, $`P_{}(n;i/\sqrt{2},0,1/\sqrt{2})`$ is not a probability measure. The above value depends on the initial qubit state, for example,
$`{\displaystyle \underset{n}{}}P_{}(n;1/\sqrt{3},1/\sqrt{3},1/\sqrt{3})=3\sqrt{6}=0.550\mathrm{},`$
$`{\displaystyle \underset{n}{}}P_{}(n;1/\sqrt{3},1/\sqrt{3},1/\sqrt{3})=(3\sqrt{6})/9=0.061\mathrm{}.`$
We should remark that in the case of the classical symmetric random walk starting from the origin, it is known that $`P_{}(n)=0`$ for any $`n,`$ therefore we have $`_nP_{}(n)=0.`$ The same conclusion can be obtained for the discrete-time and continuous-time two-state quantum walks Kon2002 ; Kon2005 ; Kon2005b .
Figure 3 shows the probability $`P_{}(n;i/\sqrt{2},0,1/\sqrt{2})`$ in log scale. The probability decreases exponentially for large $`|n|`$ and its asymptotic behavior is express by $`P_{}(n;i/\sqrt{2},0,1/\sqrt{2})c^{2|n|}`$ in the limit of $`n\pm \mathrm{}`$ as Eq. (30) indicates.
We here mention the time-dependence of the $`P(n,t)`$. The particle is observed near the origin with high probability. This, however, dose not imply that the particle can not escape form the region near the origin. Figure 4 shows the change of $`P(n,t)`$ in a space-time. The particle is observed with high-probability at the dark region. One clearly finds three dark regions. The first is a region near a center line connecting with the origin, and we can confirm the localization near the origin. The second and third regions are boundaries of the triangle. These regions show trajectories of two peaks moving outside in the space-time, which are also seen in the Hadamard walk. As a result we conclude that the particle in 3QW splits three parts in superposition.
## V Conclusions and discussions
The unique properties which are not observed in other quantum walks were found in three-state quantum walk. The particle which exists at the origin splits three pieces in superposition. Two of three leave for infinite points and the remainder stays at the origin. In contrast with the ordinary Hadamard walk, the probability of finding the particle at the origin does not vanish for large time, and its maximum probability is $`104\sqrt{6}=0.202\mathrm{}.`$
A simple reason why the 3QW is different form other quantum walks is the difference in the degree of degenerate eigenvalues. The necessary condition of the localization is the existence of the degenerate eigenvalues. And furthermore, the each degree of degeneration must be proportion to the dimension of the Hilbert space. In addition, if the degenerate eigienvalue is 1 only, then the probability of finding the particle can converges in the limit of $`t\mathrm{}`$. The Grover walk and the four-state quantum walk exhibit the localization, but the probability of finding the particle oscillates. Because there are degenerate eigenvalues with values “1” and “-1”. On the other hand, the degenerated eigenvalue in 3QW unrelated to the wave number is only “1”, therefore the probability converges. Although no experiment exists about the quantum walks yet, the stationary properties of 3QW may be advantage to comparison with theoretical results.
Finally we discuss a relation between the limit distribution for the original 3QW $`X_t`$ as time $`t\mathrm{}`$ and that of the rescaled $`X_t/t`$ in the same limit. When we consider the 3QW starting from a mixture of three pure states $`[1,0,0]^T,[0,1,0]^T,`$ and $`[0,0,1]^T`$ with probability $`1/3`$ respectively, we can obtain a weak limit probability distribution $`f(x)`$ for the rescaled 3QW $`X_t/t`$ as $`t\mathrm{}`$ in the following:
$`f(x)={\displaystyle \frac{1}{3}}\delta _0(x)+{\displaystyle \frac{\sqrt{8}I_{(1/\sqrt{3},1/\sqrt{3})}(x)}{3\pi (1x^2)\sqrt{13x^2}}},`$ (32)
for $`x[1,1]`$, where $`\delta _0(x)`$ denotes the pointmass at the origin and $`I_{(a,b)}(x)=1,`$ if $`x(a,b),=0,`$ otherwise. The above derivation is due to the method by Grimmett et al. GJS . We should note that the first term in the right-hand side of Eq. (32) corresponds to the localization for the original 3QW. In fact, from Eq. (28), we have
$`{\displaystyle \frac{1}{3}}{\displaystyle \underset{n}{}}\left[P_{}(n;1,0,0)+P_{}(n;0,1,0)+P_{}(n;0,0,1)\right]={\displaystyle \frac{1}{3}}.`$
The last value $`1/3`$ is nothing but the coefficient $`1/3`$ of $`\delta _0(x).`$ Moreover, the second term in the right-hand side of Eq. (32) has a similar form a weak limit density function for the same rescaled Hadamard walk with two inner states starting from a uniform random mixture of two pure states $`[1,0]^T`$ and $`[0,1]^T`$:
$`f_H(x)={\displaystyle \frac{I_{(1/\sqrt{2},1/\sqrt{2})}(x)}{\pi (1x^2)\sqrt{12x^2}}},`$
for $`x[1,1]`$, see Kon2002 ; Kon2005 ; GJS . This case does not have a delta measure term corresponding to a localization. More detailed study on this line will appear in our forthcoming paper.
## Appendix A Proof of Eq. (25)
Here we give an outline of the proof of Eq. (25). A direct computation gives
$`{\displaystyle \underset{j=2}{\overset{3}{}}}\left[\begin{array}{c}\mathrm{\Psi }_j(n,t;1;\alpha ,\beta ,\gamma )\\ \mathrm{\Psi }_j(n,t;2;\alpha ,\beta ,\gamma )\\ \mathrm{\Psi }_j(n,t;3;\alpha ,\beta ,\gamma )\end{array}\right]=M\left[\begin{array}{c}\alpha \\ \beta \\ \gamma \end{array}\right],`$
where $`M=(m_{ij})_{1i,j3}`$ with
$`m_{11}=3J_{n,t}+{\displaystyle \frac{1}{2}}\left\{J_{n1,t}+J_{n+1,t}+\left(K_{n1,t}K_{n+1,t}\right)\right\},`$
$`m_{33}=3J_{n,t}+{\displaystyle \frac{1}{2}}\left\{J_{n1,t}+J_{n+1,t}\left(K_{n1,t}K_{n+1,t}\right)\right\},`$
$`m_{12}=\left\{J_{n,t}+J_{n+1,t}+\left(K_{n,t}K_{n+1,t}\right)\right\},`$
$`m_{32}=\left\{J_{n,t}+J_{n1,t}+\left(K_{n,t}K_{n1,t}\right)\right\},`$
$`m_{13}=2J_{n+1,t},m_{31}=2J_{n1,t},`$
$`m_{21}=\left\{J_{n,t}+J_{n1,t}+\left(K_{n1,t}K_{n,t}\right)\right\},`$
$`m_{23}=\left\{J_{n,t}+J_{n+1,t}+\left(K_{n+1,t}K_{n,t}\right)\right\},`$
$`m_{22}=4J_{n,t},`$
and
$`J_{n,t}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _\pi ^\pi }{\displaystyle \frac{\mathrm{cos}(kn)}{5+\mathrm{cos}k}}\mathrm{cos}(\theta _kt)𝑑k,`$
$`K_{n,t}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _\pi ^\pi }{\displaystyle \frac{\mathrm{cos}(kn)}{\sqrt{(5+\mathrm{cos}k)(1\mathrm{cos}k)}}}\mathrm{sin}(\theta _kt)𝑑k.`$
From the Riemann-Lebesgue lemma, we can show that
$`\underset{t\mathrm{}}{lim}J_{n,t}=0,\underset{t\mathrm{}}{lim}(K_{n,t}K_{n+1,t})=0,`$
for any $`n`$. Therefore we have the desired conclusion. |
warning/0507/cond-mat0507524.html | ar5iv | text | # Response of thin-film SQUIDs to applied fields and vortex fields: Linear SQUIDs
## I Introduction
The research in this paper has been motivated by several important recent developments in superconductivity: (a) the fabrication of thin-film SQUIDs (superconducting quantum interference devices) made of high-$`T_c`$ superconductors, Koelle99 (b) the study of noise generated by vortices in active and passive superconducting devices, Ferrari91 ; Miklich94 ; Glyantsev96 ; Koelle99 ; Humphreys99 ; Straub01 ; Woerdenweber02 ; Doenitz04 ; Nowak05 ; Kuriki04 and (c) line-width reduction in superconducting devices to eliminate noise due to vortices trapped during cooldown in the earth’s magnetic field.ClemUnpub ; Dantsker96 ; Dantsker97 ; Mitchell03 ; Stan04 ; Du04
The fact that the London penetration depth $`\lambda `$ increases as $`T`$ increases and diverges at $`T_c`$ is an important consideration for high-$`T_c`$ SQUIDs operated at liquid-nitrogen temperature. When $`\lambda `$ is larger than the film thickness $`d`$, the physical length that enters the equations governing the spatial variation of currents and fields is the Pearl lengthPearl64 $`\mathrm{\Lambda }=\lambda ^2/d`$. Accordingly, the equations governing the behavior of active or passive thin-film superconducting devices depend upon the ratio of $`\mathrm{\Lambda }`$ to the linear dimensions of the device. In particular, the equations governing SQUIDs involve not just the magnetic flux up through a contour within the SQUID, but the effective flux $`\mathrm{\Phi }`$, which is the sum of the magnetic flux and a term proportional to the line integral of the current density around the same contour. While the effective flux $`\mathrm{\Phi }`$ is similar to London’s fluxoid, London61 which is quantized in multiples of the superconducting flux quantum $`\varphi _0=h/2e`$, we show in the next section that $`\mathrm{\Phi }`$ is not quantized. We also give in Sec. II the basic equations, valid for any value of $`\mathrm{\Lambda }`$, that govern the behavior of a dc SQUID.
A vortex trapped in the body of the SQUID during cooldown through the superconducting transition temperature $`T_c`$ in an ambient magnetic field generates a magnetic field and a screening current that together make a sizable vortex-position-dependent contribution $`\mathrm{\Phi }_\mathrm{v}`$ to the effective flux $`\mathrm{\Phi }`$. If such a vortex remains fixed in position and the temperature remains constant, this simply produces a harmless bias in $`\mathrm{\Phi }`$. On the other hand, both vortex motion due to thermal agitation and temperature fluctuations generate corresponding fluctuations in $`\mathrm{\Phi }_\mathrm{v}`$ and noise in the SQUID output. In Sec. II we show that to calculate $`\mathrm{\Phi }_\mathrm{v}`$, it is not necessary to calculate the spatial dependence of the vortex-generated fields and currents. Instead, one may determine $`\mathrm{\Phi }_\mathrm{v}`$ with the help of the sheet-current distribution of a circulating current in the absence of the vortex.
In Sec. III we apply the basic equations of Sec. II to calculate the properties of a model linear SQUID, which has the basic topology of a SQUID but is greatly stretched along one axis, such that the central portion resembles a coplanar stripline. The advantage of using such a model is that simple analytical results can be derived that closely approximate the exact numerically calculated quantities in the appropriate limits. In addition to calculating the field and current distributions for several values of $`\mathrm{\Lambda }`$, we calculate the total inductance, geometric inductance, kinetic inductance, and magnetic moment when the SQUID carries a circulating current. We calculate the field and current distributions, magnetic moment, and effective area $`A_{\mathrm{eff}}=\mathrm{\Phi }_\mathrm{f}/B_\mathrm{a}`$ when a perpendicular magnetic induction $`B_\mathrm{a}`$ is applied and the effective flux $`\mathrm{\Phi }_\mathrm{f}`$ is focused into the SQUID. Finally, we calculate the field and current distributions and the magnetic moment for the zero-fluxoid state when the junctions are short-circuited and the sample remains in the state with $`\mathrm{\Phi }=0`$ when a perpendicular magnetic induction $`B_\mathrm{a}`$ is applied.
In Sec. IV, we present a brief summary of our results.
## II Basic equations
Our purpose in this section is to derive general equations that govern the behavior of a dc SQUID consisting of thin superconducting films of thickness $`d`$ less than the weak-field London penetration depth $`\lambda `$, such that the fields and currents are governed by the two-dimensional screening length or Pearl length Pearl64 $`\mathrm{\Lambda }=\lambda ^2/d`$. We calculate both the current $`I=I_1+I_2`$ through the SQUID \[see Fig. 1\] and the circulating currentfoot1 $`I_\mathrm{d}=(I_2I_1)/2`$ and describe how to calculate the critical current $`I_\mathrm{c}`$ of the SQUID for arbitrary values of the SQUID’s inductance $`L`$. For this case, the contributions from line integrals of the current density to the effective flux in the hole cannot be neglected, and the kinetic inductance makes a significant contribution to $`L`$. When a perpendicular magnetic induction $`B_\mathrm{a}`$ is applied, we calculate how much magnetic flux is focused into the SQUID’s hole; this flux also can be expressed in terms of the effective areaKetchen85 of the hole. We also show how to calculate how much magnetic flux generated by a vortex in the main body of the SQUID is focused into the hole.
Consider a dc SQUID in the $`xy`$ plane, as sketched in Fig. 1. We suppose that the SQUID is symmetric about the $`y`$ axis, which lies along the centerline. The maximum Josephson critical current is $`I_0`$ for each of the Josephson junctions, shown as small black squares.
The currents up through the left and right sides of the SQUID can be written as $`I_1=I/2I_\mathrm{d}`$ and $`I_2=I/2+I_\mathrm{d}`$. When the magnitude of $`I`$, the total current through the SQUID, is less than the critical current $`I_\mathrm{c}`$, the equations that determine $`I=I_1+I_2`$ and the circulating current $`I_\mathrm{d}=(I_2I_1)/2`$ can be derived using a method similar to that used in Ref. Orlando91, . We begin by writing the local current density $`𝒋`$ in the superconductors (i.e., the main body of the SQUID and the counterelectrode) asLondon61
$$𝒋=(1/\mu _0\lambda ^2)[𝑨+(\varphi _0/2\pi )\gamma ],$$
(1)
where $`𝑨`$ is the vector potential, $`\varphi _0=h/2e`$ is the superconducting flux quantum, and $`\gamma `$ is the phase of the order parameter. The quantity inside the brackets, which is gauge-invariant, can be thought of as the superfluid velocity expressed in units of vector potential. From the point of view of the Ginzburg-Landau theory, implicit in the use of this London-equation approach is the assumption that the applied fields and currents are so low that the magnitude of the order parameter is not significantly reduced from its equilibrium value in the absence of fields and currents.
To obtain the SQUID equations, we integrate the vector potential around a contour $`C`$ that passes in a counterclockwise direction through both junctions, the main body of the SQUID, and the counterelectrode as shown in Fig. 1 and write the result in two ways. Since $`𝑩=\times 𝑨`$, this integral yields, on the one hand, the magnetic flux in the $`z`$ direction
$$_SB_z(x,y)𝑑S,$$
(2)
where $`S`$ is the area surrounded by the contour $`C`$ and $`B_z(x,y)`$ is the $`z`$ component of the net magnetic induction in the plane of the SQUID produced by the sum of a perpendicular applied field $`B_\mathrm{a}`$ and the self-field $`B_{sz}(x,y)`$ generated via the Biot-Savart law by the supercurrent density $`𝒋(x,y)`$. On the other hand, for those portions of the contour lying in the superconductors, we eliminate the line integrals of $`𝑨`$ in favor of line integrals of $`𝒋`$ using Eq. (1). We then express the line integrals of $`\gamma `$ in terms of the values of $`\gamma `$ at the junctions. Equating the two expressions for the line integral of $`𝑨`$, we find that the effective flux $`\mathrm{\Phi }`$ in the $`z`$ direction through the SQUID is given by
$$\mathrm{\Phi }=(\varphi _0/2\pi )(\varphi _1\varphi _2),$$
(3)
where
$$\mathrm{\Phi }=_SB_z(x,y)𝑑S+\mu _0\lambda ^2_C𝒋𝒅𝒍,$$
(4)
with the integration contour $`C`$ now passing through both superconductors but excluding the junction barriers. These equations are equivalent to Eq. (8.67) in Ref. Orlando91, . The gauge-invariant phase differences across the junctions $`b`$ and $`a`$ are, respectively,Josephson69
$`\varphi _1=\gamma _{bc}\gamma _{bs}(2\pi /\varphi _0){\displaystyle _{bc}^{bs}}𝑨𝒅𝒍,`$ (5)
$`\varphi _2=\gamma _{ac}\gamma _{as}(2\pi /\varphi _0){\displaystyle _{ac}^{as}}𝑨𝒅𝒍,`$ (6)
where $`bc`$ labels a point on the counterelectrode side of the junction $`b`$ and $`bs`$ labels the point directly across the insulator in the SQUID washer, and $`ac`$ and $`as`$ label corresponding points for junction $`a`$. According to the Josephson equations,Josephson69 the junction supercurrents are $`I_1=I_0\mathrm{sin}\varphi _1`$ and $`I_2=I_0\mathrm{sin}\varphi _2`$. In the above derivation we have assumed that the linear dimensions of the Josephson junctions are much less than the Josephson penetration depth $`\lambda _J`$,Josephson69 and that the applied fields are sufficiently small that the Josephson current densities and gauge-invariant phase differences are very nearly constant across the junction areas.
The magnetic moment $`𝒎=m\widehat{z}`$ generated by the currents in the SQUID isJackson62
$$𝒎=\frac{1}{2}𝒓\times 𝒋d^3r.$$
(7)
It can be shown with the help of the London fluxoid quantization condition,London61
$$_S^{}B_z(x,y)𝑑S+\mu _0\lambda ^2_C^{}𝒋𝒅𝒍=n\varphi _0,$$
(8)
where $`n`$ is an integer and $`C^{}`$ is a closed contour that surrounds an area $`S^{}`$ within the body of the SQUID, that if there are no vortices present (i.e., when $`n=0`$), the expression for $`\mathrm{\Phi }`$ in Eq. (4) is independent of the choice of contour $`C`$. Any convenient path can be chosen for $`C`$, provided only that the path remains in the superconducting material in the body of the SQUID and the counterelectrode. On the other hand, when there are vortices in the main body of the SQUID, the quantity $`\mathrm{\Phi }`$ increases by $`\varphi _0`$ each time the contour $`C`$ is moved from a path inside the vortex axis to one enclosing the vortex axis. Thus, without specifying the precise contour $`C`$, $`\mathrm{\Phi }`$ is determined only modulo $`\varphi _0`$. However, this is of no physical consequence, because the gauge-invariant phases $`\varphi _1`$ and $`\varphi _2`$, which also enter Eq. (3), are also determined only modulo $`2\pi `$. The final equations determining the currents $`I`$ and $`I_\mathrm{d}`$ are independent of the choice of contour $`C`$ and remain valid even when vortices are present in the main body of the SQUID.
When the thickness $`d`$ of the SQUID is much larger than $`\lambda `$, the contours $`C`$ and $`C^{}`$ can be chosen to be at the midpoint of the thickness, where $`𝒋`$ is exponentially small, such that the line integrals of $`𝒋`$ can be neglected. The resulting equations are then the familiar ones found in many reference books, such as Refs. Orlando91, ; deBruyn70, ; Solymar72, ; deBruyn77, ; VanDuzer81, ; Clarke90, ; Tinkham96, . However, we are interested here in the case for which $`d<\lambda `$, such that the fields and currents are governed by the two-dimensional screening length or Pearl length Pearl64 $`\mathrm{\Lambda }=\lambda ^2/d`$. The term in Eq. (4) involving $`𝒋`$ then must be carefully accounted for. For this case, $`𝒋`$ is very nearly constant over the thickness and it is more convenient to deal with the sheet-current density $`𝑱(x,y)=𝒋d`$, such that Eqs. (4) and (8) take the formfoot4
$$\mathrm{\Phi }=_SB_z(x,y)𝑑S+\mu _0\mathrm{\Lambda }_C𝑱𝒅𝒍$$
(9)
and
$$_S^{}B_z(x,y)𝑑S+\mu _0\mathrm{\Lambda }_C^{}𝑱𝒅𝒍=n\varphi _0.$$
(10)
For the general case when the SQUID is subject to a perpendicular applied magnetic induction $`B_\mathrm{a}`$, carries a current $`I`$ unequally divided between the two arms, $`I_1=I/2I_\mathrm{d}`$ and $`I_2=I/2+I_\mathrm{d}`$, where the circulating currentfoot1 is $`I_\mathrm{d}=(I_2I_1)/2`$, and contains a vortex at the position $`𝒓_\mathrm{v}`$ in the body of the SQUID, the effective flux $`\mathrm{\Phi }`$ in the $`z`$ direction can be written as the sum of four independent contributions:
$$\mathrm{\Phi }=\mathrm{\Phi }_I+\mathrm{\Phi }_\mathrm{d}+\mathrm{\Phi }_\mathrm{f}+\mathrm{\Phi }_\mathrm{v},$$
(11)
The first term on the right-hand side of Eq. (11) is that which would be produced by equal currents $`I/2`$ in the $`y`$ direction on the left and right sides of the SQUID shown in Fig. 1:
$$\mathrm{\Phi }_I=_SB_I(x,y)𝑑S+\mu _0\mathrm{\Lambda }_C𝑱_I𝒅𝒍,$$
(12)
where $`B_I(x,y)`$ is the $`z`$ component of the self-field generated via the Biot-Savart law by the sheet-current density $`𝑱_I(x,y)`$, subject to the condition that the same current $`I/2`$ flows through the two contacts $`a`$ and $`b`$. For a symmetric SQUID, $`J_I(x,y)`$, the $`y`$ component of $`𝑱_I(x,y)`$, is then an even function of $`x`$, and $`J_I(x,y)`$ and $`B_I(x,y)`$ are odd functions of $`x`$. As a result, both terms on the right-hand side of Eq. (12) vanish by symmetry, and $`\mathrm{\Phi }_I=0`$. Since $`𝑱_I=0`$ except at the contacts $`a`$ and $`b`$, we may write $`𝑱_I=(I/2)\times 𝑮_I`$, where $`𝑮_I=\widehat{z}G_I,`$ such that $`𝑱_I(x,y)=(I/2)\widehat{z}\times G_I(x,y)`$. The contours of the scalar stream function $`G_I(x,y)=`$ const correspond to streamlines of $`𝑱_I(x,y)`$, and we may choose $`G_I=0`$ for points $`𝒓_i=(x_i,y_i)`$ all along the inner edges of the superconductors and $`G_I=1`$ for points $`𝒓_o=(x_o,y_o)`$ all along the outer right edges and $`G_I=1`$ for points $`𝒓_o=(x_o,y_o)`$ all along the outer left edges.
The second term on the right-hand side of Eq. (11) is due to the circulating currentfoot1 $`I_\mathrm{d}=(I_2I_1)/2`$ in the counterclockwise direction when unequal currents flow in the two sides of the SQUID shown in Fig. 1:
$$\mathrm{\Phi }_\mathrm{d}=_SB_\mathrm{d}(x,y)𝑑S+\mu _0\mathrm{\Lambda }_C𝑱_\mathrm{d}𝒅𝒍,$$
(13)
where $`B_\mathrm{d}(x,y)`$ is the $`z`$ component of the self-field generated via the Biot-Savart law by the circulating sheet-current density $`𝑱_\mathrm{d}(x,y)`$ when a current $`I_\mathrm{d}`$ flows through contact $`a`$ from the counterelectrode into the body of the SQUID, passes around the central hole, and flows through contact $`b`$ back into the counterelectrode. The magnetic moment $`m_d`$ generated by the circulating current is proportional to $`I_d`$, as can be seen from Eq. (7). Since $`𝑱_\mathrm{d}=0`$ except at the contacts $`a`$ and $`b`$, we may write $`𝑱_\mathrm{d}=I_\mathrm{d}\times 𝑮_\mathrm{d}`$, where $`𝑮_\mathrm{d}=\widehat{z}G_\mathrm{d}`$, such that $`𝑱_\mathrm{d}(x,y)=I_\mathrm{d}\widehat{z}\times G_\mathrm{d}(x,y)`$. The contours of the scalar stream function $`G_\mathrm{d}(x,y)=`$ const correspond to streamlines of $`𝑱_\mathrm{d}(x,y)`$, and we may choose $`G_\mathrm{d}=0`$ for points $`𝒓_i=(x_i,y_i)`$ all along the inner edges of the superconductors and $`G_\mathrm{d}=1`$ for points $`𝒓_o=(x_o,y_o)`$ all along the outer edges. Once a numerical result for $`\mathrm{\Phi }_\mathrm{d}`$ is found, the result can be used to determine the inductance $`L`$ of the SQUID via $`L=\mathrm{\Phi }_\mathrm{d}/I_\mathrm{d}`$, as was done for a circular ring in Ref. Brandt04, . The resulting expression for $`L`$ is the sum of the geometric and kinetic inductances.
The third term on the right-hand side is a flux-focusing term due to the applied field:
$$\mathrm{\Phi }_\mathrm{f}=_SB_\mathrm{f}(x,y)𝑑S+\mu _0\mathrm{\Lambda }_C𝑱_\mathrm{f}𝒅𝒍,$$
(14)
where $`B_\mathrm{f}(x,y)`$ is the $`z`$ component of the net magnetic induction in the plane of the SQUID produced by the sum of a perpendicular applied field $`B_\mathrm{a}`$ and the $`z`$ component of the self-field $`B_{s\mathrm{f}}(x,y)`$ generated via the Biot-Savart law by the sheet-current density $`𝑱_\mathrm{f}(x,y)`$ induced in response to $`B_\mathrm{a}`$, subject to the condition that no current flows through the junctions $`a`$ and $`b`$. In other words, the desired fields are those that would appear in response to $`B_\mathrm{a}`$ if the junctions $`a`$ and $`b`$ were open-circuited. Since $`𝑱_\mathrm{f}=0`$, we may write $`𝑱_\mathrm{f}=\times 𝑮_\mathrm{f}`$, where $`𝑮_\mathrm{f}=\widehat{z}G_\mathrm{f}`$, such that $`𝑱_\mathrm{f}(x,y)=\widehat{z}\times G_\mathrm{f}(x,y)`$. The contours of the scalar stream function $`G_\mathrm{f}(x,y)=`$ const correspond to streamlines of $`𝑱_\mathrm{f}(x,y)`$, and we may chose $`G_\mathrm{f}=0`$ for all points $`(x,y)`$ along the inner and outer edges of the superconductor. Once a numerical result for $`\mathrm{\Phi }_\mathrm{f}`$ is found, the result can be used to determine the effective areaKetchen85 of the SQUID’s central hole, $`A_{\mathrm{eff}}=\mathrm{\Phi }_\mathrm{f}/B_\mathrm{a}`$, as was done for a circular ring in Ref. Brandt04, .
To prove that the effective area is also given by $`A_{\mathrm{eff}}=m_d/I_d`$,foot2 we consider the electromagnetic energy cross term $`E_{\mathrm{fd}}=(𝑩_\mathrm{f}𝑩_\mathrm{d}/\mu _0+\mu _0\lambda ^2𝒋_\mathrm{f}𝒋_\mathrm{d})d^3r`$, where the integral extends over all space. Here, $`𝑩_\mathrm{f}(𝒓)=𝑩_\mathrm{a}(𝒓)+𝑩_{s\mathrm{f}}(𝒓)=\times 𝑨_\mathrm{f}(𝒓),`$ where $`𝒋_\mathrm{a}(𝒓)=\times 𝑩_\mathrm{a}(𝒓)/\mu _0`$ is the current density in the distant coil that produces a nearly uniform field $`B_a`$ in the vicinity of the SQUID, $`𝒋_\mathrm{f}=\times 𝑩_{s\mathrm{f}}/\mu _0`$ is the induced current density in the SQUID, and $`𝑩_{s\mathrm{f}}`$ is the corresponding self-field under the conditions of flux focusing, i.e., when $`𝒋_\mathrm{f}=0`$ through the junctions. Also, $`𝑩_\mathrm{d}=\times 𝑨_\mathrm{d}`$ is the dipole-like field distribution generated by the circulating current $`I_\mathrm{d}`$ with density $`𝒋_\mathrm{d}`$ in the SQUID; at large distances from the SQUIDJackson62 $`𝑨_\mathrm{d}=\mu _0𝒎_\mathrm{d}\times 𝒓/4\pi r^3`$. We evaluate $`E_{\mathrm{fd}}`$ in two ways, making use of the vector identities $`(𝑨\times 𝑩)=𝑩(\times 𝑨)𝑨(\times 𝑩)`$ and $`(\gamma 𝒋)=\gamma 𝒋+\gamma 𝒋`$, and applying the divergence theorem, first with $`𝑨=𝑨_\mathrm{d}`$, $`𝑩=𝑩_\mathrm{f}`$, $`\gamma =\gamma _\mathrm{d}`$, and $`𝒋=𝒋_\mathrm{f}`$, from which we obtain $`E_{\mathrm{fd}}=B_\mathrm{a}m_\mathrm{d}`$, and then with $`𝑨=𝑨_\mathrm{f}`$, $`𝑩=𝑩_\mathrm{d}`$, $`\gamma =\gamma _\mathrm{f}`$, and $`𝒋=𝒋_\mathrm{d}`$, from which we obtain $`E_{\mathrm{fd}}=\mathrm{\Phi }_\mathrm{f}I_\mathrm{d}`$ with the help of Eq. (3). Since $`\mathrm{\Phi }_\mathrm{f}=B_\mathrm{a}A_{\mathrm{eff}}`$, the effective area obeys $`A_{\mathrm{eff}}=m_d/I_d`$.
The fourth term on the right-hand side of Eq. (11) is due to a vortex at position $`𝒓_\mathrm{v}=\widehat{x}x_\mathrm{v}+\widehat{y}y_\mathrm{v}`$ in the body of the SQUID:
$$\mathrm{\Phi }_\mathrm{v}(𝒓_\mathrm{v})=_SB_\mathrm{v}(x,y)𝑑S+\mu _0\mathrm{\Lambda }_C𝑱_\mathrm{v}𝒅𝒍,$$
(15)
where $`B_\mathrm{v}(x,y)`$ is the $`z`$ component of the self-field generated by the vortex’s sheet-current density $`𝑱_\mathrm{v}(x,y)`$ via the Biot-Savart law when no current flows through the junctions $`a`$ and $`b`$. The desired fields are those that would appear in response to the vortex if the junctions $`a`$ and $`b`$ were open-circuited. Since $`𝑱_\mathrm{v}=0`$, it is possible to express $`𝑱_\mathrm{v}(x,y)`$ in terms of a scalar stream function, as we did for $`𝑱_\mathrm{f}(x,y)`$ and $`𝑱_\mathrm{d}(x,y)`$. However, as shown below, it is possible to use energy arguments to express $`\mathrm{\Phi }_\mathrm{v}(x,y)`$ in terms of the stream function $`G_\mathrm{d}(x,y)`$.foot3
To obtain $`\mathrm{\Phi }_\mathrm{v}(𝒓)`$ when a vortex is at the position $`𝒓=\widehat{x}x+\widehat{y}y`$, imagine disconnecting the counterelectrode in Fig. 1 and attaching leads from a power supply to the contacts $`a`$ and $`b`$. The power supply provides a constant current $`I_\mathrm{d}`$ in the counterclockwise direction, and the sheet-current distribution through the body of the SQUID is given by $`𝑱_\mathrm{d}(x,y)=I_\mathrm{d}\widehat{z}\times G_\mathrm{d}(x,y)`$, as discussed above. We also imagine attaching leads from a high-impedance voltmeter to the contacts $`a`$ and $`b`$. If the vortex moves, the effective flux $`\mathrm{\Phi }_\mathrm{v}`$ changes with time, and the voltage read by the voltmeter will beClem70 $`V_{ab}=d\mathrm{\Phi }_\mathrm{v}/dt`$. The power delivered by the power supply can be expressed in terms of the Lorentz force on the vortex, $`𝑱_\mathrm{d}\times \widehat{z}\varphi _0=I_\mathrm{d}\varphi _0G_\mathrm{d}`$; i.e., the rate at which work is done on the moving vortex is $`I_\mathrm{d}\varphi _0G_\mathrm{d}d𝒓/dt`$. Equating this to the power $`P=I_\mathrm{d}d\mathrm{\Phi }_\mathrm{v}/dt=I_\mathrm{d}\mathrm{\Phi }_\mathrm{v}d𝒓/dt`$ delivered by the power supply to maintain constant current, we obtain the equation $`\mathrm{\Phi }_\mathrm{v}(𝒓)=\varphi _0G_\mathrm{d}(𝒓)`$. Thus $`\mathrm{\Phi }_\mathrm{v}(𝒓)=\varphi _0G_\mathrm{d}(𝒓)+`$const, where the constant can have one of two possible values depending upon whether the integration contour $`C`$ is chosen to run inside or outside the vortex axis at $`𝒓=\widehat{x}x+\widehat{y}y`$. Choosing $`C`$ to run around the outer boundary of the SQUID, we obtain
$$\mathrm{\Phi }_\mathrm{v}(𝒓)=\varphi _0G_\mathrm{d}(𝒓).$$
(16)
Since $`G_\mathrm{d}(𝒓_o)=1`$ for points $`𝒓=𝒓_o`$ on the outer edges of the SQUID and $`G_\mathrm{d}(𝒓_i)=0`$ for points $`𝒓=𝒓_i`$ on the inner edges (at the perimeter of the central hole or along the edges of the slit), we have $`\mathrm{\Phi }_\mathrm{v}(𝒓_o)=\varphi _0`$ and $`\mathrm{\Phi }_\mathrm{v}(𝒓_i)=0`$. The derivation of Eq. (16 ) implicitly assumes that the vortex-core radius is much smaller than the linear dimensions of the SQUID.
We now return to the problem of how to find the currents $`I`$ and $`I_\mathrm{d}`$ in the SQUID, as well as the critical current $`I_\mathrm{c}`$. As discussed above, we have $`\mathrm{\Phi }_I=0`$ for a symmetric SQUID. For simplicity, we assume first that there are no vortices in the body of the SQUID, such that $`\mathrm{\Phi }_\mathrm{v}=0`$ and $`\mathrm{\Phi }=\mathrm{\Phi }_\mathrm{f}+\mathrm{\Phi }_\mathrm{d}`$ in Eq. (3), where $`\mathrm{\Phi }_\mathrm{d}=LI_\mathrm{d}`$. From the sum and the difference of $`I_1`$ and $`I_2`$ we obtain
$`I=2I_0\mathrm{cos}\left({\displaystyle \frac{\pi \mathrm{\Phi }_\mathrm{f}}{\varphi _0}}+{\displaystyle \frac{\pi LI_\mathrm{d}}{\varphi _0}}\right)\mathrm{sin}\overline{\varphi },`$ (17)
$`I_\mathrm{d}=I_0\mathrm{sin}\left({\displaystyle \frac{\pi \mathrm{\Phi }_\mathrm{f}}{\varphi _0}}+{\displaystyle \frac{\pi LI_\mathrm{d}}{\varphi _0}}\right)\mathrm{cos}\overline{\varphi },`$ (18)
where $`\overline{\varphi }=(\varphi _1+\varphi _2)/2`$ is determined experimentally by how much current is applied to the SQUID. When $`\overline{\varphi }=0`$, the current $`I`$ is zero. As $`\overline{\varphi }`$ increases, the magnitude of $`I`$ increases and reaches its maximum value $`I_\mathrm{c}`$ at a value of $`\overline{\varphi }`$ that must be determined by numerically solving Eqs. (17) and (18). A simple solution is obtained for arbitrary $`\mathrm{\Phi }_\mathrm{f}`$ only in the limit $`\pi LI_0/\varphi _00`$, for which $`I=I_\mathrm{c}=2I_0|\mathrm{cos}(\pi \mathrm{\Phi }_\mathrm{f}/\varphi _0)|`$ and $`I_\mathrm{d}=0`$ at the critical current. For values of $`\pi LI_0/\varphi _0`$ of order unity, as is the case for practical SQUIDs, one may obtain $`I_\mathrm{c}`$ for any value of $`\mathrm{\Phi }_\mathrm{f}`$ by solving Eq. (18) self-consistently for $`I_\mathrm{d}`$ for a series of values of $`\overline{\varphi }`$ and by substituting the results into Eq. (17) to determine which value of $`\overline{\varphi }`$ maximizes $`I`$. Equations (17) and (18) have been solved numerically by de Bruyn Ouboter and de Waele,deBruyn70 (some of their results are also shown by Orlando and DelinOrlando91 ), who showed that at $`I_\mathrm{c}`$
$`I_\mathrm{c}(\mathrm{\Phi }_\mathrm{f})`$ $`=`$ $`I_\mathrm{c}(\mathrm{\Phi }_\mathrm{f}+n\varphi _0)=I_\mathrm{c}(\mathrm{\Phi }_\mathrm{f}),`$ (19)
$`I_\mathrm{d}(\mathrm{\Phi }_\mathrm{f})`$ $`=`$ $`I_\mathrm{d}(\mathrm{\Phi }_\mathrm{f}+n\varphi _0)=I_\mathrm{d}(\mathrm{\Phi }_\mathrm{f}),`$ (20)
$`I_1(\mathrm{\Phi }_\mathrm{f})`$ $`=`$ $`I_1(\mathrm{\Phi }_\mathrm{f}+n\varphi _0)=I_2(\mathrm{\Phi }_\mathrm{f}),`$ (21)
$`I_2(\mathrm{\Phi }_\mathrm{f})`$ $`=`$ $`I_2(\mathrm{\Phi }_\mathrm{f}+n\varphi _0)=I_1(\mathrm{\Phi }_\mathrm{f}),`$ (22)
where $`n`$ is an integer. Hence all the physics is revealed by displaying $`I_\mathrm{c}(\mathrm{\Phi }_\mathrm{f})`$ over the interval $`0\mathrm{\Phi }_\mathrm{f}\varphi _0/2`$, as shown in Fig. 2.
When a vortex is present, Eqs. (17) and (18) still hold, except that $`\mathrm{\Phi }_\mathrm{f}`$ in these equations is replaced by the sum $`\mathrm{\Phi }_\mathrm{f}+\mathrm{\Phi }_\mathrm{v}`$. Thermally agitated motion of vortices in the body of the SQUID can produce flux noise via the term $`\mathrm{\Phi }_\mathrm{v}(𝒓_\mathrm{v})`$ and the time dependence of the vortex position $`𝒓_\mathrm{v}.`$ From Eq. (16) we see that the sensitivity of $`I_c`$ to vortex-position noise is proportional to the magnitude of $`\mathrm{\Phi }_\mathrm{v}(𝒓)=\varphi _0G_\mathrm{d}=𝑱_\mathrm{d}\times \widehat{z}\varphi _0/I_\mathrm{d}.`$ Thus $`I_\mathrm{c}`$ is most sensitive to vortex-position noise when the vortices are close to the inner or outer edges of the SQUID, where the magnitude of $`𝑱_\mathrm{d}`$ is largest. These equations provide more accurate results for the vortex-position sensitivity than the approximations given in Refs. Ferrari91, and Sun94, .
So far, we have investigated how the general equations governing the behavior of a dc SQUID are altered when the contributions arising from line integrals of the current density are included. As we have shown in Ref. Brandt04, , these additional contributions are important when the Pearl length $`\mathrm{\Lambda }`$ is an appreciable fraction of the linear dimensions of the SQUID. We have found that the basic SQUID equations, Eqs. (17) and (18), are unaltered, except that the magnetic flux (sometimes called $`\mathrm{\Phi }_{\mathrm{ext}}`$Orlando91 ) generated in the SQUID’s central hole by the externally applied field in the absence of a vortex is replaced by the effective flux $`\mathrm{\Phi }_\mathrm{f}`$, given in Eq. (14). Similarly, we have shown that the total inductance $`L`$ of the SQUID has contributions both from the magnetic induction (geometric inductance) and the associated supercurrent (kinetic inductance). We also have shown in principle how to calculate the effect of the return flux from a vortex at position $`𝒓_\mathrm{v}`$ in the body of the SQUID, and we have found that the effective flux arising from the vortex is $`\mathrm{\Phi }_\mathrm{v}(𝒓_\mathrm{v})`$, given in Eqs. (15) and (16). To demonstrate that all the above quantities can be calculated numerically for arbitrary values of $`\mathrm{\Lambda }`$, we next examine the behavior of a model SQUID as decribed in Sec. III.
## III Long SQUID in a perpendicular magnetic field
Here we consider a long SQUID whose thickness $`d`$ is less than the London penetration depth $`\lambda `$ and whose topology is like that of Fig. 1 but which is stretched to a large length $`l`$ in the $`y`$ direction, as sketched in Fig. 3. SQUIDs of similar geometry have been investigated experimentally in Refs. Pegrum99, ; Eulenburg99, ; Kuriki99, ; Mitchell02, ; Mitchell03, . We treat here the case for which the length $`l`$ is much larger than the width $`2w`$ of the body of the SQUID, and we focus on the current and field distributions in and near the left ($`w<x<a`$) and right ($`a<x<w`$) arms and near the center of the SQUID, where to a good approximation the current density $`j_y`$ is uniform across the thickness and depends only upon $`x`$, and the magnetic induction $`𝑩=\times 𝑨`$ depends only upon $`x`$ and $`z`$. In the equations that follow, we deal with the sheet-current density, whose component in the $`y`$ direction is $`J_y(x)=j_y(x)d`$.
The self-field magnetic induction generated by $`J_y(x)`$ is $`𝑩_J(x,z)=\times 𝑨_J(x,z)`$, where $`A_J(x,z)`$, the $`y`$ component of the vector potential obtained from Ampere’s law, is
$$A_J(x,z)=\frac{\mu _0}{2\pi }J_y(x^{})\mathrm{ln}\frac{C}{\sqrt{(xx^{})^2+z^2}}dx^{}.$$
(23)
The integration here and in the following equations is carried out only over the strips, and $`C`$ is a constant with dimensions of length remaining to be determined. In the presence of a perpendicular applied field $`𝑩_\mathrm{a}=\widehat{z}B_\mathrm{a}=\times 𝑨_\mathrm{f}`$, the total vector potential is $`𝑨=𝑨_J+𝑨_\mathrm{f}`$, where $`𝑨_\mathrm{f}=\widehat{y}B_\mathrm{a}x.`$
### III.1 Formal solutions
We now use the approach of Ref. Brandt01, to calculate the in-plane magnetic-induction and sheet-current distributions appearing in Eqs. (11)-(14) in Sec. II. For all of these contributions we shall take into account the in-plane ($`z=0`$) self-field contribution $`A_J(x)=A_J(x,0)`$ to the $`y`$ component of the vector potential, where
$$A_J(x)=\frac{\mu _0}{2\pi }J_y(x^{})\mathrm{ln}\frac{C}{|xx^{}|}dx^{}.$$
(24)
We first examine the equal-current case and consider the contributions $`B_I(x)`$, the $`z`$ component of $`𝑩_I(x)`$, and $`J_I(x)`$, the $`y`$ component of $`𝑱_I(x)`$, due to equal currents $`I/2`$ in the left and right sides of the SQUID. Since $`J_I(x)=J_I(x)`$, the corresponding $`y`$ component of the vector potential is also a symmetric function of $`x`$: $`A_I(x)=A_I(x)`$, where the subcripts $`I`$ refer to the equal-current case. There are no flux quanta between the strips $`(\mathrm{\Phi }_I=0`$), and the second term in the brackets on the right-hand side of Eq. (1) vanishes; $`(\varphi _0/2\pi )\gamma =0`$. However, the constant $`C`$ must be chosen such that $`J_I(x)=A_I(x)/\mu _0\mathrm{\Lambda }`$ in the superconductor. Combining this equation with Eq. (24), making use of the symmetry $`J_I(x)=J_I(x)`$, and noting that $`I=2_a^wJ_I(x)𝑑x`$, we obtain
$$\frac{I}{2\pi }\mathrm{ln}\frac{b}{C}=_a^w\left[\frac{1}{2\pi }\mathrm{ln}\frac{b^2}{|x^2x^{}_{}{}^{}2|}+\mathrm{\Lambda }\delta (xx^{})\right]J_I(x^{})𝑑x^{}$$
(25)
for $`a<x<w`$. Here $`b`$ can be chosen to be any convenient length, such as the length $`l`$ or $`w`$, but not $`C`$. We now define the inverse integral kernel $`K^{\mathrm{sy}}(x,x^{})`$ for the symmetric-current case via
$`{\displaystyle _a^w}K^{\mathrm{sy}}(x,x^{\prime \prime })\left[{\displaystyle \frac{1}{2\pi }}\mathrm{ln}{\displaystyle \frac{b^2}{|x^{{}_{}{}^{\prime \prime }2}x^{}_{}{}^{}2|}}+\mathrm{\Lambda }\delta (x^{\prime \prime }x^{})\right]𝑑x^{\prime \prime }`$
$`=\delta (xx^{}).`$ (26)
Applying this kernel to Eq. (25), we obtain
$$J_I(x)=\frac{I}{2\pi }\mathrm{ln}\left(\frac{b}{C}\right)_a^wK^{\mathrm{sy}}(x,x^{})𝑑x^{}.$$
(27)
Since $`I=2_a^wJ_I(x)𝑑x`$, we find
$$\mathrm{ln}\left(\frac{b}{C}\right)=\pi /_a^w_a^wK^{\mathrm{sy}}(x,x^{})𝑑x𝑑x^{},$$
(28)
such that
$$J_I(x)=\frac{I}{2}_a^wK^{\mathrm{sy}}(x,x^{})𝑑x^{}/_a^w_a^wK^{\mathrm{sy}}(x^{},x^{\prime \prime })𝑑x^{}𝑑x^{\prime \prime }.$$
(29)
For $`axw`$, the stream function is
$$G_I(x)=\frac{2}{I}_a^xJ_I(x^{})𝑑x^{},$$
(30)
and for $`wxa`$, $`G_I(x)=G_I(x)`$. The corresponding $`z`$ component of the magnetic induction $`B_I(x)`$ can be obtained from the Biot-Savart law or, since $`B_I(x)=dA_I(x)/dx,`$ from Eq. (24). Note that $`B_I(x)=B_I(x)`$. Although the kernel $`K^{\mathrm{sy}}(x,x^{})`$ depends upon $`b`$, we find numerically that $`J_I(x)`$ and $`B_I(x)`$ are independent of $`b`$.
We next examine the circulating-current case and consider the contributions $`B_\mathrm{d}(x)`$, the $`z`$ component of $`𝑩_\mathrm{d}(x)`$, and $`J_\mathrm{d}(x)`$, the $`y`$ component of $`𝑱_\mathrm{d}(x)`$, due to a circulating current $`I_\mathrm{d}`$; the current in the $`y`$ direction on the right side of the SQUID is $`I_2=I_\mathrm{d}`$ and that on the left side is $`I_1=I_\mathrm{d}`$. The vector potential is still given by Eq. (24), except that we add subscripts $`\mathrm{d}`$. However, since $`J_\mathrm{d}(x)=J_\mathrm{d}(x)`$, the vector potential is now an antisymmetric function of $`x`$; $`A_\mathrm{d}(x)=A_\mathrm{d}(x)`$. The circulating current is generated by the fluxoid $`\mathrm{\Phi }_\mathrm{d}`$ \[see Eq. (13)\] associated with a nonvanishing gradient of the phase $`\gamma `$ around the loop. The second term inside the bracket on the right-hand side of Eq. (1), $`(\varphi _0/2\pi )\gamma `$, is $`\widehat{y}\mathrm{\Phi }_\mathrm{d}/2l`$ for $`a<x<w`$ and $`\widehat{y}\mathrm{\Phi }_\mathrm{d}/2l`$ for $`w<x<a`$. Thus, $`J_\mathrm{d}(x)=[A_\mathrm{d}(x)\mathrm{\Phi }_\mathrm{d}/2l]/\mu _0\mathrm{\Lambda }`$ for $`a<x<w`$. Combining this equation with that for $`A_\mathrm{d}(x)`$, noting that the inductance of the SQUID is $`L=\mathrm{\Phi }_\mathrm{d}/I_\mathrm{d}`$, and making use of the symmetry $`J_\mathrm{d}(x)=J_\mathrm{d}(x)`$, we obtain
$$\frac{LI_\mathrm{d}}{2l}=\mu _0_a^w\left[\frac{1}{2\pi }\mathrm{ln}\left|\frac{x+x^{}}{xx^{}}\right|+\mathrm{\Lambda }\delta (xx^{})\right]J_\mathrm{d}(x^{})𝑑x^{}$$
(31)
for $`a<x<w`$. We now define the inverse integral kernel $`K^{\mathrm{as}}(x,x^{})`$ for the asymmetric-current case via
$`{\displaystyle _a^w}K^{\mathrm{as}}(x,x^{\prime \prime })\left[{\displaystyle \frac{1}{2\pi }}\mathrm{ln}\left|{\displaystyle \frac{x^{\prime \prime }+x^{}}{x^{\prime \prime }x^{}}}\right|+\mathrm{\Lambda }\delta (x^{\prime \prime }x^{})\right]𝑑x^{\prime \prime }`$
$`=\delta (xx^{}).`$ (32)
Applying this kernel to Eq. (31) and noting that $`I_\mathrm{d}=_a^wJ_\mathrm{d}(x)𝑑x`$, we obtain
$$J_\mathrm{d}(x)=\alpha I_\mathrm{d}_a^wK^{\mathrm{as}}(x,x^{})𝑑x^{}$$
(33)
and
$$L=2\alpha \mu _0l,$$
(34)
where
$$\alpha =1/_a^w_a^wK^{\mathrm{as}}(x,x^{})𝑑x𝑑x^{}$$
(35)
is a dimensionless function of $`a,b,w,`$ and $`\mathrm{\Lambda }`$, which we calculate numerically in the next section. For $`axw`$, the stream function is
$$G_\mathrm{d}(x)=\frac{1}{I_\mathrm{d}}_a^xJ_\mathrm{d}(x^{})𝑑x^{},$$
(36)
and for $`wxa`$, $`G_\mathrm{d}(x)=G_\mathrm{d}(x)`$.
In this formulation, as in Ref. Brandt04, , $`L=L_\mathrm{m}+L_\mathrm{k}`$ is the total inductance. The geometric inductance contribution $`L_\mathrm{m}=2E_\mathrm{m}/I_\mathrm{d}^2`$, where $`E_\mathrm{m}=l_a^wJ_\mathrm{d}(x)A_\mathrm{d}(x)𝑑x`$ is the stored magnetic energy, and the kinetic contribution $`L_\mathrm{k}=2E_\mathrm{k}/I_\mathrm{d}^2`$, where $`E_\mathrm{k}=\mu _0\mathrm{\Lambda }l_a^wJ_\mathrm{d}^2(x)𝑑x`$ is the kinetic energy of the supercurrent, can be calculated using Eq. (33) fromBrandt04 ; Meservey69 ; Yoshida92
$`L_\mathrm{m}`$ $`=`$ $`{\displaystyle \frac{\mu _0l}{\pi I_\mathrm{d}^2}}{\displaystyle _a^w}{\displaystyle _a^w}\mathrm{ln}\left|{\displaystyle \frac{x+x^{}}{xx^{}}}\right|J_\mathrm{d}(x)J_\mathrm{d}(x^{})𝑑x𝑑x^{},`$
$`L_\mathrm{k}`$ $`=`$ $`{\displaystyle \frac{2\mu _0l\mathrm{\Lambda }}{I_\mathrm{d}^2}}{\displaystyle _a^w}J_\mathrm{d}^2(x)𝑑x={\displaystyle \frac{2\mu _0l\mathrm{\Lambda }}{(wa)}}{\displaystyle \frac{<J_\mathrm{d}^2>}{<J_\mathrm{d}>^2}},`$ (38)
where the brackets ($`<>`$) denote averages over the film width. We can show that $`L_\mathrm{m}+L_\mathrm{k}=L`$ with the help of Eqs. (32) and (35).
The $`z`$ component of the magnetic induction $`B_\mathrm{d}(x)`$ generated by $`J_\mathrm{d}(x)`$ can be obtained from the Biot-Savart law, or, since $`B_\mathrm{d}(x)=dA_\mathrm{d}(x)/dx,`$ from Eq. (24). Note that $`B_\mathrm{d}(x)=B_\mathrm{d}(x)`$.
When $`lw`$, the magnetic moment in the $`z`$ direction generated by the circulating current $`I_\mathrm{d}`$ is (to lowest order in $`w/l`$)
$$m_\mathrm{d}=2l_a^wxJ_\mathrm{d}(x)𝑑x,$$
(39)
where the factor 2 accounts for the factBrandt94 that the currents along the $`y`$ direction and those along the $`x`$ direction at the ends (U-turn) give exactly the same contribution to $`m_\mathrm{d}`$, even in the limit $`l\mathrm{}`$. To next higher order in $`w/l`$ one has to replace $`l`$ in Eq. (39) by $`l(wa)q`$, where $`(wa)q/2`$ is the distance of the center of gravity of the $`x`$-component of the currents near each end from this end. For a single strip of width $`wa`$, one has, e.g., $`q=1/3`$ for the Bean critical state (with rectangular current stream lines) and $`q=0.47`$ for ideal screening ($`\mathrm{\Lambda }w`$).EHB1
We next examine flux focusing. As discussed in Sec. II, to calculate the effective area of the slot, we need to calculate the fields produced in response to a perpendicular applied magnetic induction $`B_\mathrm{a}`$, subject to the condition that no current flows through either junction. Since this is equivalent to having both junctions open-circuited, the problem reduces to finding the fields produced in the vicinity of a pair of long superconducting strips connected by a superconducting link at only one end, i.e., when the slot of width $`2a`$ between the two strips is open at one end. However, the desired fields may be regarded as the superposition of the solutions of two separate problems when the slot has closed ends: (a) the fields generated in response to $`B_\mathrm{a}`$, when $`\mathrm{\Phi }=0`$ (the zero-fluxoid case) and a clockwise screening current flows around the slot \[second term on the right-hand side of Eq. (41) below\], and (b) the fields generated in the absence of $`B_\mathrm{a}`$, when flux quanta in the amount of $`\mathrm{\Phi }_\mathrm{f}`$ are in the slot and a counterclockwise screening current flows around the slot \[first term on the right-hand side of Eq. (41)\]. The desired flux-focusing solution is obtained by setting the net circulating current equal to zero.
The equations describing the fields in the flux-focusing case are derived as follows. The $`z`$ component of the net magnetic induction $`B_\mathrm{f}(x)`$ is the sum of $`B_\mathrm{a}`$ and the self-field $`B_{s\mathrm{f}}(x)`$ generated by $`J_\mathrm{f}(x)`$. The vector potential $`A_y(x)`$ is the sum of $`B_\mathrm{a}x`$, which describes the applied magnetic induction, and the self-field contribution given by Eq. (24) but with subscripts f. Since $`J_\mathrm{f}(x)=J_\mathrm{f}(x)`$, the vector potential is again an antisymmetric function of $`x`$; $`A_\mathrm{f}(x)=A_\mathrm{f}(x)`$. The fluxoid $`\mathrm{\Phi }_\mathrm{f}`$ \[see Eq. (14)\] contributes a nonvanishing gradient of the phase $`\gamma `$ around the loop, such that the second term inside the brackets on the right-hand side of Eq. (1), $`(\varphi _0/2\pi )\gamma `$, is $`\widehat{y}\mathrm{\Phi }_\mathrm{f}/2l`$ for $`a<x<w`$ and $`\widehat{y}\mathrm{\Phi }_\mathrm{f}/2l`$ for $`w<x<a`$. Equation (1) yields $`J_\mathrm{f}(x)=[B_\mathrm{a}x+A_\mathrm{f}(x)\mathrm{\Phi }_\mathrm{f}/2l]/\mu _0\mathrm{\Lambda }`$ for $`a<x<w`$. Combining this equation with that for $`A_\mathrm{f}(x)`$ \[Eq. (24)\], making use of the symmetry $`J_\mathrm{f}(x)=J_\mathrm{f}(x)`$, and introducing the effective area via $`\mathrm{\Phi }_\mathrm{f}=B_\mathrm{a}A_{\mathrm{eff}}`$ \[see Sec. II\], we obtain
$`B_\mathrm{a}(A_{\mathrm{eff}}2lx)/2l=`$
$`\mu _0{\displaystyle _a^w}\left[{\displaystyle \frac{1}{2\pi }}\mathrm{ln}\left|{\displaystyle \frac{x+x^{}}{xx^{}}}\right|+\mathrm{\Lambda }\delta (xx^{})\right]J_\mathrm{f}(x^{})𝑑x^{}`$ (40)
for $`a<x<w`$. We again use the inverse integral kernel $`K^{\mathrm{as}}(x,x^{})`$ for the asymmetric-current case \[Eq. (32)\] to obtain
$$J_\mathrm{f}(x)=\frac{B_\mathrm{a}}{\mu _0}_a^w\left(\frac{A_{\mathrm{eff}}}{2l}x^{}\right)K^{\mathrm{as}}(x,x^{})𝑑x^{}.$$
(41)
The effective area of the SQUID $`A_{\mathrm{eff}}`$ is found from the condition that the net current around the loop is zero \[$`_a^wJ_\mathrm{f}(x)𝑑x=0`$\], which yields
$$A_{\mathrm{eff}}=2\alpha l_a^w_a^wx^{}K^{\mathrm{as}}(x,x^{})𝑑x𝑑x^{}.$$
(42)
For $`axw`$, the stream function is
$$G_\mathrm{f}(x)=_a^xJ_\mathrm{f}(x^{})𝑑x^{},$$
(43)
and for $`wxa`$, $`G_\mathrm{f}(x)=G_\mathrm{f}(x)`$. The spatial distribution of the resulting $`z`$ component of the in-plane magnetic induction is given by $`B_\mathrm{f}(x)=B_\mathrm{a}+B_{s\mathrm{f}}(x)`$, where $`B_{s\mathrm{f}}(x)`$ can be obtained from the Biot-Savart law or by substituting $`J_\mathrm{f}(x)`$ into Eq. (24) and making use of $`B_{s\mathrm{f}}(x)=dA_\mathrm{f}(x)/dx`$. Note that $`B_\mathrm{f}(x)=B_\mathrm{f}(x)`$. The resulting magnetic moment $`m_\mathrm{f}`$ in the $`z`$ direction can be calculated by replacing $`J_\mathrm{d}`$ by $`J_\mathrm{f}`$ in Eq. (39).
In the next section we also present numerical results for field and current distributions in the zero-fluxoid case, in which $`I=0`$, $`\mathrm{\Phi }_\mathrm{v}=0`$, and the effective flux is zero: $`\mathrm{\Phi }=\mathrm{\Phi }_\mathrm{d}+\mathrm{\Phi }_\mathrm{f}=0`$. Such a case could be achieved by short-circuiting the Josephson junctions in Fig. 1, cooling the device in zero field such that initially $`\mathrm{\Phi }=0`$, and then applying a small perpendicular magnetic induction $`B_\mathrm{a}`$. A circulating current $`J(x)`$, given by the second term on the right-hand side of Eq. (41), would spontaneously arise in order to keep $`\mathrm{\Phi }=0,`$ as in the Meissner state.
### III.2 Numerical solutions
In the previous section we have presented formal solutions for the sheet-current density $`J_y(x)`$ in Eqs. (29), (33), and (41), which are expressed as integrals involving the geometry-dependent inverse kernels $`K^{\mathrm{sy}}(x,x^{})`$ and $`K^{\mathrm{as}}(x,x^{})`$. As in Ref. Brandt04, for thin rings, these integrals are evaluated on a grid $`x_i`$ ($`i=1,2,\mathrm{},N`$) spanning only the strip (but avoiding the edges, where the integrand may have infinities), $`a<|x_i|<w`$, such that for any function $`f(x)`$ one has $`_a^wf(x)𝑑x=_{i=1}^Nw_if(x_i)`$. Here $`w_i`$ are the weights, approximately equal to the local spacing of the $`x_i`$; the weights obey $`_{i=1}^Nw_i=wa`$. We have chosen the grid such that the weights $`w_i`$ are narrower and the grid points $`x_i`$ more closely spaced near the edges $`a`$ and $`w`$, where $`J_y(x)`$ varies more rapidly. We have accomplished this by choosing some appropriate continuous function $`x(u)`$ and an auxiliary discrete variable $`u_ii\frac{1}{2}`$, such that $`w_i=x^{}(u_i)(u_2u_1)`$. We can choose $`x(u)`$ such that its derivative $`x^{}(u)`$ vanishes (or is reduced) at the strip edges to give a denser grid there. By choosing an appropriate substitution function $`x(u)`$ one can make the numerical error of this integration method arbitrarily small, decreasing rapidly with any desired negative power of the grid number $`N`$, e.g., $`N^2`$ or $`N^3`$.
For the equal-current case, Eq. (25) becomes
$$\frac{I}{2\pi }\mathrm{ln}\frac{b}{C}=\underset{j=1}{\overset{N}{}}(w_jQ_{ij}^{\mathrm{sy}}+\mathrm{\Lambda }\delta _{ij})J_I(x_j)$$
(44)
where $`\delta _{ij}=0`$ for $`ij`$, $`\delta _{ii}=1`$, and
$`Q_{ij}^{\mathrm{sy}}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\mathrm{ln}{\displaystyle \frac{b^2}{|x_i^2x_j^2|}},ij,`$
$`Q_{ii}^{\mathrm{sy}}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\mathrm{ln}{\displaystyle \frac{\pi b^2}{x_iw_i}}.`$ (45)
The optimum choice of the diagonal term $`Q_{ii}^{\mathrm{sy}}`$ (i.e., with $`|x_i^2x_j^2|`$ replaced by $`x_iw_i/\pi `$ for $`i=j`$, which reduces the numerical error from order $`N^1`$ to $`N^2`$ or higher depending on the grid) is discussed in Eq. (3.12) of Ref. Brandt94, for strips and in Eq. (18) of Ref. Brandt04, for disks and rings. The superscript (sy) is a reminder that this is for a symmetric current distribution \[$`J_I(x)=J_I(x)`$\]. Defining $`K_{ij}^{\mathrm{sy}}=(w_jQ_{ij}^{\mathrm{sy}}+\mathrm{\Lambda }\delta _{ij})^1`$, such that
$$\underset{k=1}{\overset{N}{}}K_{ik}^{\mathrm{sy}}(w_jQ_{kj}^{\mathrm{sy}}+\mathrm{\Lambda }\delta _{kj})=\delta _{ij},$$
(46)
and applying it to Eq. (44), we obtain
$$J_I(x_i)=\frac{I}{2\pi }\mathrm{ln}\left(\frac{b}{C}\right)\underset{j=1}{\overset{N}{}}K_{ij}^{\mathrm{sy}}.$$
(47)
Since $`I=2_{i=1}^Nw_iJ_I(x_i)`$, we find
$$\mathrm{ln}\left(\frac{b}{C}\right)=\pi /\underset{i=1}{\overset{N}{}}\underset{j=1}{\overset{N}{}}w_iK_{ij}^{\mathrm{sy}},$$
(48)
such that
$$J_I(x_i)=\frac{I}{2}\underset{j=1}{\overset{N}{}}K_{ij}^{\mathrm{sy}}/\underset{k=1}{\overset{N}{}}\underset{l=1}{\overset{N}{}}w_kK_{kl}^{\mathrm{sy}}.$$
(49)
It is remarkable that although the parameter $`b`$ appears in Eq. (45), the final result for $`J_I(x_i)`$ in Eq. (49) does not depend upon $`b`$. The stream function $`G_I(x)`$ can be evaluated asnote1
$$G_I(x_i)=\underset{j=1}{\overset{i}{}}\underset{k=1}{\overset{N}{}}w_jK_{jk}^{\mathrm{sy}}/\underset{l=1}{\overset{N}{}}\underset{m=1}{\overset{N}{}}w_lK_{lm}^{\mathrm{sy}}.$$
(50)
Shown in Fig. 4 are plots of $`J_I(x)`$, $`G_I(x)`$, and the corresponding magnetic induction $`B_I(x)`$ vs $`x`$ for $`a/w=0.3`$ and various values of $`\mathrm{\Lambda }/w=`$ 0, 0.03, 0.1, 0.3, 1. The curves for $`\mathrm{\Lambda }=0`$ exactly coincide with the analytic expressions of Appendix A. The magnetic field lines for this case are depicted in Fig. 5.
For the circulating-current case, Eq. (31) becomes
$$\frac{LI_\mathrm{d}}{2l}=\mu _0\underset{j=1}{\overset{N}{}}(w_jQ_{ij}^{\mathrm{as}}+\mathrm{\Lambda }\delta _{ij})J_\mathrm{d}(x_j)$$
(51)
where
$`Q_{ij}^{\mathrm{as}}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\mathrm{ln}{\displaystyle \frac{x_i+x_j}{|x_ix_j|}},ij,`$
$`Q_{ii}^{\mathrm{as}}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\mathrm{ln}{\displaystyle \frac{4\pi x_i}{w_i}},`$ (52)
The superscript (as) is a reminder that this is for an asymmetric current distribution \[$`J_\mathrm{d}(x)=J_\mathrm{d}(x)`$\]. Defining $`K_{ij}^{\mathrm{as}}=(w_jQ_{ij}^{\mathrm{as}}+\mathrm{\Lambda }\delta _{ij})^1`$, such that
$$\underset{k=1}{\overset{N}{}}K_{ik}^{\mathrm{as}}(w_jQ_{kj}^{\mathrm{as}}+\mathrm{\Lambda }\delta _{kj})=\delta _{ij},$$
(53)
applying it to Eq. (51), and noting that $`I_\mathrm{d}=_{i=1}^Nw_iJ_\mathrm{d}(x_i)`$, we obtain
$$J_\mathrm{d}(x_i)=\alpha I_\mathrm{d}\underset{j=1}{\overset{N}{}}K_{ij}^{\mathrm{as}}$$
(54)
and
$$L=2\alpha \mu _0l,$$
(55)
where
$$\alpha =1/\underset{i=1}{\overset{N}{}}\underset{j=1}{\overset{N}{}}w_iK_{ij}^{\mathrm{as}}.$$
(56)
The stream function $`G_\mathrm{d}(x)`$ can be evaluated asnote1
$$G_\mathrm{d}(x_i)=\alpha \underset{j=1}{\overset{i}{}}\underset{k=1}{\overset{N}{}}w_jK_{jk}^{\mathrm{sy}}.$$
(57)
Shown in Fig. 6 are plots of $`J_\mathrm{d}(x)`$, $`G_\mathrm{d}(x)`$, and $`B_\mathrm{d}(x)`$ vs $`x`$ for $`a/w=0.3`$ and various values of $`\mathrm{\Lambda }/w`$. Note that these curves look similar to those in Fig. 4, but they all have opposite parity, as can be seen from the different profiles $`B(x)`$ near $`x=0`$. The magnetic field lines for this case are shown in Fig. 7.
As discussed in Sec. II, when a vortex is present in the region $`a<|x|<w`$, the sensitivity of the SQUID’s critical current $`I_c`$ is proportional to the magnitude of $`d\mathrm{\Phi }_\mathrm{v}/dx=\varphi _0dG_\mathrm{d}/dx=\varphi _0J_\mathrm{d}(x)/I_\mathrm{d}`$. From Fig. 6 we see that when $`\mathrm{\Lambda }w`$, this sensitivity is greatly enhanced when the vortex is close to the edges $`a`$ and $`w`$ but that when $`\mathrm{\Lambda }w`$, the sensitivity is nearly independent of position.
Shown as the solid curves in Fig. 8(a) are plots of the inductance $`L`$ vs $`a/w`$ for various values of $`\mathrm{\Lambda }/w=`$ 0, 0.03, 0.1, 0.3, and 1. The solid curves in Fig. 9(a) show the same $`L`$ vs $`\mathrm{\Lambda }/w`$ (range 0.0045 to 2.2) for several values of $`a/w=`$ 0.01, 0.1, 0.4, 0.8, 0.95, and 0.99.
The geometric and kinetic contributions $`L_\mathrm{m}`$ and $`L_\mathrm{k}`$ can be calculated separately from Eqs. (LABEL:Lm) and (38)
$`L_\mathrm{m}`$ $`=`$ $`{\displaystyle \frac{2\mu _0l}{I_\mathrm{d}^2}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{j=1}{\overset{N}{}}}w_iw_jQ_{ij}^{\mathrm{as}}J_\mathrm{d}(x_i)J_\mathrm{d}(x_j),`$ (58)
$`L_\mathrm{k}`$ $`=`$ $`{\displaystyle \frac{2\mu _0l\mathrm{\Lambda }}{I_\mathrm{d}^2}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{j=1}{\overset{N}{}}}w_iJ_\mathrm{d}^2(x_i),`$ (59)
using Eqs. (54) and (56). We can show that $`L_\mathrm{m}+L_\mathrm{k}=L`$ using the property of inverse matrices that $`𝑴𝑴^1=𝑴^1𝑴=𝑰,`$ where $`𝑰`$ is the identity matrix.
Shown as solid curves in Fig. 8(b) are $`L_\mathrm{m}`$ and $`L_\mathrm{k}`$ vs $`a/w`$. For $`\mathrm{\Lambda }=0`$, when $`L_\mathrm{k}=0`$, $`L=L_\mathrm{m}`$ exactly coincides with Eq. (A10), which may be approximated by Eq. (A11) for $`a/w<0.7`$ \[open circles in Fig. 8(a)\] and by Eq. (A12) for $`a/w>0.7`$ \[open squares in Fig. 8(a)\]. The dotted curves in Fig. 8(b) for $`L_\mathrm{m}`$ and $`L_\mathrm{k}`$ are those of Eqs. (C3) and (C6) in the limit $`\mathrm{\Lambda }/w\mathrm{}`$, when the circulating current density is uniform. The dotted curves in Fig. 8(a) are obtained from $`L=L_\mathrm{m}+L_\mathrm{k}`$ using the approximations of Eqs. (C3) and (C6); they are an excellent approximation to $`L`$ for $`\mathrm{\Lambda }/w0.03`$ except for small values of $`a/w`$. Improved agreement for small values of $`\mathrm{\Lambda }/w`$ and $`a/w`$ is shown by the dashed curves in Fig. 8(a), which show the approximation of Eq. (B7) for $`L`$, and in Fig. 8(b), which show the approximation of Eq. (B12) for $`L_\mathrm{k}`$.
The solid curves in Fig. 9(b) show $`L_\mathrm{m}`$ and $`L_\mathrm{k}`$ vs $`\mathrm{\Lambda }/w`$. The geometric inductance $`L_\mathrm{m}`$ depends upon $`\mathrm{\Lambda }`$ but only weakly, varying slowly between its $`\mathrm{\Lambda }=0`$ asymptote \[Eq. (A10), horizontal dot-dashed line\] and its $`\mathrm{\Lambda }=\mathrm{}`$ asymptote \[Eq. (C3), horizontal dotted line\]. For larger values of of $`a/w`$, $`L_\mathrm{m}`$ is nearly independent of $`\mathrm{\Lambda }`$. On the other hand, the kinetic inductance $`L_\mathrm{m}`$, is approximately proportional to $`\mathrm{\Lambda }/w`$. The straight dotted lines in Fig. 9(b), calculated from the large-$`\mathrm{\Lambda }`$ approximation given in Eq. (C6), are a good approximation to $`L_\mathrm{m}`$ except for small values of $`\mathrm{\Lambda }/w`$ and $`a/w`$. The dotted curves in Fig. 9(a) are obtained from $`L=L_\mathrm{m}+L_\mathrm{k}`$ using the approximations of Eqs. (C3) and (C6). Improved agreement for small values of $`\mathrm{\Lambda }/w`$ and $`a/w`$ is shown by the dashed curves in Fig. 9(a), which show the approximation of Eq. (B7) for $`L`$, and in Fig. 9(b), which show the approximation of Eq. (B12) for $`L_\mathrm{k}`$.
In Eq. (3) of Ref. Yoshida92, , Yoshida et al. derived an approximate expression for the kinetic inductance when $`\mathrm{\Lambda }/w1`$. We have found that their expression for $`L_\mathrm{k}`$ is not an accurate approximation to our exact numerical results. To eliminate the logarithmic divergences due to the inverse square-root dependence of the current density near the edges, Yoshida et al. followed an approach used by Meservey and Tedrow,Meservey69 and chose a cutoff length of the order of $`d`$, the film thickness. When $`d<\lambda `$, however, this approach cannot be correct, because the equations describing the fields and currents in superconducting strips contain only the two-dimensional screening length $`\mathrm{\Lambda }=\lambda ^2/d`$. The cutoff length therefore must instead be chosen to be of the order of $`\mathrm{\Lambda }`$, as we have done in Appendix B.
The magnetic moment in the $`z`$ direction generated by the circulating current is, from Eq. (39),
$$m_\mathrm{d}=2\underset{i=1}{\overset{N}{}}w_ix_iJ_\mathrm{d}(x_i).$$
(60)
As shown in Fig. 10, this magnetic moment vanishes very slowly when the gap width and $`\mathrm{\Lambda }`$ go to zero, $`a/w0`$ and $`\mathrm{\Lambda }/w0`$. This can be explained by the fact that for $`\mathrm{\Lambda }=0`$ and $`a<xw`$ one has $`J_d(x)1/x`$, Eq. (A5). The contribution of these small $`x`$ to $`m_d`$, Eq. (60), stays finite due to the factor $`x`$, but the total current $`I_d`$ to which $`m_d`$ is normalized, diverges when $`a/w0`$, thus suppressing the plotted ratio $`m_d/I_d`$. Interestingly, the curves in Fig. 10 coincide with those in Fig. 13; see below. Expressions for $`m_\mathrm{d}`$ in the limits $`\mathrm{\Lambda }/w0`$ and $`\mathrm{\Lambda }/w\mathrm{}`$ are given in Eqs. (A13) and (C7).
For the flux-focusing case, Eq. (40) becomes
$$B_\mathrm{a}(A_{\mathrm{eff}}/2lx_i)=\mu _0\underset{j=1}{\overset{N}{}}(w_jQ_{ij}^{\mathrm{as}}+\mathrm{\Lambda }\delta _{ij})J_\mathrm{f}(x_j).$$
(61)
Applying Eq. (53), we obtain
$$J_\mathrm{f}(x_i)=\frac{B_\mathrm{a}}{\mu _0}\left(\frac{A_{\mathrm{eff}}}{2l}\underset{j=1}{\overset{N}{}}K_{ij}^{\mathrm{as}}\underset{j=1}{\overset{N}{}}K_{ij}^{\mathrm{as}}x_j\right),$$
(62)
where, since $`_{i=1}^Nw_iJ_\mathrm{f}(x_i)=0`$, the effective area is
$$A_{\mathrm{eff}}=2\alpha l\underset{i=1}{\overset{N}{}}\underset{j=1}{\overset{N}{}}w_iK_{ij}^{\mathrm{as}}x_j.$$
(63)
The stream function $`G_\mathrm{f}(x)`$ can be evaluated asnote1
$$G_\mathrm{f}(x_i)=\underset{j=1}{\overset{i}{}}w_jJ_\mathrm{f}(x_j).$$
(64)
Shown in Fig. 11 are plots of the flux-focusing $`J_\mathrm{f}(x)`$, $`G_\mathrm{f}(x)`$, and $`B_\mathrm{f}(x)`$ vs $`x`$ for $`a/w=0.3`$ and $`\mathrm{\Lambda }/w=`$ 0, 0.03, 0.1, 0.3, and 1. The first term in Eq. (62) equals the circulating-current sheet-current density, Eq. (54), with appropriate weight factor such that the total circulating current vanishes, $`I_1=I_2=0`$. The corresponding magnetic field lines are depicted in Fig. 12. Shown in Figs. 13 and 14 are plots of the effective area $`A_{\mathrm{eff}}(a/w,\mathrm{\Lambda }/w)`$ versus $`a/w`$ and $`\mathrm{\Lambda }/w`$, respectively, in units of the maximum possible area $`2wl`$. In the limit $`\mathrm{\Lambda }/w0`$, $`A_{\mathrm{eff}}`$ is given by Eq. (A19), and when $`\mathrm{\Lambda }/w\mathrm{}`$, $`A_{\mathrm{eff}}=l(w+a).`$ Note in Fig. 14 that $`A_{\mathrm{eff}}`$ increases with increasing $`\mathrm{\Lambda }`$, particularly for small gap widths $`2a`$. Flux focusing is reflected by the fact that for small $`a/w0`$ the effective area $`A_{\mathrm{eff}}`$ of the gap tends to a constant, except in the limit $`\mathrm{\Lambda }0`$, where it vanishes very slowly, $`A_{\mathrm{eff}}/2wl(\pi /2)/\mathrm{ln}(4w/a)`$ \[Eq. (A19)\]. When $`a/w0`$, the enhancement factor $`A_{\mathrm{eff}}/2al\mathrm{}`$ and thus diverges even for $`\mathrm{\Lambda }=0`$. In the limit $`a/w1`$, $`A_{\mathrm{eff}}(0,\mathrm{\Lambda }/w)`$ tends to a universal function \[see Fig. 14\]. Interestingly, Figs. 13 and 10 show identical curves; this is because the identity $`A_{\mathrm{eff}}=m_d/I_d`$ holds for all values of $`a/w`$ and $`\mathrm{\Lambda }/w`$, as proved in general in Sec. II.
Figures 15 and 16 show the minimum of the magnetic induction in the flux-focusing case, $`B_\mathrm{f}(0)=B_\mathrm{f}(x=0)`$ \[see Fig. 11\], plotted versus $`a/w`$ and $`\mathrm{\Lambda }/w`$, respectively. The ratio $`B_\mathrm{f}(0)/B_a1`$ tends to unity for $`a/w1`$ and for $`\mathrm{\Lambda }/w1`$, and it diverges for $`a/w0`$ when $`\mathrm{\Lambda }=0`$. The curve for $`\mathrm{\Lambda }=0`$ exactly coincides with the analytic expression $`B_f(0)/B_a=w𝑬(k^{})/a𝑲(k^{})`$ obtained from Eq. (A17). For $`a/w1`$ this yields $`B_f(0)/B_a(w/a)/\mathrm{ln}(4w/a)`$, which is a good approximation for $`0<a/w0.3`$. MB1 The magnetic moment $`m_\mathrm{f}`$ for the flux-focusing case, calculated from Eq. (60) but with $`J_\mathrm{d}(x_i)`$ replaced by $`J_\mathrm{f}(x_i)`$, is shown in Fig. 17. Expressions for $`m_\mathrm{f}`$ in the limits $`\mathrm{\Lambda }/w0`$ and $`\mathrm{\Lambda }/w\mathrm{}`$ are given in Eqs. (A20) and (C9)
Shown in Fig. 18 are profiles for the zero-fluxoid case with plots of $`J(x)`$ and the corresponding $`B(x)`$ generated by an applied magnetic induction $`B_a>0`$ when the junctions are short-circuited such that $`\mathrm{\Phi }=0`$ and $`I_1=I_2>0`$; for comparison see analogous profiles in Sec. 2.5 of Ref. Babaei02, for two parallel strips and in Sec. IV of Ref. Babaei03, and Sec. 4 of Ref. Brandt04, for rings. That the current density $`J(x)`$ in the zero-fluxoid case is given by the second term on the right-hand sides of Eqs. (41) and (62), can be seen by setting $`\mathrm{\Phi }_\mathrm{f}=B_\mathrm{a}A_{\mathrm{eff}}=0`$ in Eqs. (40), (41), (61), and (62). Depicted in Fig. 18 are the examples $`a/w=0.3`$ with $`\mathrm{\Lambda }/w=`$ 0, 0.03, 0.1, 0.3, and 1. Figure 19 shows the magnetic field lines for this case and Fig. 20 the magnetic moment $`m`$. Expressions for $`J`$, $`B`$, and $`m`$ for the zero-fluxoid case in the limits $`\mathrm{\Lambda }/w0`$ and $`\mathrm{\Lambda }/w\mathrm{}`$ are given in Appendixes A and C.
## IV Summary
In Sec. II of this paper we have presented general equations governing the static behavior of a thin-film dc SQUID for all values of the Pearl length $`\mathrm{\Lambda }=\lambda ^2/d`$, where the London penetration depth $`\lambda `$ is larger than $`d`$, the film thickness. The SQUID’s critical current $`I_c`$ depends upon the effective flux $`\mathrm{\Phi }`$, which is the sum of the magnetic flux up through a contour surrounding the central hole and a term proportional to the line integral of the current density around this contour. For a symmetric SQUID there are three important contributions to $`\mathrm{\Phi }`$: a circulating-current term $`\mathrm{\Phi }_\mathrm{d}`$, a vortex-field term $`\mathrm{\Phi }_\mathrm{v}`$, and a flux-focusing term $`\mathrm{\Phi }_\mathrm{f}`$, all of which depend upon $`\mathrm{\Lambda }`$. Since $`\mathrm{\Lambda }`$ is a function of temperature, an important consequence is that all of the contributions to $`\mathrm{\Phi }`$ are temperature-dependent.
The circulating-current term $`\mathrm{\Phi }_\mathrm{d}`$ can be expressed in terms of the SQUID inductance $`L`$ and the circulating current $`I_\mathrm{d}`$ via $`\mathrm{\Phi }_\mathrm{d}=LI_\mathrm{d}`$. The SQUID inductance has two contributions, $`L=L_\mathrm{m}+L_\mathrm{k}`$, where the first term is the geometric inductance (associated with the energy stored in the magnetic field) and the second is the kinetic inductance (associated with the kinetic energy of the circulating supercurrent). Both contributions are functions of $`\mathrm{\Lambda }`$, since they both depend on the spatial distribution of the current density. However, $`L_\mathrm{m}`$ depends only weakly upon $`\mathrm{\Lambda }`$, because for the same circulating current $`I_\mathrm{d}`$, the energy stored in the magnetic field does not vary greatly as $`\mathrm{\Lambda }`$ ranges from zero to infinity. On the other hand, because the kinetic energy density is proportional to $`\mathrm{\Lambda }`$, $`L_\mathrm{k}`$ is also nearly proportional to $`\mathrm{\Lambda }`$, with deviations from linearity occurring only for small values of $`\mathrm{\Lambda }/w`$.
The vortex-field term can be written as $`\mathrm{\Phi }_\mathrm{v}=\varphi _0G_\mathrm{d}`$, where $`G_\mathrm{d}`$ is a dimensionless stream function describing the circulating sheet-current density $`𝑱_d`$. Roughly speaking, when $`\mathrm{\Lambda }`$ is small, $`I_c`$ is most strongly dependent upon the vortex position when the vortex is close to the edges of the film, but when $`\mathrm{\Lambda }`$ is large, $`I_c`$ is equally sensitive to the vortex position wherever the vortex is. Recent experimentsDoenitz05 have used the relationship $`\mathrm{\Phi }_\mathrm{v}=\varphi _0G_\mathrm{d}`$ to determine the vortex-free sheet-current density $`𝑱_d(x,y)`$ from vortex images obtained via low-temperature scanning electron microscopy.Straub01 ; Doenitz04 The experimental data obtained in magnetic fields up to 40 $`\mu `$T are in excellent agreement with numerical calculations of $`𝑱_d(x,y)`$, confirming the validity of the above relationship, even in the presence of many (up to 200) vortices in the SQUID washer.
The flux-focusing term can be expressed as $`\mathrm{\Phi }_\mathrm{f}=B_\mathrm{a}A_{\mathrm{eff}}`$, where $`B_\mathrm{a}`$ is the applied magnetic induction and $`A_{\mathrm{eff}}`$ is the effective area of the central hole of the SQUID. Although $`A_{\mathrm{eff}}`$ is primarily determined by the dimensions of the SQUID, it also depends upon the value of $`\mathrm{\Lambda }.`$
To illustrate the $`\mathrm{\Lambda }`$ dependence of the above quantities, in Sec. III of this paper we analyzed in detail the behavior of a long SQUID whose central region resembles a coplanar stripline. We numerically calculated the profiles of the sheet-current density, stream function, and magnetic induction in the equal-current, circulating-current, flux-focusing, and zero-fluxoid cases for various representative values of $`\mathrm{\Lambda }`$. We presented plots of the inductances $`L,`$ $`L_\mathrm{m}`$, and $`L_\mathrm{k}`$, the effective area $`A_{\mathrm{eff}}`$, and the magnetic moments for these cases. Useful analytic approximations are provided for the $`\mathrm{\Lambda }/w0`$ limit in Appendix A, for small $`\mathrm{\Lambda }/w`$ and $`a/w`$ in Appendix B, and for the $`\mathrm{\Lambda }/w\mathrm{}`$ limit in Appendix C.
We are in the process of applying the above theory to square and circular SQUIDs, using the numerical method of Ref. Brandt05, .
###### Acknowledgements.
We thank D. Koelle for stimulating discussions. This work was supported in part by Iowa State University of Science and Technology under Contract No. W-7405-ENG-82 with the U.S. Department of Energy and in part by the German Israeli Research Grant Agreement (GIF) No G-705-50.14/01.
## Appendix A the limit $`\mathrm{\Lambda }/w=0`$
In the ideal-screening limit $`\mathrm{\Lambda }/w=0`$, the $`y`$ component of the sheet-current density in the strips ($`a<|x|<w`$) for the equal-current case isBabaei02
$$J_I(x)=\frac{I}{\pi }\frac{|x|}{[(x^2a^2)(w^2x^2)]^{1/2}},$$
(65)
and the $`z`$ component of the magnetic induction in the plane $`z=0`$ of the strips is
$`B_I(x)`$ $`=`$ $`{\displaystyle \frac{\mu _0I}{2\pi }}{\displaystyle \frac{x}{[(x^2a^2)(x^2w^2)]^{1/2}}},|x|>w,`$ (66)
$`=`$ $`0,a<|x|<w,`$ (67)
$`=`$ $`{\displaystyle \frac{\mu _0I}{2\pi }}{\displaystyle \frac{x}{[(a^2x^2)(w^2x^2)]^{1/2}}},|x|<a,`$ (68)
and the constant $`C`$ in Eqs. (24), (25), (27), and (28) is $`C=\sqrt{w^2a^2}/2.`$
For the circulating-current case, the $`y`$ component of the sheet-current density in the strips ($`a<|x|<w`$) in the limit $`\mathrm{\Lambda }/w=0`$ isBabaei02
$$J_\mathrm{d}(x)=\frac{2B_0}{\mu _0}\frac{x}{|x|}\frac{w^2}{[(x^2a^2)(w^2x^2)]^{1/2}},$$
(69)
and the $`z`$ component of the magnetic induction in the plane $`z=0`$ of the strips is
$`B_\mathrm{d}(x)`$ $`=`$ $`B_0{\displaystyle \frac{w^2}{[(x^2a^2)(x^2w^2)]^{1/2}}},|x|>w,`$ (70)
$`=`$ $`0,a<|x|<w,`$ (71)
$`=`$ $`B_0{\displaystyle \frac{w^2}{[(a^2x^2)(w^2x^2)]^{1/2}}},|x|<a,`$ (72)
where the parameter $`B_0`$, the magnetic flux $`\mathrm{\Phi }_\mathrm{d}`$ in the $`z`$ direction in the slot, the circulating current $`I_\mathrm{d}`$, and the geometric inductance $`L_\mathrm{m}`$ are related by
$$\mathrm{\Phi }_\mathrm{d}=L_\mathrm{m}I_\mathrm{d}=2B_0lw𝑲(k)$$
(73)
and
$$L_\mathrm{m}=\mu _0l𝑲(k)/𝑲(k^{}),$$
(74)
where $`𝑲(k)`$ is the complete elliptic integral of the first kind of modulus $`k=a/w`$ and complementary modulus $`k^{}=\sqrt{1k^2}`$. The geometric inductance is well approximated for small $`a/w`$ by
$$L_\mathrm{m}=(\pi \mu _0l/2)/\mathrm{ln}(4w/a),$$
(75)
neglecting corrections proportional to $`a^2/w^2`$, and for small $`(wa)/w`$ by
$$L_\mathrm{m}=(\mu _0l/\pi )\mathrm{ln}[16/(1a^2/w^2)],$$
(76)
neglecting corrections proportional to $`1a^2/w^2`$. In the limit that $`\mathrm{\Lambda }=0`$, the kinetic inductance vanishes ($`L_\mathrm{k}=0)`$, and the inductance in Eq. (74) becomes the total inductance: $`L=L_\mathrm{m}`$. The magnetic moment \[see Eq. (39)\] can be obtained from Eqs. (13)-(16) of Ref. Babaei02, :
$$m_\mathrm{d}=[\pi lw/𝑲(k^{})]I_\mathrm{d}.$$
(77)
For the flux-focusing case, the $`y`$ component of the sheet-current density in the strips ($`a<|x|<w`$) in the limit $`\mathrm{\Lambda }/w=0`$ isBabaei02
$$J_\mathrm{f}(x)=\frac{2B_\mathrm{a}}{\mu _0𝑲(k^{})}\frac{x}{|x|}\frac{𝑬(k^{})w^22𝑲(k^{})x^2}{[(x^2a^2)(w^2x^2)]^{1/2}},$$
(78)
and the $`z`$ component of the magnetic induction in the plane $`z=0`$ of the strips is
$`B_\mathrm{f}(x)`$ $`=`$ $`{\displaystyle \frac{B_\mathrm{a}}{𝑲(k^{})}}{\displaystyle \frac{𝑬(k^{})w^22𝑲(k^{})x^2}{[(x^2a^2)(x^2w^2)]^{1/2}}},`$ (79)
$`|x|>w,`$
$`=`$ $`0,a<|x|<w,`$ (81)
$`=`$ $`{\displaystyle \frac{B_\mathrm{a}}{𝑲(k^{})}}{\displaystyle \frac{𝑬(k^{})w^22𝑲(k^{})x^2}{[(a^2x^2)(w^2x^2)]^{1/2}}},`$
$`|x|<a,`$
where $`𝑬(k^{})`$ is the complete elliptic integral of the second kind of complementary modulus $`k^{}=\sqrt{1k^2}`$ and modulus $`k=a/w`$. The magnetic flux in the $`z`$ direction in the slot is
$$\mathrm{\Phi }_\mathrm{f}=\pi B_\mathrm{a}lw/𝑲(k^{}),$$
(82)
and the effective area $`A_{\mathrm{eff}}`$ of the slot is
$$A_{\mathrm{eff}}=\mathrm{\Phi }_\mathrm{f}/B_\mathrm{a}=\pi lw/𝑲(k^{}).$$
(83)
Note that $`A_{\mathrm{eff}}=m_\mathrm{d}/I_\mathrm{d}.`$ The magnetic moment generated by $`J_\mathrm{f}(x)`$ is
$$m_\mathrm{f}=\pi l[w^2+a^22w^2𝑬(k^{})/𝑲(k^{})]B_\mathrm{a}/\mu _0.$$
(84)
For the zero-fluxoid case, the $`y`$ component of the sheet-current density in the strips can be obtained from Sec. 2.5 of Ref. Babaei02, :
$$J(x)=\frac{2B_\mathrm{a}}{\mu _0}\frac{x}{|x|}\frac{x^2[1𝑬(k)/𝑲(k)]w^2}{[(x^2a^2)(w^2x^2)]^{1/2}}.$$
(85)
The corresponding $`z`$ component of the magnetic induction isBabaei02
$`B(x)`$ $`=`$ $`B_\mathrm{a}{\displaystyle \frac{x^2[1𝑬(k)/𝑲(k)]w^2}{[(x^2a^2)(x^2w^2)]^{1/2}}},|x|>w,`$ (86)
$`=`$ $`0,a<|x|<w,`$ (87)
$`=`$ $`B_\mathrm{a}{\displaystyle \frac{[1𝑬(k)/𝑲(k)]w^2x^2}{[(a^2x^2)(w^2x^2)]^{1/2}}},|x|<a.`$ (88)
The magnetic moment generated by $`J(x)`$ is
$$m=\pi l[2w^2𝑬(k)/𝑲(k)w^2+a^2]B_\mathrm{a}/\mu _0.$$
(89)
## Appendix B Behavior for small $`\mathrm{\Lambda }`$ and small $`a`$
In this section we present some expressions for $`L`$, $`L_\mathrm{k}`$, and $`L_\mathrm{m}`$ that follow from approximating the circulating-current distribution for small values of $`\mathrm{\Lambda }`$ and $`a`$.
When the slot is very narrow ($`a/w1`$), we approximate the sheet-current density in the region $`a<x<w`$ generated by the fluxoid $`\mathrm{\Phi }_\mathrm{d}`$ via
$$J_\mathrm{d}(x)=I_0\frac{w}{\sqrt{(x^2a^2+\delta ^2)(w^2x^2+\delta ^2)}},$$
(90)
where $`I_0=2\mathrm{\Phi }_\mathrm{d}/\pi \mu _0l`$ and $`\delta `$ is a quantity of order $`\mathrm{\Lambda }=\lambda ^2/d`$ determined as follows. When $`l\mathrm{}`$ and then $`a0`$ and $`w\mathrm{}`$, an exact calculation yields for $`x>0`$
$$J_y(x)=\frac{2\mathrm{\Phi }_\mathrm{d}}{\pi \mu _0l\mathrm{\Lambda }}_0^{\mathrm{}}\frac{e^{xt/\mathrm{\Lambda }}dt}{t^2+1}.$$
(91)
We find $`_0^bJ_y(x)𝑑x=(2\mathrm{\Phi }_\mathrm{d}/\pi \mu _0l)\mathrm{ln}(\gamma b/2\mathrm{\Lambda })`$ when $`b\mathrm{\Lambda }`$, where $`\gamma =e^C=1.781\mathrm{}`$, and $`C=0.577\mathrm{}`$ is Euler’s constant. From Eq. (90) we find $`_0^bJ_\mathrm{d}(x)𝑑x=(2\mathrm{\Phi }_\mathrm{d}/\pi \mu _0l)\mathrm{ln}(2b/\delta )`$ when $`a=0`$ and $`\delta bw`$. Comparing these two integrals we obtain $`\delta =(4/\gamma )\mathrm{\Lambda }=2.246\mathrm{\Lambda }`$. Integrating Eq. (90) from $`a`$ to $`w`$ to obtain $`I_\mathrm{d}`$, we find
$$I_\mathrm{d}=I_0\frac{w}{\sqrt{w^2+\delta ^2}}[F(\lambda _a,q)F(\lambda _w,q)],$$
(92)
where $`F(\varphi ,k)`$ is the elliptic integral of the first kind and
$`\lambda _a`$ $`=`$ $`\mathrm{arcsin}\sqrt{{\displaystyle \frac{w^2a^2+\delta ^2}{w^2a^2+2\delta ^2}}},`$ (93)
$`\lambda _w`$ $`=`$ $`\mathrm{arcsin}{\displaystyle \frac{\delta }{\sqrt{w^2a^2+2\delta ^2}}},`$ (94)
$`q`$ $`=`$ $`\sqrt{{\displaystyle \frac{w^2a^2+2\delta ^2}{w^2+\delta ^2}}}.`$ (95)
Expanding Eq. (92) for $`aw`$ and $`\delta w`$, using $`I_0=2\mathrm{\Phi }_\mathrm{d}/\pi \mu _0l`$, and neglecting terms of order $`a^2/w^2`$ and $`\delta ^2/w^2`$, we obtain
$$L=\mathrm{\Phi }_\mathrm{d}/I_\mathrm{d}=(\pi \mu _0l/2)/\{\mathrm{ln}[4w/(a+\delta )]\delta /w\},$$
(96)
where $`\delta =2.246\mathrm{\Lambda }`$. Note that Eq. (B7) reduces to Eq. (A11) when $`\mathrm{\Lambda }=0`$.
From Eq. (90) we obtain the approximation
$$_a^wJ_\mathrm{d}^2(x)𝑑x=I_0^2\frac{w}{(w^2a^2+\delta ^2)}(f_a+f_w),$$
(97)
where
$`f_a`$ $`=`$ $`{\displaystyle \frac{w}{\sqrt{\delta ^2a^2}}}\mathrm{tan}^1{\displaystyle \frac{(wa)\sqrt{\delta ^2a^2}}{a(wa)+\delta ^2}},`$ (99)
$`a<\delta ,`$
$`=`$ $`{\displaystyle \frac{w}{\sqrt{a^2\delta ^2}}}\mathrm{tanh}^1{\displaystyle \frac{(wa)\sqrt{a^2\delta ^2}}{a(wa)+\delta ^2}},`$
$`a>\delta ,`$
$`f_w`$ $`=`$ $`{\displaystyle \frac{w}{\sqrt{w^2+\delta ^2}}}\mathrm{tanh}^1{\displaystyle \frac{(wa)\sqrt{w^2+\delta ^2}}{w(wa)+\delta ^2}}.`$ (100)
Using Eqs. (92) and (97), we obtain from Eq. (38)
$$L_\mathrm{k}=2\mu _0l\frac{\mathrm{\Lambda }}{w}\frac{(w^2+\delta ^2)}{(w^2a^2+2\delta ^2)}\frac{(f_a+f_w)}{[F(\lambda _a,q)F(\lambda _w,q)]^2},$$
(101)
where $`\delta =2.246\mathrm{\Lambda }`$. Although our intention in using the ansatz of Eq. (90) initially was to obtain an improved approximation to $`L_\mathrm{k}`$ for small values of $`\mathrm{\Lambda }`$ and $`a`$, we see from Figs. 8(b) and 9(b) that Eq. (101) provides a reasonably good approximation for all values of $`\mathrm{\Lambda }`$ and $`a`$.
## Appendix C the limit $`\mathrm{\Lambda }/w\mathrm{}`$
In the weak-screening limit $`\mathrm{\Lambda }/w\mathrm{}`$, $`K^{\mathrm{sy}}(x,x^{})=\mathrm{\Lambda }^1\delta (xx^{})`$, the $`y`$ component of the sheet-current density in the strips ($`a<|x|<w`$) in the equal-current case is uniform, $`J_I=I/2(wa)`$, the $`z`$ component of the magnetic induction in the plane of the strips obtained from the Biot-Savart law is
$$B_I(x)=\frac{\mu _0I}{4\pi (wa)}\mathrm{ln}\left|\frac{(xw)(x+a)}{(x+w)(xa)}\right|,$$
(102)
and the constant $`C`$ in Eqs. (24), (25), (27), and (28) is $`C=w\mathrm{exp}[\pi \mathrm{\Lambda }/(wa)].`$
For the circulating-current case in the limit $`\mathrm{\Lambda }/w\mathrm{}`$, $`K^{\mathrm{as}}(x,x^{})=\mathrm{\Lambda }^1\delta (xx^{}),`$ $`\alpha =\mathrm{\Lambda }/(wa)`$, the $`y`$ component of the sheet-current density in the strips is again uniform, $`J_\mathrm{d}=I_\mathrm{d}/(wa)`$ for $`a<x<w`$, and the $`z`$ component of the magnetic induction in the plane of the strips obtained from the Biot-Savart law is
$$B_\mathrm{d}(x)=\frac{\mu _0I_\mathrm{d}}{2\pi (wa)}\mathrm{ln}\left|\frac{x^2w^2}{x^2a^2}\right|.$$
(103)
The geometric inductance is, from Eq. (LABEL:Lm)
$`L_\mathrm{m}`$ $`=`$ $`{\displaystyle \frac{\mu _0l}{\pi (wa)^2}}[w^2\mathrm{ln}\left({\displaystyle \frac{4w^2}{w^2a^2}}\right)`$ (104)
$`2aw\mathrm{ln}\left({\displaystyle \frac{w+a}{wa}}\right)+a^2\mathrm{ln}\left({\displaystyle \frac{4a^2}{w^2a^2}}\right)],`$
which is independent of $`\mathrm{\Lambda }`$. Equation (C3) is well approximated for small $`a/w`$ by
$$L_\mathrm{m}=(\mu _0l/\pi )(1+2a/w)\mathrm{ln}4,$$
(105)
neglecting corrections proportional to $`a^2/w^2`$, and for small $`(wa)/w`$ by
$$L_\mathrm{m}=(\mu _0l/\pi )\left[\mathrm{ln}\frac{2}{1a/w}+\frac{3}{2}\frac{1}{2}(1a/w)\right],$$
(106)
neglecting corrections proportional to $`(1a/w)^2`$. From Eq. (38) we obtain the kinetic inductance
$$L_\mathrm{k}=2\mu _0l\mathrm{\Lambda }/(wa).$$
(107)
When $`\mathrm{\Lambda }w`$, the total inductance $`L`$ is dominated by the kinetic inductance ($`L_\mathrm{k}L_\mathrm{m}`$), such that $`LL_\mathrm{k}`$. Since $`J_\mathrm{d}`$ is uniform, the magnetic moment is easily found from Eq. (39) to be
$$m_\mathrm{d}=l(w+a)I_\mathrm{d}.$$
(108)
For the flux-focusing case in the limit $`\mathrm{\Lambda }/w\mathrm{}`$, $`K^{\mathrm{as}}(x,x^{})=\mathrm{\Lambda }^1\delta (xx^{}),`$ the $`y`$ component of the sheet-current density is $`J_\mathrm{f}(x)=B_\mathrm{a}(w+a2x)/2\mu _0\mathrm{\Lambda }`$ for $`a<x<w`$, the effective area is $`A_{\mathrm{eff}}=l(w+a)=m_\mathrm{d}/I_\mathrm{d}`$, and the $`z`$ component of the magnetic induction in the plane of the strips is $`B_\mathrm{f}(x)=B_\mathrm{a}+B_{s\mathrm{f}}(x)`$, where from the Biot-Savart law
$`B_{s\mathrm{f}}(x)`$ $`=`$ $`{\displaystyle \frac{B_\mathrm{a}}{4\pi \mathrm{\Lambda }}}\left[(w+a)\mathrm{ln}\right|{\displaystyle \frac{x^2w^2}{x^2a^2}}|`$ (109)
$`+`$ $`2x\mathrm{ln}\left|{\displaystyle \frac{(x+w)(xa)}{(xw)(x+a)}}\right|4(wa)].`$
The magnetic moment generated by $`J_\mathrm{f}(x)`$ in this limit is
$$m_\mathrm{f}=[l(wa)^3/6\mathrm{\Lambda }](B_\mathrm{a}/\mu _0).$$
(110)
For the zero-fluxoid case in the limit $`\mathrm{\Lambda }/w\mathrm{}`$, the applied field is only weakly screened, and the $`z`$ component of the magnetic flux density is nearly equal to the applied magnetic induction, $`B(x)B_\mathrm{a}`$. The $`y`$ component of the vector potential is approximately given by $`A(x)=B_\mathrm{a}x`$, and the $`y`$ component of the induced sheet-current density, obtained from Eq. (1) with $`\gamma =0`$, is $`J(x)=(B_\mathrm{a}/\mu _0\mathrm{\Lambda })x`$. To the next order of approximation, $`B(x)=B_\mathrm{a}+B_s(x),`$ where the self-field $`B_s`$ is found from the Biot-Savart law
$$B_s(x)=\frac{B_\mathrm{a}}{2\pi \mathrm{\Lambda }}\left[x\mathrm{ln}\left|\frac{(x+w)(xa)}{(xw)(x+a)}\right|2(wa)\right].$$
(111)
The magnetic moment generated by $`J(x)`$ in this limit is
$$m=[2l(w^3a^3)/3\mathrm{\Lambda }](B_\mathrm{a}/\mu _0).$$
(112) |
warning/0507/q-bio0507021.html | ar5iv | text | # Spontaneous polarization in eukaryotic gradient sensing: A mathematical model based on mutual inhibition of frontness and backness pathways
## 1 Introduction
When motile cells are exposed to a chemoattractant gradient, they develop a morphological polarity consisting of a distinct front and back . The formation of the morphological polarity is driven by the spatial segregation of distinct sets of intracellular molecules to the front and the rear of the cell (see Figure 1a and refs. ). The *frontness* molecules, which include Cdc42, Rac, PI3P, PI3K, Arp2/3, and F-actin, localize to the front of the cell where they coordinate the extension of an actin-rich extension. The *backness* molecules, which include Rho, Rho kinase, PTEN, and myosin II, migrate to the rear of the cell where they are thought to activate cell contraction.
In general, the spatial segregation of the frontness/backness molecules and the resultant morphological polarization occurs even if the cells are exposed to a uniform chemoattractant profile (see Figure 1b and ). This phenomenon has been called *spontaneous polarization* to emphasize the fact that the cells polarize despite the absence of a perceptible external cue .
The existence of spontaneous polarization is reminiscent of the Turing instability in reaction-diffusion systems . Thus, it has led to several models that view spontaneous polarization, either explicitly or implicitly, as the onset of a Turing bifurcation .<sup>2</sup><sup>2</sup>2Alternative models of gradient sensing in which spontaneous polarization plays no role have also been proposed . For the most part, they are motivated by experimental systems such as PDGF-stimulated fibroblasts and latrunculin-treated *Dictyostelium* cells , which do not display spontaneous polarization. In these models, polarization occurs only in the presence of external gradients. According to these models, the cell is in a stable homogeneous steady state in the absence of chemoattractant. However, exposure of the cell to a sufficiently large uniform chemoattractant concentration pushes it past a Turing bifurcation point, where the homogeneous steady state is unstable with respect to certain non-homogeneous perturbations. The inevitable presence of noise is then sufficient to drive the cell towards a nonhomogeneous steady state corresponding to the polarized state of the cell.
As attractive as these models of spontaneous polarization may be, there is a significant gap between the theory and experiments. Indeed, all the models consist of one or more slow-diffusing *activators* whose synthesis is autocatalytic, and which in consequence, tend to grow and spread across the entire cell. The unconstrained growth and dispersion of the activator is restricted by hypothesizing the existence of a diffusible *inhibitor* (which is a by-product of activator synthesis and impedes the growth of the activator) or *substrate* (which is consumed during activator synthesis and stimulates the growth of the activator). The diffusible activator/substrate ensures that the growth of the activator(s) remains confined to a localized region of the cell membrane. This region is identified with the front of the cell, and the rest of the cell membrane, suffering from an activator deficit, is presumed to constitute the back of the cell. The predictions of these models are partially consistent with experiments involving the activation or inhibition of the frontness molecules. Indeed, if PI3K is inhibited, the gradient of the frontness molecules such as PI3P progressively decreases, until at sufficiently high levels of inhibition, there is no gradient at all (Figure 2a). On the other hand, when Rac is overexpressed, high levels of frontness molecules, PI3P and Rac, are found all over the cell membrane (Figure 2b). Both these features are reproduced by the activator-inhibitor class of models . However, these models cannot explain the spatial distribution of the backness molecules in response to inhibition of the frontness molecules. Specifically, if the activity of Cdc42 is suppressed, Rho is found not only at the back but also at the front of the cell (Figure 3a). Likewise, if the Gi proteins are inhibited with pertussis toxin (PTX), a uropod-like structure forms at the up-gradient edge of the cell (Figure 3b). The models also offer no insight into experiments involving activation or inhibition of the backness molecules. For instance, when Rho is activated, PI3P fails to polarize (Figure 4a). Conversely, when Rho is inhibited, PI3P spreads all over the membrane, and the cell extends a single broad pseudopod or multiple pseudopods (Figure 4b). These results are beyond the scope of the activator-inhibitor and activator-substrate models because the backness molecules are not even acknowledged as legitimate variables — they are implicitly assumed to somehow settle down in regions uninhabited by the frontness molecules. This stands in sharp contrast to the experimental data which shows that backness pathways downregulate the frontness pathways: Inhibition of Rho kinase increases Rac activity 2–3 fold .
One hypothesis regarding the interaction between the frontness and backness pathways is that they inhibit each other . Based on extensive experiments with neutrophils and neutrophil-like HL-60 cells , Bourne and coworkers arrived at the kinetic scheme shown in Figure 5a, which they describe as follows
> Briefly, the attractant binds to a G protein-coupled receptor (R), which in turn activates different trimeric G proteins to generate two divergent, opposing signaling pathways, which promote polarized frontness and backness, respectively. In the frontness pathway, Gi, PI3Ps, and Rac promote de novo formation of actin polymers. One or more positive feedback loops in this first pathway mediate localized increases in sensitivity to attractant: one of these requires polymerized actin, while Rac or Cdc 42 may in addition enhance PI3P accumulation more directly, via an alternative pathway (dotted curved line in Figure 7 \[reproduced here as Figure 5a\]). Backness signals, generated by G12 and G13, depend on activation of a Rho-dependent pathway that stimulates activation of myosin II, formation of contractile actin-myosin complexes, and myosin-dependent inhibition of Rac- and PI3P-dependent responses. Backness signals inhibit frontness signals, and vice versa (dashed straight lines in Figure 7).
They go on to explain spontaneous polarization in terms of this kinetic scheme as follows
> The more or less symmetrically distributed actin ruffles and PI3P accumulation seen at early times (e.g., 30 s) after application of a uniform stimulus presumably mask a fine-textured mosaic of interspersed backness and frontness signals, some triggering activation of PI3Ks, Rac, and actin polymerization, others promoting activation of Rho and myosin. Localized mechanical incompatibility of the two cytoskeletal responses, combined with the ability of each to damp signals that promote the other (dashed inhibitor lines in Figure 7), then gradually drive them to separate into distinct domains of the membrane.
The goal of this work is to formulate a mathematical model of spontaneous polarization based on the foregoing mechanism, namely, mutual inhibition of the frontness and backness pathways.
It turns out that the Bourne model imposes two requirements that are incompatible with Turing instabilities in a two-component system consisting of frontness and backness pathways. Specifically, the mutual inhibition between the frontness and backness pathways must be sufficiently strong to ensure that the homogeneous steady state is unstable in the presence of diffusion, and yet weak enough to guarantee its stability in the absence of diffusion. These two requirements cannot be satisfied simultaneously in a two-component system, a point that will be discussed in Section 3.
Thus, we are led to consider the modification of the Bourne model shown in Figure 5b. It contains two variables, $`U_2,U_3`$, representing the frontness and backness pathways, respectively, whose activation is driven by receptor ligation, and which inhibit each other in a concentration-dependent manner. It is assumed furthermore that the mutual inhibition between $`U_2`$ and $`U_3`$ is so strong that they are mutually incompatible (in a sense that will be made mathematically precise below). We refer to these two variables as *activators.* The model differs from the Bourne picture inasmuch as it assumes the existence of a diffusible *inhibitor*, denoted $`U_1`$, whose synthesis is promoted by both activators, but which, in turn, inhibits both activators. The inhibitor serves to stabilize the coexistence of the otherwise incompatible activators: Transient increases in activator concentrations at localized “hotspots” are efficiently suppressed by the concominant increase in the concentration of the inhibitor. However, the mobility of the inhibitor imposes constraints upon its stabilizing effect because it tends to diffuse away from the “hotspot.” This is not an issue at low chemoattractant concentrations, for under these conditions, the activator concentrations, and hence, their mutual incompatibility, is so small that despite the high diffusibility of the inhibitor, it successfully damps fluctuations of the activator levels. However, at high chemoattractant concentrations, the activator concentrations and their mutual incompatibility are so large that the diffusible inhibitor fails to suppress the two activators sufficiently. The mutual incompatibility of the two activators now overcomes the mollifying effect of the inhibitor, and the activators segregate spatially into separate domains.
The mathematical model in this paper quantifies the above physical argument. We show that the model yields spontaneous polarization at sufficiently high active receptor levels. Moreover, if the frontness or backness pathways are inhibited or activated, they redistribute spatially in a manner consistent with the experiments described above. Finally, we show the model displays steady states corresponding to both chemoattraction and chemorepulsion. Interest in this phenomenon is motivated by the fact that when neurons are exposed to activators of the cGMP pathway, chemorepulsion turns into chemorepulsion . Based on the analysis of the model, we suggest that such transitions can be triggered by altering the balance of power in the mutually inhibitory interactions between the frontness and backness pathways.
The mechanism of spontaneous polarization in this model is distinct from that in activator-inhibitor or activator-substrate models. Unlike these models, the spatial segregation of the two activators is driven entirely by their mutual inhibition — positive feedback plays no role. In reality, both positive feedback and mutual inhibition cooperate to produce symmetry-breaking (Figure 5a). The contribution of this work is to highlight the distinct role of mutual inhibition, a feature that emphasized in the experimental literature , but absent from previous activator-inhibitor and activator-substrate models.
The paper is organized as follows. In Section 2, we define the model and simulate the experiments described above. In Section 3, we elaborate on the physics underlying the model. Finally, we summarize the conclusions.
## 2 Results
### 2.1 The model
We assume that
1. Both activators promote the synthesis of the inhibitor, which in turn, degrades by a first-order process. Thus, the net rate of synthesis of $`U_1`$ is
$$R_1u_1+A_{12}u_2+A_{13}u_3.$$
2. The synthesis of $`U_2`$ is receptor-mediated. It is autocatalytic at low concentrations and self-limiting at high concentrations, i.e.,
$$R_2ru_2A_{22}u_2^2.$$
where $`r`$ denotes the receptor activity, and the second term represents a self-limiting process that prevents synthesis rate from increasing beyond bounds.
3. The synthesis of $`U_2`$ is inhibited in a concentration-dependent manner by the common inhibitor, $`U_1`$, and the other activator, $`U_3`$. Assuming that these interactions follow bimolecular kinetics, the net rate of synthesis of $`U_2`$ is
$$R_2ru_2A_{22}u_2^2A_{21}u_1u_2A_{23}u_2u_3.$$
4. The synthesis of $`U_3`$ follows kinetics similar to those of $`U_2`$, i.e., its net rate of synthesis is
$$R_3ru_3A_{33}u_3^2A_{31}u_1u_3A_{32}u_2u_3$$
Given these assumptions, we arrive at the equations
$`{\displaystyle \frac{u_1}{T}}`$ $`=D_1{\displaystyle \frac{u_1}{X^2}}R_1u_1+A_{12}u_2+A_{13}u_3`$ (1)
$`{\displaystyle \frac{u_2}{T}}`$ $`=D_2{\displaystyle \frac{u_2}{X^2}}+\left(R_2rA_{21}u_1A_{23}u_3A_{22}u_2\right)u_2`$ (2)
$`{\displaystyle \frac{u_3}{T}}`$ $`=D_3{\displaystyle \frac{u_3}{X^2}}+\left(R_3rA_{31}u_1A_{32}u_2A_{33}u_3\right)u_3`$ (3)
where $`X`$ denotes the spatial coordinate, $`T`$ denotes time, and $`u_1,u_2,u_3`$ denote the concentrations of $`U_1,U_2,U_3`$, respectively. We assume the Neumann boundary conditions
$$\frac{u_i}{X}=0\text{ at }x=0,L$$
since they imply no interaction with the environment, which is consistent with our desire to study spontaneous (autonomous) polarization.
It is convenient to rescale the equations. If we define
$`t=a_{11}T,`$ $`x={\displaystyle \frac{X}{L}},`$ $`d_i={\displaystyle \frac{D_i/L^2}{R_1}},`$
$`a_{ij}={\displaystyle \frac{A_{ij}}{R_i}},`$ $`\rho _2={\displaystyle \frac{R_2}{R_1}},`$ $`\rho _3={\displaystyle \frac{R_3}{R_1}},`$
we obtain the scaled equations
$`{\displaystyle \frac{u_1}{t}}`$ $`=d_1{\displaystyle \frac{u_1}{x^2}}u_1+a_{12}u_2+a_{13}u_3`$ (4)
$`{\displaystyle \frac{u_2}{t}}`$ $`=d_2{\displaystyle \frac{u_2}{x^2}}+\rho _2\left(ra_{21}u_1a_{22}u_2a_{23}u_3\right)u_2`$ (5)
$`{\displaystyle \frac{u_3}{t}}`$ $`=d_3{\displaystyle \frac{u_3}{x^2}}+\rho _3\left(ra_{31}u_1a_{32}u_2a_{33}u_3\right)u_3`$ (6)
and boundary conditions
$$\frac{u_i}{x}=0\text{ at }x=0,1.$$
Our goal is to study the variation of the steady states as a function of the receptor activity, $`r`$. We shall show, in particular, that the model has a homogeneous steady state at which $`u_2`$ and $`u_3`$ coexist, but it becomes Turing unstable at a sufficiently large $`r`$. The non-homogeneous steady state emerging from the Turing bifurcation is such that $`u_2`$ and $`u_3`$ are spatially segregated.
### 2.2 The homogeneous steady states
A Turing instability occurs when a non-homogeneous steady state bifurcates from a homogeneous steady state that is stable in the absence of diffusion. It is therefore useful to characterize the conditions for stability of the homogeneous steady states in the absence of diffusion.
The model has three types of homogeneous steady states: the trivial steady state, $`u_1=u_2=u_3=0`$, denoted $`E_{000}`$; the semitrivial steady states, $`u_1,u_2>0,u_3=0`$, and $`u_1,u_3>,u_2=0`$, denoted $`E_{110}`$ and $`E_{101}`$, respectively; and the nontrivial steady state, $`u_1,u_2,u_3>0`$, denoted $`E_{111}`$. The key results, depicted graphically in Figure 6a, are as follows (see Appendix A for details):
1. The trivial steady state, $`E_{000}`$, exists for all $`r>0`$, and is always unstable.
2. The semi-trivial steady state, $`E_{110}`$, which lies on the $`u_1u_2`$-plane, exists for all $`r>0`$. It is stable if and only if
$$\beta _{32}\frac{\alpha _{32}}{\alpha _{22}}>1.$$
Here, $`\alpha _{22}a_{22}+a_{21}a_{12}`$ is a measure of the extent to which $`U_2`$ inhibits itself both directly ($`a_{22}`$) and indirectly through production of $`U_1`$ ($`a_{21}a_{12}`$). Likewise, $`\alpha _{32}a_{32}+a_{31}a_{12}`$ is a measure of the extent to which $`U_2`$ inhibits $`U_3`$ directly ($`a_{32}`$) and indirectly through production of $`U_1`$ ($`a_{31}a_{12})`$. The foregoing stability condition then says that $`E_{110}`$ is stable if and only if $`U_2`$ inhibits $`U_3`$ more than it inhibits itself.
3. The semi-trivial steady state, $`E_{101}`$, which lies on the $`u_1u_3`$-plane, exists for all $`r>0`$. It is stable if and only if
$$\beta _{23}\frac{\alpha _{23}}{\alpha _{33}}>1,$$
where $`\alpha _{23}a_{23}+a_{21}a_{13},\alpha _{33}a_{33}+a_{31}a_{13}`$ are measures of the extent to which $`U_3`$ inhibits $`U_2`$ and itself, respectively, and $`\beta _{23}`$ is measure of the cross-inhibition relative to the self-inhibition.
4. The coexistence steady state, $`E_{111}`$, exists if and only if
$$\beta _{23},\beta _{32}<1\text{ or }\beta _{23},\beta _{32}>1.$$
It is stable only if the mutual inhibition of $`U_2`$ and $`U_3`$ is sufficiently weak, i.e., $`\beta _{23}\beta _{32}<1`$.
Figure 6a is reminiscent of the bifurcation diagram for the Lotka-Volterra model . This is not surprising because $`U_2`$ and $`U_3`$ obey Lotka-Volterra dynamics in the absence of the inhibitor. Indeed, the equations obtained by letting $`u_1=0`$ in equations (56) are identical to the Lotka-Volterra model for two competing species. It follows that the dynamics of $`U_2`$ and $`U_3`$ in the absence of the inhibitor are described by a bifurcation diagram very similar to Figure 6a, the only difference being that $`\beta _{23}`$ and $`\beta _{32}`$ must now be replaced by $`b_{23}a_{23}/a_{22},b_{32}a_{32}/a_{33}`$ (see Figure 6b). We shall appeal to this fact below.
### 2.3 Turing instability of homogeneous steady states
Since $`E_{000}`$ is always unstable, it can never undergo a Turing bifurcation. However, the semi-trivial and non-trivial steady states are stable for all sufficiently small $`r>0`$. The question then arises whether these steady states can undergo a Turing bifurcation. It is shown in Appendix B that
1. The semi-trivial steady states cannot undergo a Turing bifurcation.
2. The nontrivial steady state, $`E_{111}`$, can undergo a Turing bifurcation, but this is so only if
$$b_{23}b_{32}>1,$$
(7)
i.e., the mutual inhibition between $`U_2`$ and $`U_3`$ must be sufficiently strong — so strong, in particular, that they cannot coexist in the absence of the inhibitor (see Figure 6b).
Both conclusions are a consequence of the following fact which will be discussed in Section 3: In this model, the only destabilizing mechanism driving the Turing instability is mutual inhibition of $`U_2`$ and $`U_3`$, which requires the existence of *both* activators. Thus, the semi-trivial steady states fail to undergo a Turing bifurcation because one of the two activators is absent at such steady states. The non-trivial steady state, which is characterized by positive concentrations of $`U_2`$ and $`U_3`$, allows for a Turing instability, but only if their mutual inhibition is sufficiently strong.
Although the mutual inhibition must be sufficiently strong to ensure that $`E_{111}`$ undergoes a Turing instability, intuition suggests that it cannot be too strong, lest the two activators become incompatible even in the presence of the inhibitor. This is indeed the case. To see this, observe that (7) can be satisfied in three different ways (see Figure 6b):
1. $`b_{23},b_{32}>1`$, i.e., the mutual inhibition between $`U_2`$ and $`U_3`$ is so strong that they display bistable dynamics in the absence of the inhibitor.
2. $`b_{23}<1,b_{32}>1`$, i.e., in the absence of the inhibitor, $`U_2`$ inhibits $`U_3`$ much more than $`U_3`$ inhibits $`U_2`$, so that $`U_2`$ ultimately prevails over $`U_3`$.
3. $`b_{23}>1,b_{32}<1`$, i.e., in the absence of the inhibitor, $`U_3`$ inhibits $`U_2`$ much more than $`U_2`$ inhibits $`U_3`$, so that $`U_3`$ ultimately prevails over $`U_2`$.
It turns out that in the first case, when both $`U_2`$ and $`U_3`$ inhibit each other strongly, the necessary condition (7) is satisfied. Yet, the Turing instability cannot be realized because a stable coexistence steady state fails to exist even in the presence of the inhibitor. This follows from the fact that $`E_{111}`$ exists and is stable only if the mutual inhibition is sufficiently weak in the presence of the inhibitor, i.e.,
$$\beta _{23}=\frac{a_{23}+a_{21}a_{13}}{a_{33}+a_{31}a_{13}}<1,\beta _{32}=\frac{a_{32}+a_{31}a_{12}}{a_{22}+a_{21}a_{12}}<1$$
which can be recast in the form
$$\frac{a_{22}}{a_{12}}\left(b_{32}1\right)<a_{21}a_{31}<\frac{a_{33}}{a_{13}}\left(1b_{23}\right).$$
(8)
Clearly, (8) cannot be satisfied if $`b_{23},b_{32}>1`$. Under these conditions, the mutual inhibition of $`U_2`$ and $`U_3`$ is so strong that they cannot coexist stably even in the presence of the inhibitor.
We conclude that $`E_{111}`$ exists and bifurcates via a Turing instability only if the interaction between $`U_2`$ and $`U_3`$ in the absence of the inhibitor is such that only one of them prevails ultimately. Furthermore
1. If $`U_2`$ prevails over $`U_3`$ in the absence of the inhibitor ($`b_{23}<1,b_{32}>1`$), then a Turing instability obtains only if $`b_{23}b_{32}>1`$ and
$$0<\frac{a_{22}}{a_{12}}\left(b_{32}1\right)<a_{21}a_{31}<\frac{a_{33}}{a_{13}}\left(1b_{23}\right).$$
(9)
2. Conversely, if $`U_3`$ prevails over $`U_2`$ in the absence of the inhibitor ($`b_{23}>1,b_{32}<1`$), then a Turing instability obtains only if $`b_{23}b_{32}>1`$ and
$$\frac{a_{22}}{a_{12}}\left(b_{32}1\right)<a_{21}a_{31}<\frac{a_{33}}{a_{13}}\left(1b_{23}\right)<0.$$
(10)
These conditions have a simple physical interpretation. Consider, for instance, the condition (9). It says that if $`U_2`$ prevails over $`U_3`$ in the absence of the inhibitor, a Turing instability obtains only if (a) $`a_{21}a_{31}>0`$, i.e., $`U_1`$ inhibits $`U_2`$ more than it inhibits $`U_3`$. The stronger inhibition of $`U_2`$ by $`U_1`$ is necessary in order to compensate for its intrinsic superiority over $`U_3`$. (b) The magnitude of this difference must be neither too small nor too large to prevent under- or over-compensation that would preclude the coexistence of $`U_2`$ and $`U_3`$ even in the presence of the inhibitor.
It is shown in Appendix B that the conditions (9) or (10) are not only necessary but almost sufficient for $`E_{111}`$ to undergo a Turing instability. It suffices to impose the additional condition that $`d_1`$ be sufficiently larger than $`d_2`$ and $`d_3`$.
### 2.4 Simulation of experiments
To simulate the data shown in Figure 1b, wherein the frontness and backness components segregate spontaneously in response to a uniform chemoattractant profile, we chose parameter values satisfying (9), shown in the column labeled “$`U_2`$ wins” of Table 1. Linear stability analysis shows that given these parameter values, the coexistence steady state, $`E_{111}`$, undergoes a Turing instability at $`r0.6`$ and wavenumber $`k1`$ (Figure 11a). Computations with the continuation software package CONTENT confirm the existence of this instability (Figure 7a). The homogeneous steady state is stable for $`0<r0.6`$, and Turing unstable therafter. The concentration profiles of the non-homogeneous steady state created at the Turing bifurcation point are consistent with the data. Figure 7b shows that the non-homogeneous steady state at $`r=1`$ is such that $`U_2`$ and $`U_1`$ are in phase, and $`U_2`$ and $`U_3`$ are out of phase. The latter is consistent with the data in Figure 1b. We refer to this steady state as the *chemoattraction* steady state for the following reason. To a first degree of approximation, the steady state obtained in the presence of a receptor gradient, $`r(x)=1+ϵ\mathrm{cos}(\pi x),0<ϵ1`$, is identical to the steady state shown in Figure 7b. Thus, the profile of the frontness pathways is in phase with the distribution of the active receptors, which is characteristic of chemoattraction.
Figures 2 and 3 show the redistribution of the frontness and backness pathways in response to inhibition and activation of the frontness pathways. To simulate these experiments, we computed the variation of the chemoattraction steady state at $`r=1`$ with respect to $`R_2`$. Figure 8a shows that as $`R_2`$ decreases, so does $`u_2^{\text{max}}`$ until it becomes zero, i.e., the chemoattraction steady state ceases to exist, and is replaced by the homogeneous steady state, $`E_{101}`$, in which $`u_2=0`$ throughout the spatial domain. The variation of the concentration profile for the frontness pathway is consistent with the data shown in Figure 2. Figures 8b shows that as $`R_2`$ decreases, so does the gradient of the frontness molecules until it vanishes at $`R_21.5`$, which is in qualitative agreement with the data shown in Figure 2a. Conversely, at sufficiently large values of $`R_2`$, the non-homogeneous steady state is replaced by the homogeneous steady state, $`E_{110}`$, characterized by high levels of the frontness molecules uniformly distributed throughout the spatial domain (see Figure 2b). The variation of the concentration profile for the backness pathway is also consistent with the data. Figures 8c shows that as $`R_2`$ decreases, the backness pathway progressively advances into the pre-existing front until it occupies not only the back but also the front of the cell. This is consistent with the data shown in Figure 3a.
Simulations in which $`r`$ is held at 1.0 and $`R_3`$ is progressively changed give analogous results (simulations not shown). As $`R_3`$ increases, the frontness pathways occupy a progressively smaller portion of the cell until there is no frontness throughout the cell. The simulations did not yield multiple peaks corresponding to the multiple pseudopods shown in Figure 4b. Numerical calculations of the critical wavenumber, $`k_0`$, show that it increases as $`R_3`$ decreases, which is the correct trend. However, the steady state ceases to exist before $`k_0`$ can exceed 2.
To investigate the existence of steady states corresponding to chemorepulsion, we chose parameter values satisfying (10) so that $`U_3`$ prevails over $`U_2`$ in the absence of the inhibitor (column labeled “$`U_3`$ wins” of Table 1). Here, the coexistence steady state, $`E_{111}`$, undergoes a Turing bifurcation at $`r0.65`$ and $`k=1`$ (Figures 11b). The non-homogeneous steady state that emerges from this bifurcation is such that $`U_2`$ and $`U_3`$ are once again out of phase,<sup>3</sup><sup>3</sup>3It is shown in Appendix B that regardless of the parameter values, $`U_2`$ and $`U_3`$ will always be out-of-phase in the non-homogeneous steady state bifurcating from $`E_{111}`$. but it is $`U_3`$, instead of $`U_2`$, that is in phase with $`U_1`$ (Figure 9b). This non-homogeneous steady state represents the *chemorepulsion* steady state, since it approximates the steady state that would be obtained in the presence of a receptor gradient, $`r(x)=1+ϵ\mathrm{cos}(\pi x),0<ϵ1`$, and the profile of the frontness pathways are, in this case, out of phase with the active receptor distribution. Given the symmetry of the equations, it is not surprising that the model admits a chemorepulsion steady state. For, if choose the parameters corresponding to the column “$`U_2`$ wins” in Table 1, but interchange $`u_2`$ and $`u_3`$, the chemoattraction steady state in Figure 7b turns into a chemorepulsion steady state. However, the analysis of the model suggests a mechanism for the transition. Specifically, the conditions (910) imply that transitions from chemorepulsion to chemattraction, such as those observed in neurons , can be triggered by pharmacological agents that selectively inhibit the backness pathway so that the frontness pathway becomes intrinsically superior in the absence of the inhibitor.
## 3 Discussion
We have shown that a model based on mutual inhibition of the frontness and backness pathways can yield spontaneous polarization involving spatial segregation of the two pathways. The frontness and backness pathways are mutually incompatible in the absence of the inhibitor (eq. 7), but coexist in the presence of the inhibitor because it suppresses the growth, and hence, the mutual incompatibility, of the frontness and backness pathways. Since the inhibitor is diffusible, it can suppress the nominal incompatibility manifested at low receptor levels. However, it fails to achieve this goal at high receptor levels, resulting in spatial segregation of the frontness and backness pathways (see Figures 1 and 7).
Since the model takes due account of the backness pathways, we could explore the effect of frontness inhibition or activation on the backness pathways, and the effect of backness modification on the frontness and backness pathways. The simulations in Section 2 yielded results that are in qualitative agreement with the data shown in Figures 14. These results are beyond the scope of activator-inhibitor or activator-substrate models of gradient sensing since they do not account for the backness pathways.
The particular manner in which the two pathways segregate depends on their behavior in the absence of the tempering influence of the inhibitor. If the frontness pathways dominate over the backness pathways, their spatial segregation is consistent with chemoattraction. Conversely, if the backness pathways dominate, the emergent spatial pattern corresponds to chemorepulsion. The model therefore suggests that transitions between chemoattraction and chemorepulsion can be provoked by agents that shift the balance of power between the frontness and backness pathways.
The model proposed here is the simplest possible model of spontaneous polarization driven by mutual inhibition. By this, we mean the following
1. A model containing a smaller number of variables will not yield spontaneous polarization.
2. The model is simpler than any other model in which the symmetry-breaking is driven entirely by mutual inhibition.
In what follows, we justify these conclusions, since they are required to explain certain points raised above.
To see that a model containing 2 variables will not yield spontaneous polarization driven by mutual inhibition, let $`u=[u_1,u_2]^t`$ be a homogeneous steady state of a two-component reaction-diffusion system
$$\frac{u_i}{t}=D_i\frac{u_i}{x^2}+f_i(u_1,u_2),i=1,2$$
and let $`J(u)`$ denote the Jacobian at $`u`$. Then, $`u`$ can undergo a Turing instability only if
$$J_{11}+J_{22}<0,J_{11}J_{22}J_{12}J_{21}>0,D_1J_{22}+D_2J_{11}>0.$$
(11)
The first and third conditions imply that $`J_{11}J_{22}<0`$, i.e., one of the components, say $`U_1`$, must inhibit its own synthesis ($`J_{11}<0`$), and the other, say $`U_2`$, must activate its own synthesis ($`J_{22}>0`$). Now, if the two components inhibit each other, then $`J_{12},J_{21}<0`$, so that $`J_{12}J_{21}>0`$ and $`J_{11}J_{22}J_{12}J_{21}<0`$, which violates the second condition for Turing stability. Thus, a two-component model in which each component inhibits the other component cannot yield spontaneous polarization. Indeed, such a model can yield a Turing instability only if $`J_{12}J_{21}<0`$, i.e., the interaction between $`U_1`$ and $`U_2`$ is asymmetric. If $`U_1`$ inhibits $`U_2`$, then $`U_2`$ must activate $`U_1`$; this corresponds to the *activator-inhibitor* model; conversely, if $`U_1`$ activates $`U_2`$, then $`U_2`$ must inhibit $`U_1`$, which corresponds to the *activator-substrate* model. No other cross-interactions are admissible. It follows that if spontaneous polarization is to be driven by mutual inhibition, there must be at least three components. Thus, we were led to modify the Bourne model by postulating the existence of an additional component, namely, a diffusible inhibitor.
To show that the model is simpler than any other model in which symmetry-breaking is driven entirely by mutual inhibition, we appeal to the fact that in a general reaction-diffusion system, a homogeneous steady state cannot become Turing instable unless it contains an *unstable subsytem* .<sup>4</sup><sup>4</sup>4A *subsystem* refers to any proper subset of the model variables. For example, the two-component model discussed above has two subsystems, $`\{U_1,U_2\}`$. A subsystem is said to be *unstable* if it becomes unstable as soon as the remaining components of the system are somehow rendered constant. In other words, without the stabilizing effects of the other variables, the subsystem diverges from its values at the homogeneous steady state. In the 2-component model, the activator is an unstable subsystem: If the inhibitor or substrate level is somehow held constant, the activator diverges from the homogeneous steady state due to positive feedback ($`J_{22}>0`$). As shown above, the existence of such an unstable subsystem is necessary for the full system to undergo a Turing instability. The three-component model considered here has three 1-component subsystems, $`\{U_1,U_2,U_3\}`$, and three 2-component subsystems, $`\{U_1U_2,U_2U_3,U_1U_3\}`$. It turns out that in the neighborhood of the non-trivial steady state, $`E_{111}`$, all three 1-component subsystems are stable, i.e. $`J_{11},J_{22},J_{33}<0`$ (see 13). In other words, the self-interactions of the individual components are stabilizing. Notably, $`U_2`$ and $`U_3`$ are stable subsystems even though their synthesis was assumed to be autocatalytic. This is because the kinetics of their synthesis are such that the feedback is positive only when $`u_2`$ and $`u_3`$ are small. However, in the neighborhood of $`E_{111}`$, $`u_2`$ and $`u_3`$ are so large that the self-limiting effect dominates, and the feedback is, in fact, negative. Now, since the 1-component subsystems are stable, $`E_{111}`$ can be Turing unstable only if at least one of the 2-component subsystems is unstable, i.e., the cross-interaction of at least one 2-component system is destabilizing. Such destabilizing cross-interactions can occur only if the interaction between the 2 components is mutually inhibitory or synergistic, i.e.,
$$J_{12}J_{21}>0\text{ or }J_{13}J_{31}>0\text{ or }J_{23}J_{32}>0.$$
The first two conditions cannot be satisfied since the interaction between $`U_1`$ and each of the two activators, $`U_2`$ and $`U_3`$ is neither mutually inhibitory nor mutually synergistic — $`U_1`$ inhibits $`U_2`$ and $`U_3`$, but $`U_2`$ and $`U_3`$ activate $`U_1`$. Thus, $`\{U_1,U_2\}`$ and $`\{U_1,U_3\}`$ are stable subsystems, and $`\{U_2,U_3\}`$ is the only unstable subsystem in the model. Consequently, the spatial segregation of the frontness and backness pathways is driven *entirely* by mutual inhibition between $`U_2`$ and $`U_3`$ — positive feedback plays no role since $`J_{11},J_{22},J_{33}<0`$.
It should be noted that reciprocal distribution of the frontness and backness pathways (specifically, PI3K and PTEN) is a central feature of the local-excitation-global-inhibition model proposed by Iglesias and coworkers . However, the mechanism for spatial segregation of PI3K and PTEN is such that polarization occurs only in the presence of a chemoattractant gradient. Indeed, according to this model
1. Exposure to chemoattractant generates binding sites for PI3K and PTEN at the membrane.
2. The rate of generation of binding sites for PI3K and PTEN is proportional to $`r^2`$ and $`r`$, respectively, where $`r`$ denotes the active receptor concentration.
3. The concentration of PI3P is determined by the relative concentrations of membrane-bound PI3K and PTEN — the higher the PI3K:PTEN ratio, the larger the concentration of PI3P.
It follows from (2) that when a cell is exposed to a gradient, the leading edge develops a high PI3K:PTEN ratio, while the trailing edge has a high PTEN:PI3K ratio. Then, (3) implies that the PI3P concentration at the leading edge is higher than that at the trailing edge. Importantly, the polarized distribution of PI3P arises only in the presence of a gradient. For, if $`r`$ is constant, the PI3K:PTEN ratio, and hence, the concentration of PI3P, is constant all over the cell membrane. Thus, the local-excitation-global-inhibition model cannot capture the spontaneous polarization observed in the absence of gradients.
Although our model captures spontaneous polarization driven by mutual inhibition of frontness and backness pathways, it is missing two important features of the experimental data. First, it is known that both negative *and* positive feedback operate in chemoattractant-mediated polarization (see Figure 5a and refs. ). Although we assumed the existence of positive feedback in the model, it played no role in the spontaneous polarization since $`J_{22},J_{33}<0`$ in the neighborhood of $`E_{111}`$. Second, the model does not account for adapation, the process that enables the cell to ultimately return to the homogeneous steady state even in the presence of chemoattractant . A modified model taking due account of positive feedback and long-term adaptation is currently under investigation.
Two of the model variables ($`U_2,U_3`$) correspond to the frontness and backness pathways of the Bourne scheme. However, the model hypothesizes the existence of an additional variable, namely, a diffusible inhibitor required to suppress the mutual incompatibility of the frontness and backness pathways. It is, by no means, necessary that this diffusible component be an inhibitor. It could just as well be a diffusible *substrate* that is required for the synthesis of the two activators, but is consumed in process of activator synthesis. This would reverse the signs of $`J_{12},J_{13},J_{21},J_{31}`$, but the key inequalities in the above arguments, $`J_{13}J_{31},J_{12}J_{21}<0`$, would be preserved. A diffusible inhibitor was assumed only because there is some evidence of their existence.
1. Luo and coworkers have shown that cytosolic inositol phosphates, whose levels surge immediately after chemoattractant stimulation, inhibit the frontness pathways in both *Dictyostelium* and neutrophil-like HL-60 cells .
2. Nimnual *et al* have shown that Rac-mediated production of reactive oxygen species downregulates the backness component, Rho .
It remains to be seen if these diffusible inhibitors play a critical role in the maintenance of polarity, as required by the diffusible inhibitor of the proposed model.
## 4 Conclusions
We have formulated a model that provides a mathematical realization of the Bourne scheme for spontaneous polarization by mutual inhibition of frontness and backness pathways . The model predicts several experimentally observed features that are outside the scope of prevailing models of gradient sensing.
1. Mutual inhibition of the frontness and backness pathways plays a critical role in generating the spontaneous polarization.
2. If the frontness pathway is suppressed, the backness pathways occupy a progressively larger proportion of the cell. Conversely, if the backness pathway is suppressed, the frontness pathways invade and occupy the previous back of the cell.
3. Depending on the parameter values, the model yields steady states corresponding to both chemoattraction and chemorepulsion. Analysis of the model suggests that chemorepulsion-to-chemoattraction transitions observed in neurons can be triggered by agents that inhibit the backness pathways more than the frontness pathways.
The model provides a useful starting point for formulating models that account for the positive and negative feedback effects, both of which have been shown to play a role in gradient sensing.
#### Acknowledgments
I would like to thank Prof. D. A. Lauffenburger for helpful comments on the manuscript.
## Appendix A Existence/stability of homogeneous steady states
The homogeneous steady states of (46) satisfy the equations
$`u_1+a_{12}u_2+a_{13}u_3`$ $`=0`$
$`\left(ra_{21}u_1a_{22}u_2a_{23}u_3\right)u_2`$ $`=0`$
$`\left(ra_{31}u_1a_{32}u_2a_{33}u_3\right)u_3`$ $`=0.`$
It follows that there are 4 possible steady states: $`u_1=u_2=u_3=0`$, denoted $`E_{000}`$; $`u_1,u_2>0,u_3=0`$, denoted $`E_{110}`$; $`u_1,u_3>0,u_2=0`$, denoted $`E_{101}`$; and $`u_1,u_2,u_3>0`$, denoted $`E_{111}`$. In what follows, we show the existence and stability criteria for all 4 steady states. To this end, it is convenient to note that the Jacobian at any point, $`u_1,u_2,u_3`$ is
$$J=\left[\begin{array}{ccc}1& a_{12}& a_{13}\\ \rho _2u_2a_{21}& \rho _2u_2a_{22}+\rho _2f_2& \rho _2u_2a_{31}\\ \rho _3u_3a_{31}& \rho _3u_3a_{32}& \rho _3u_3a_{33}+\rho _3f_3\end{array}\right]$$
where
$`f_2`$ $`ra_{21}u_1a_{22}u_2a_{23}u_3`$
$`f_3`$ $`ra_{31}u_1a_{32}u_2a_{33}u_3`$
1. $`E_{000}=[0,0,0]^t`$: It is evident that this steady state always exists. The Jacobian yields no information regarding the stability since it is singular at $`E_{000}`$. However, we can infer that it is unstable for all $`r>0`$ by observing that the $`u_2`$-axis is an invariant manifold. The motion along this invariant manifold is given by
$$\frac{du_2}{dt}=\rho _2(ra_{22}u_2)u_2,$$
which drives the system away from $`E_{000}`$ for all $`r>0`$.
2. $`E_{110}=[u_1,u_2,0]^t,u_1,u_2>0`$: It is easy to check that this steady state exists for all $`r>0`$, and is given by
$$u_3=0,u_2=\frac{r}{\alpha _{22}},u_1=a_{12}u_2,$$
where $`\alpha _{22}a_{22}+a_{21}a_{12}`$, and the Jacobian at $`E_{110}`$ is
$$\left[\begin{array}{ccc}1& a_{12}& a_{13}\\ \rho _2u_2a_{21}& \rho _2u_2a_{22}& \rho _2u_2a_{31}\\ 0& 0& \rho _3(ra_{31}u_1a_{32}u_2)\end{array}\right].$$
(12)
It follows that one of the eigenvalues is
$$\lambda _1=\rho _3(ra_{31}u_1a_{32}u_2)=\rho _3r\left(1\frac{\alpha _{32}}{\alpha _{33}}\right),$$
where $`\alpha _{32}a_{32}+a_{31}a_{12}`$. The other two eigenvalues satisfy
$`\lambda _2+\lambda _3`$ $`=1\rho _2a_{22}ru_2<0,`$
$`\lambda _2\lambda _3`$ $`=\rho _2a_{22}ru_2+\rho _2u_2a_{21}a_{12}>0,`$
so that the real parts of $`\lambda _2,\lambda _3`$ are negative. Hence, $`E_{110}`$ is stable if and only if $`\beta _{32}\alpha _{32}/\alpha _{22}>1`$.
3. $`E_{101}=[u_1,0,u_3]^t,u_1,u_3>0`$: The calculations are analogous to those for $`E_{101}`$.
4. $`E_{111}=[u_1,u_2,u_3]^t,u_1,u_2,u_3>0`$: This steady state satisfies the equations
$`u_1`$ $`=a_{12}u_2+a_{13}u_3`$
$`r`$ $`=\alpha _{22}u_2+\alpha _{23}u_3`$
$`r`$ $`=\alpha _{32}u_2+\alpha _{33}u_3.`$
It follows that $`E_{111}`$ exists for all $`r>0`$ if and only if $`[1,1]^t`$ lies between $`[\alpha _{22},\alpha _{32}]^t`$ and $`[\alpha _{23},\alpha _{33}]^t`$, i.e., $`\beta _{23}\alpha _{23}/\alpha _{33},\beta _{32}\alpha _{32}/\alpha _{22}`$ satisfy
$$\beta _{23},\beta _{32}<1\text{ or }\beta _{23},\beta _{32}>1.$$
In both cases, the steady state is given by
$`u_2`$ $`={\displaystyle \frac{1}{\alpha _{22}}}{\displaystyle \frac{1\beta _{23}}{1\beta _{23}\beta _{32}}}r,`$
$`u_3`$ $`={\displaystyle \frac{1}{\alpha _{33}}}{\displaystyle \frac{1\beta _{32}}{1\beta _{23}\beta _{32}}}r,`$
$`u_1`$ $`=a_{12}u_2+a_{13}u_3.`$
Evidently, $`u_1,u_2,u_3`$ increase linearly with $`r`$.
It turns out that $`E_{111}`$ is unstable if $`\beta _{23},\beta _{32}>1`$, and stable for sufficiently small $`r>0`$ if $`\beta _{23},\beta _{32}>1`$. To see this, observe that the Jacobian at $`E_{111}`$ is
$$\left[\begin{array}{ccc}1& a_{12}& a_{13}\\ \rho _2u_2a_{21}& \rho _2u_2a_{22}& \rho _2u_2a_{31}\\ \rho _3u_3a_{31}& \rho _3u_3a_{32}& \rho _3u_3a_{33}\end{array}\right],$$
(13)
so that
$`\text{tr }J`$ $`=1\rho _2a_{22}u_2\rho _3a_{33}u_3,`$
$`detJ`$ $`=\rho _2\rho _3u_2u_3\alpha _{22}\alpha _{23}(1\beta _{23}\beta _{32}),`$
$`\mathrm{\Sigma }J`$ $`=\rho _2\rho _3u_2u_3(a_{22}a_{33}a_{23}a_{32})+\rho _2u_2\alpha _{22}`$
$`+\rho _3u_3\alpha _{33}.`$
If $`\beta _{23},\beta _{32}>1`$, then $`detJ>0`$, and $`E_{111}`$ is unstable. Hence, $`E_{111}`$ is stable only if $`\beta _{23},\beta _{32}<1`$. We are particularly interested in the case when $`a_{22}a_{33}a_{23}a_{32}>0`$, since this is necessary for Turing instability (see below). Under these conditions, $`E_{111}`$ is stable for all sufficiently small $`r>0`$, since $`\text{tr }J<0`$, $`detJ<0`$ for all $`r>0`$, and $`(\text{tr }J)(\mathrm{\Sigma }J)detJ<0`$ for all sufficiently small $`r>0`$ (Figure 10).
## Appendix B Turing instability of homogeneous steady states
### B.1 Conditions for Turing instability
A homogeneous steady state, $`u=[u_1(r),u_2(r),u_3(r)]^t`$, is Turing unstable if it is stable in the absence of diffusion and unstable in the presence of diffusion. Now the stability of the homogeneous steady state, $`u`$, in the presence of diffusion is determined by the linearized equation
$$\frac{v}{t}=D\frac{v}{x^2}+J(u)v$$
(14)
where $`v`$ denotes the deviation from $`u`$, $`J(u)`$ denotes the Jacobian at $`u`$, and $`D=\text{diag}(d_1,d_2,d_3)`$. Let $`\mu _k`$ and $`\varphi _k(x)`$ denote the eigenvalues and eigenfunctions of the differential operator, $`\frac{d^2}{dx^2}`$, with Neumann boundary conditions, i.e.,
$$\mu _k=(k\pi )^2,\varphi _k(x)=\mathrm{cos}(k\pi x),k=0,1,2,\mathrm{}$$
If the steady state is perturbed along any eigenfunction, $`\varphi _k(x)`$, the subsequent evolution of the perturbation is of the form $`v(t,x)=s_k(t)\varphi _k(x)`$ where $`s_k(t)^3`$. Substituting in (14) yields
$$\frac{ds_k}{dt}=C_ks_k,C_k(u)J(u)(k\pi )^2D.$$
It follows that the homogeneous steady state is stable if and only if the eigenvalues of $`C_k(u)`$ have negative real parts for all $`k0`$. It is Turing unstable if and only if
1. The eigenvalues of $`C_0(u)=J(u)`$ have negative real parts (which ensures that the homogeneous steady state is stable in the absence of diffusion).
2. The eigenvalues of $`C_k(u)`$ have positive real part for some $`k>0`$ (which ensures that the homogeneous steady state is unstable in the presence of diffusion).
The Turing bifurcation occurs at the critical value, $`r=r_0>0`$, such that exactly one of the eigenvalues of $`C_k(u)`$ becomes zero, i.e.,
$$detC_k(u)=0,\frac{d}{dk^2}detC_k(u)=0$$
for some $`k>0`$ (Figure 11).
### B.2 Turing instability of the homogeneous steady states
It follows from the definition of $`C_k`$ that
$`detC_k(u)`$ $`=\mathrm{}_{123}(d_1\mathrm{}_{23}+d_2\mathrm{}_{13}+d_3\mathrm{}_{13})k^2`$
$`+(d_1d_2\mathrm{}_3+d_2d_3\mathrm{}_1+d_1d_3\mathrm{}_2)k^4`$
$`(d_1d_2d_3)k^6`$ (15)
where $`\mathrm{}_i`$, $`\mathrm{}_{ij},ij`$, $`\mathrm{}_{123}`$ are the determinants of the 1-, 2-, and 3-dimensional subsystems of $`J`$, i.e., $`\mathrm{}_i=J_{ii}(u)`$,
$`\mathrm{}_{ij}`$ $`=det\left[\begin{array}{cc}J_{ii}& J_{ij}\\ J_{ji}& J_{jj}\end{array}\right],ij`$
and $`\mathrm{}_{123}=detJ(u)`$. These results are sufficient for investigating the possibility that the homogeneous steady states undergo a Turing bifurcation.
We consider each of the homogeneous steady states.
1. $`E_{000}=[0,0,0]^t`$: Since $`E_{000}`$ is unstable even in the absence of diffusion, it cannot entertain Turing instability.
2. $`E_{110}=[u_1,u_2,0]^t`$: This steady state is stable in the absence of diffusion provided $`\beta _{32}>1`$. Thus, it can undergo a Turing instability if and only if there is an $`r>0`$ such that (15) has a positive root. It turns out that such a root does not exist. To see this, observe that since $`E_{110}`$ is stable in the absence of diffusion, $`\mathrm{}_{123}=detJ<0`$. Inspection of the Jacobian (12) shows that the $`\mathrm{}_i`$’s are always negative, and the $`\mathrm{}_{ij}`$’s are always positive. Hence, all the coefficients of the polynomial (15) are negative for all $`r>0`$. Such a polynomial cannot have positive roots ($`k>0`$). We conclude that $`E_{110}`$ cannot undergo a Turing instability.
3. $`E_{101}=[u_1,0,u_3]^t`$: This steady state is stable in the absence of diffusion provided $`\beta _{23}<1`$. An argument analogous to that for $`E_{110}`$ shows that it cannot undergo a Turing instability.
4. $`E_{111}=[u_1,u_2,u_3]^t`$: This steady state is stable in the absence of diffusion provided $`\beta _{23},\beta _{32}<1`$ and $`r`$ is sufficiently small. Under these conditions, $`\mathrm{}_{123}<0`$. Inspection of the Jacobian (13) shows that no matter what the value of $`r`$, all the $`\mathrm{}_i`$’s are negative, and $`\mathrm{}_{12},\mathrm{}_{13}`$ are positive. Hence, $`E_{111}`$ can undergo a Turing instability only if $`\mathrm{}_{23}<0`$. This condition is not only necessary, but also sufficient. For, if $`d_2`$ and $`d_3`$ are sufficiently small compared to $`d_1`$, the coefficient of $`k^2`$ is approximately $`d_1\mathrm{}_{23}<0`$, and the polynomial (15) has a positive root.
The non-homogeneous steady state that bifurcates from $`E_{111}`$ is such that $`U_2`$ and $`U_3`$ are out-of-phase. To see this, it suffices to observe that to a first degree of approximation, the non-homogeneous steady state near the Turing bifurcation point has the form $`v_0\mathrm{cos}(k_0x)`$, where $`k_0`$ be the critical wavenumber, and $`v_0`$ is the eigenvector corresponding to the zero eigenvalue of $`C_{k_0}`$, i.e., $`C_{k_0}v_0=0`$. The null space of the rank-2 matrix, $`C_{k_0}`$, is spanned by the vector, $`[c_1,c_2,1]`$, where
$`c_2`$ $`{\displaystyle \frac{C_{k_0,22}C_{k_0,12}C_{k_0,21}/C_{k_0,11}}{C_{k_0,23}C_{k_0,21}C_{k_0,13}/C_{k_0,11}}}.`$
Substitution of the coefficients of $`C_{k_0}`$ shows that $`c_2`$ is always negative, which implies that $`U_2`$ and $`U_3`$ are out-of-phase. |
warning/0507/hep-ph0507215.html | ar5iv | text | # A Closer Look at the Analysis of NLL BFKL
## Introduction
The Balitsky–Fadin–Kuraev–Lipatov$`^\mathrm{?}`$ (BFKL) framework systematically resums the class of logarithms originating from the kinematics that dominate the total cross section in the Regge limit of scattering amplitudes, where the centre of mass energy $`\sqrt{\widehat{s}}`$ is large and the momentum transfer $`\sqrt{\widehat{t}}`$ is fixed. In this limit the scattering of two gluons $`p_Ap_Bp_A^{}p_B^{}`$ will be dominated by multi–particle production leading to final states described by momenta $`k_0=p_A^{},k_1,\mathrm{},k_n,k_{n+1}=p_B^{}`$ satisfying
$`s2k_{i1}k_it_i=q_i^2,q_i=p_A{\displaystyle \underset{r_0}{\overset{i1}{}}}k_r,{\displaystyle \underset{i=1}{\overset{n+1}{}}}s_i=s{\displaystyle \underset{i=1}{\overset{n}{}}}𝐤_i^2,k_{}^2=𝐤^2,|k_i||p_A^{}|`$ (1)
The Regge limit is therefore suitable for describing the production of multiple hard partons from e.g. gluon scattering (and in fact the large-rapidity limit of any process that includes a $`t`$-channel gluon exchange). We will in this talk focus entirely on processes within the multi-Regge kinematics of Eq. (1). In this limit the partonic cross section can be approximated by
$`\widehat{\sigma }(\mathrm{\Delta })`$ $`=`$ $`{\displaystyle \frac{d^2𝐤_a}{2\pi 𝐤_a^2}\frac{d^2𝐤_b}{2\pi 𝐤_b^2}\mathrm{\Phi }_A(𝐤_a)f(𝐤_a,𝐤_b,\mathrm{\Delta })\mathrm{\Phi }_B(𝐤_b)},`$ (2)
where $`\mathrm{\Phi }_{A,B}`$ are the impact factors characteristic of the particular scattering process, and $`f(𝐤_a,𝐤_b,\mathrm{\Delta })`$ is the gluon Green’s function describing the interaction between two Reggeised gluons exchanged in the $`t`$–channel with transverse momenta $`𝐤_{a,b}`$, spanning a rapidity interval of length $`\mathrm{\Delta }`$. The leading and next-to-leading logarithmic contributions to this gluon Green’s function can be resummed by solving the BFKL equation to the required accuracy
$`\omega f_\omega (𝐤_a,𝐤_b)`$ $`=`$ $`\delta ^{(2+2ϵ)}\left(𝐤_a𝐤_b\right)+{\displaystyle d^{2+2ϵ}𝐤𝒦(𝐤_a,𝐤+𝐤_a)f_\omega (𝐤+𝐤_a,𝐤_b)},`$ (3)
where $`w`$ is the Mellin-conjugated variable to $`\mathrm{\Delta }`$, and the BFKL kernel $`𝒦(𝐤_a,𝐤+𝐤_a)`$ is presently known to next-to-leading logarithmic accuracy.
## 1 Solutions of the BFKL equation
The solution to integral equations of the form in Eq. (3) can be written in terms of the eigenfunctions $`\varphi _i(k_a)`$ and eigenvalues $`\lambda _i`$ as
$`f_\omega (k_a,k_b)={\displaystyle \underset{i}{}}{\displaystyle \frac{\varphi _i(k_a)\varphi _i^{}(k_b)}{\omega \lambda _i}}`$ (4)
leading to
$`f(k_a,k_b,\mathrm{\Delta })={\displaystyle \underset{i}{}}{\displaystyle \frac{1}{2\pi i}e^{\mathrm{\Delta }\lambda _i}\varphi _i(k_a)\varphi _i^{}(k_b)}`$ (5)
### 1.1 Leading Logarithmic Accuracy
At leading logarithmic accuracy the BFKL kernel is conformal invariant, since the running of the coupling only enters at higher logarithmic orders. The eigenfunctions of the angular averaged kernel are of the form $`k_{}^{2}{}_{}{}^{(\gamma 1)},\gamma =1/2+i\nu `$, which means that to this accuracy, the BFKL evolution can be solved analytically, with the transverse momentum of emitted gluons integrated to infinity, by analysing the Mellin transform of the kernel. One finds
$`{\displaystyle \mathrm{d}^{D2}𝐤𝒦^{\mathrm{LL}}(𝐤_a,𝐤)\left(𝐤^2\right)^{\gamma 1}}={\displaystyle \frac{\alpha _sN}{\pi }}\chi ^{\mathrm{LL}}(\gamma )\left(𝐤_a^2\right)^{\gamma 1},`$ (6)
with $`N`$ being the number of colours and
$`\chi ^{\mathrm{LL}}(\gamma )=2\psi (1)\psi (\gamma )\psi (1\gamma ),\psi (\gamma )=\mathrm{\Gamma }^{}(\gamma )/\mathrm{\Gamma }(\gamma ).`$ (7)
Since both the eigenfunctions and eigenvalues are known, the angular averaged gluon Green’s function can now be obtained according to Eq. (5) as
$`\overline{f}(k_a,k_b,\mathrm{\Delta })={\displaystyle \frac{1}{k_b^2}}{\displaystyle _{\frac{1}{2}i\mathrm{}}^{\frac{1}{2}+i\mathrm{}}}{\displaystyle \frac{\mathrm{d}\gamma }{2\pi i}}e^{\mathrm{\Delta }\omega ^{\mathrm{LL}}(\gamma )}\left({\displaystyle \frac{k_b^2}{k_a^2}}\right)^\gamma ,`$ (8)
where
$`\omega ^{\mathrm{LL}}(\gamma ){\displaystyle \mathrm{d}^{D2}𝐤𝒦^{\mathrm{LL}}(𝐤_a,𝐤)\left(\frac{𝐤^2}{𝐤_a^2}\right)^{\gamma 1}}={\displaystyle \frac{\alpha _s(k_a^2)N}{\pi }}\chi ^{\mathrm{LL}}(\gamma ).`$ (9)
We stress that at LL the coupling is formally fixed, and so the regularisation scale is completely arbitrary.
In Fig. 1 we have plotted $`\omega ^{\mathrm{LL}}(\frac{1}{2}+i\nu )`$ for $`\alpha _s=0.2`$ and it is seen that there is a maximum at $`\nu =0`$. Therefore, the behaviour of the gluon Green’s function in the limit of $`\mathrm{\Delta }\mathrm{}`$ is determined by the value of $`\omega ^{\mathrm{LL}}(\frac{1}{2})=4\mathrm{ln}2\alpha _sN/\pi `$. A saddle point approximation based on the second order Taylor polynomial around $`\nu =0`$ will correctly describe the asymptotic exponential growth in $`\mathrm{\Delta }`$, since the polynomial attains the correct value at the maximum.
### 1.2 Next-to-Leading Logarithmic Accuracy
When trying to extend this analysis to NLL accuracy one is immediately faced with the complications introduced by the breaking of the conformal invariance by the running coupling terms. This effect will necessarily change the eigenfunctions, and thus far the eigenfunctions for the full NLL kernel in QCD have not been constructed. Traditionally, the kernel at NLL has been studied using the projection on the Born level eigenfunctions as in Eq. (9). One finds $`^\mathrm{?}`$
$`\begin{array}{cc}\hfill \omega ^{\mathrm{NLL}}(\gamma )& {\displaystyle \mathrm{d}^{D2}𝐤𝒦^{\mathrm{NLL}}(𝐤_a,𝐤)\left(\frac{𝐤^2}{𝐤_a^2}\right)^{\gamma 1}}\hfill \\ \hfill =& {\displaystyle \frac{\alpha _s(𝐤_a^2)N}{\pi }}\left(\chi ^{\mathrm{LL}}(\gamma )+\chi ^{\mathrm{NLL}}(\gamma ){\displaystyle \frac{\alpha _s(𝐤_a^2)N}{\pi }}\right)\hfill \end{array}`$ (10)
with
$`\begin{array}{cc}\hfill \chi ^{\mathrm{NLL}}(\gamma )=& {\displaystyle \frac{1}{4}}[({\displaystyle \frac{11}{3}}{\displaystyle \frac{2}{3}}{\displaystyle \frac{n_f}{N}}){\displaystyle \frac{1}{2}}(\chi ^{\mathrm{LL}}(\gamma )\psi ^{}(\gamma )+\psi ^{}(1\gamma ))\hfill \\ & 6\zeta (3)+{\displaystyle \frac{\pi ^2\mathrm{cos}(\pi \gamma )}{\mathrm{sin}^2(\pi \gamma )(12\gamma )}}\left(3+\left(1+{\displaystyle \frac{n_f}{N^3}}\right){\displaystyle \frac{2+3\gamma (1\gamma )}{(32\gamma )(1+2\gamma )}}\right)\hfill \\ & ({\displaystyle \frac{67}{9}}{\displaystyle \frac{\pi ^2}{3}}{\displaystyle \frac{10}{9}}{\displaystyle \frac{n_f}{N}})\chi ^{\mathrm{LL}}\psi ^{\prime \prime }(\gamma )\psi ^{\prime \prime }(1\gamma ){\displaystyle \frac{\pi ^3}{\mathrm{sin}(\pi \gamma )}}+4\varphi (\gamma )],\hfill \end{array}`$ (11)
where
$`\begin{array}{cc}\hfill \varphi (\gamma )=& {\displaystyle _0^1}{\displaystyle \frac{\mathrm{d}x}{1+x}}(x^{\gamma 1}+x^\gamma ){\displaystyle _x^1}{\displaystyle \frac{\mathrm{d}t}{t}}\mathrm{ln}(1t)\hfill \\ \hfill =& {\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}(1)^n\left[{\displaystyle \frac{\psi (n+1+\gamma )\psi (1)}{(n+\gamma )^2}}+{\displaystyle \frac{\psi (n+2\gamma )\psi (1)}{(n+1\gamma )^2}}\right].\hfill \end{array}`$ (12)
An approximation to the gluon Green’s function at NLL can then be constructed by use of $`\omega ^{\mathrm{NLL}}`$ in place of $`\omega ^{\mathrm{LL}}`$ in Eq. (8). We have in Fig. 2 plotted the real and imaginary part of $`\omega ^{\mathrm{NLL}}`$ compared with $`\omega ^{\mathrm{LL}}`$ for $`\alpha _s=0.2`$. The double hump structure of the real part of $`\omega ^{\mathrm{NLL}}`$ is potentially a disaster for the gluon Green’s function. At asymptotically large $`\mathrm{\Delta }`$ the behaviour of the gluon Green’s function is determined by the position and value of the maxima of $`\omega ^{\mathrm{NLL}}`$ only. Since there are two such distinct maxima located at $`\gamma _1=1/2i\nu _0,\gamma _2=1/2+i\nu _0`$, the asymptotic estimate of the NLL gluon Green’s function based on the LL eigenfunctions will have an oscillatory behaviour in $`\mathrm{ln}k_a/k_b`$. This is a problem of matching to the DGLAP evolution, and is a problem strictly outside the Regge kinematics.
The initial observation of the large NLL corrections was based on the difference $`\omega ^{\mathrm{LL}}\omega ^{\mathrm{NLL}}`$ evaluated at $`\nu =0`$. Indeed, for reasonable values of the coupling, the real part of $`\omega ^{\mathrm{NLL}}`$ is even negative at this point. However, this is not what determines the intercept. If indeed the solution of the BFKL equation at NLL was obtained using $`\omega ^{\mathrm{NLL}}`$ of Eq. (10), then the asymptotic intercept at NLL would again be determined by the maximum value attained by the real part of $`\omega ^{\mathrm{NLL}}(\gamma )`$ along the contour $`\gamma =1/2+i\nu `$. We see from Fig. 2 that this maximum value is roughly halved compared to the LL asymptotic intercept.
However, even within the Regge kinematics, this analysis leads to severe problems in the very limit, where the resummed logarithmic terms are meant to dominate the scattering matrix. The non-zero imaginary part of $`\omega ^{\mathrm{NLL}}(\frac{1}{2}+i\nu _0)`$ leads to oscillations with increasing rapidity for any choice of $`k_a`$ and $`k_b`$! This clearly signals a breakdown of the approach in the very limit it is meant to describe.
The solution to this apparent problem is the realisation that what one has studied with this method is indeed not the true solution to the BFKL equation at NLL. The LL eigenfunctions are *not* the eigenfunctions at NLL (for non-zero $`\beta _0`$). Indeed, the only contribution to the troublesome imaginary part of $`\omega ^{\mathrm{NLL}}(\gamma )`$ for $`\gamma =1/2+i\nu `$ stems from the term $`\psi ^{}(\gamma )+\psi (1\gamma )`$ in Eq. (11), which contributes to $`\omega ^{\mathrm{NLL}}`$ with a factor proportional to $`\beta _0`$ (that is, it vanishes in the limit where the LL eigenfunctions diagonalises the NLL kernel). It was observed already in Ref.$`^\mathrm{?}`$ that this part of the NLL corrections is the only one that is not symmetric under $`\gamma 1\gamma `$, and that if one expands the kernel on the LL eigenfunctions rescaled by a square root of the coupling, i.e. $`(k^2)^{\gamma 1}\left(\frac{\alpha _s(k^s)}{\alpha _s(\mu ^2)}\right)^{1/2}`$ then these and only these terms would disappear from Eq. (11). What was perhaps not realised is that since this is the only contribution to the imaginary part of $`\omega ^{\mathrm{NLL}}`$, this would simultaneously cure the problem of oscillations within the Regge-limit. It should be emphasised though that these rescaled functions still are not the true eigenfunctions at NLL, but it is straightforward to check numerically that the “off-diagonal elements” of $`\omega ^{\mathrm{NLL}}`$ (i.e. those obtained with a different $`\gamma `$ for $`k_a`$ and $`k`$ in Eq. (10)) are far smaller in this case than when using the pure LL eigenfunctions. With the advance of new approaches to the solution of the BFKL equation at full NLL accuracy$`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$ it has also been possible to check explicitly how well the two approximations compare with the full solution. We find that the approximation using the rescaled LL eigenfunctions is much closer to the true solution than the one using the pure LL ones. It should be noted that a saddle point approximation based around $`\nu =0`$ would have to use an extremely large order approximation in order to describe correctly the asymptotic intercept. Even the 16th order Taylor polynomial would fail to reach the maximum value for $`\omega ^{\mathrm{NLL}}`$ in Fig. 2. It would be far better to base a saddle point approximation around $`\nu _0`$.
Using the guess obtained from these rescaled eigenfunctions it is possible to calculate the intercept as the logarithmic derivative of the gluon Green’s function for fixed $`k_a`$ and $`k_b`$, as a function of the rapidity. This is depicted on Fig. 3. We see that although the NLL correction amounts to roughly a factor of two, it is stable. Also, it should be remembered that the study of both the LL and NLL intercept here has been performed without constraining the phase space of the BFKL resummation to such which is attainable at a given collider. The effects of such a constrain are known to be large$`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$ and reduce the LL evolution significantly (and presumably the NLL intercept to a slightly lesser extent).
## Conclusions
We have shown that the NLL corrections to the BFKL intercept are large but stable within the Regge kinematics. The instability with respect to the evolution in rapidity observed in initial analyses is a direct result of using the conformal set of leading log eigenfunctions as if they were eigenfunctions at NLL. Indeed, for $`\beta _0=0`$ this instability disappear even in this analysis, and this the case also for non-zero $`\beta _0`$ if one applies rescaled eigenfunctions. Although these still do not diagonalise the kernel at NLL, the results obtained using the rescaled eigenfunctions describe the full solution obtained numerically much better than the approximation obtain using just the LL eigenfunctions. In the conformal limit of $`\beta _0=0`$, where the NLL corrections can be studied exactly using the projections on the LL eigenfunctions, any modification of the NLL kernel to match better to the DGLAP region (like the ones of Ref.$`^{\mathrm{?},\mathrm{?}}`$) must move the position of the maximum of $`\omega ^{\mathrm{NLL}}`$ to $`\nu =0`$ while not changing the maximum value itself significantly (since this would lead to a change in the asymptotic intercept obtained within the Regge kinematic).
## Acknowledgments
I would like to thank Agustín Sabio Vera for lively discussions. This research was supported by PPARC (postdoctoral fellowship PPA/P/S/2003/00281).
## References |
warning/0507/gr-qc0507012.html | ar5iv | text | # Cylindrical solutions in braneworld gravity
## I Introduction
Since Einstein’s formulation of general theory of relativity , extra dimensions have always played a profound role in studies trying to unify gravity with other forces of nature. The very well known example is the Kaluza-Klein 5-dimensional theory whose main goal was the unification of gravitational and electromagnetic interactions via compactified fifth dimension. In recent years string inspired non-compactified extra dimensions have raised a lot of enthusiasm in the so called brane world gravity scenarios. In these scenarios, extra dimensions being large need not be compactified, instead we expect imprints of embedding and free gravitational field (in the 5-dimensional bulk spacetime of scenarios with one extra dimension) to show up in modified dynamical equations on the brane which are different from the usual four dimensional Einstein equations. Consequently, when we confine ourselves to the 3-brane, exact solutions to these modified field equations are expected to be either different from those in the usual four-dimensional general relativity or, if they are formally the same, to be reinterpreted in terms of some of the characteristics of the extra dimension. In 5-dimensional case, exact radial solutions with one and two parameters, are studied in where it is shown that there are different solutions according to the different equations of state for dark radiation and dark pressure. Gravitational collapse of matter, producing black holes on a brane, is studied in . There it is shown that a five dimensional black string solution intersects the brane in a Schwarzschild black hole. In the same scenario it is also shown that the two parameter Reissner-Nordstrom type solution could be interpreted in terms of mass and another parameter attributed to a characteristic of the bulk space . In this paper, using the modified Einstein equations on the brane, we study exact solutions on the brane with cylindrical symmetry. Actually what we will show is another manifestation of the fact that any solution of the Einstein-Maxwell equations (with traceless energy-momentum tensor) in 4-dimensional general relativity could be interpreted as a vacuum solution of brane-world 5-dimensional gravity .
## II Einstein field equations on the brane
A covariant generalization of Randall-Sundrum model is given in (the so called SMS braneworld <sup>*</sup><sup>*</sup>*Note that there are reformulations of the SMS braneworld in the literature with respect to different aspects of that formulation . ) where the 4-dimensional world is described by a 3-brane $`(M,q_{\mu \nu })`$ as a fixed point of $`Z_2`$ symmetry in a 5-dimensional space-time $`(,g_{\mu \nu })`$. Choosing the 5-th coordinate $`\chi `$ such that the brane coincides with the hypersurface $`\chi =0`$, the 5-dimensional metric could be written as;
$$ds^2=(n_\mu n_\nu +q_{\mu \nu })dx^\mu dx^\nu $$
(1)
where $`n_\mu `$ is the unit vector normal to the brane i.e $`n_\mu dx^\mu =d\chi `$ and $`q_{\mu \nu }`$ is the induced metric on the brane. Assuming that the matter is confined to the brane, the 5-dimensional Einstein equations read;
$`{}_{}{}^{(5)}G_{\mu \nu }^{}`$ $`=`$ $`\kappa _{(5)}^2{}_{}{}^{(5)}T_{\mu \nu }^{}`$ (2)
$`{}_{}{}^{(5)}T_{\mu \nu }^{}`$ $`=`$ $`\mathrm{\Lambda }_{(5)}g_{\mu \nu }+\delta (\chi )[\lambda _bq_{\mu \nu }+{}_{}{}^{(4)}T_{\mu \nu }^{}]`$ (3)
where $`\mathrm{\Lambda }_{(5)}`$ is the bulk cosmological constant and $`\lambda _b`$ is the vacuum energy on the brane or the brane tension in the bulk space. Since we assumed that the matter is confined to the brane, the 4-dimensional matter field Lagrangian determines the $`{}_{}{}^{(4)}T_{\mu \nu }^{}`$ such that $`{}_{}{}^{(4)}T_{\mu \nu }^{}n^\mu =0`$. Using the Gauss-Codacci equations, along with the $`Z_2`$ symmetry, one could project 5-dimensional curvature equations along the 4-dimensional hypersurface coincident with the brane. The effective Einstein field equations on the brane are found to be ;
$${}_{}{}^{(4)}G_{\mu \nu }^{}=\mathrm{\Lambda }q_{\mu \nu }+\kappa _{(4)}^2{}_{}{}^{(4)}T_{\mu \nu }^{}+\kappa _{(5)}^4S_{\mu \nu }_{\mu \nu }$$
(4)
where
$`S_{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{1}{12}}{}_{}{}^{(4)}T{}_{}{}^{(4)}T_{\mu \nu }^{}{\displaystyle \frac{1}{4}}{}_{}{}^{(4)}T_{\mu }^{\alpha }{}_{}{}^{(4)}T_{\alpha \nu }^{}+{\displaystyle \frac{1}{24}}q_{\mu \nu }(3{}_{}{}^{(4)}T_{}^{\alpha \beta }{}_{}{}^{(4)}T_{\alpha \beta }^{}{}_{}{}^{(4)}T_{}^{2})`$ (5)
$`\kappa _{(4)}^2`$ $`=`$ $`8\pi G_N`$ (6)
$`G_N`$ $`=`$ $`{\displaystyle \frac{\kappa _{(5)}^4\lambda _b}{48\pi }}`$ (7)
$`\mathrm{\Lambda }`$ $`=`$ $`{\displaystyle \frac{1}{2}}\kappa _{(5)}^2(\mathrm{\Lambda }_{(5)}+{\displaystyle \frac{1}{6}}\kappa _{(5)}^2\lambda _b^2)`$ (8)
and
$$_{\mu \nu }={}_{}{}^{(5)}C_{\beta \gamma \delta }^{\alpha }n_\alpha n^\gamma q_\mu ^\beta q_\nu ^\delta $$
(9)
is the transmitted projection of the bulk Weyl tensor To remind the reader, the n-dimensional ($`n3`$) Weyl tensor is defined by, $`{}_{}{}^{(n)}C_{\alpha \beta \gamma \delta }^{}={}_{}{}^{(n)}R_{\alpha \beta \gamma \delta }^{}+{\displaystyle \frac{1}{n2}}(g_{\alpha \delta }{}_{}{}^{(n)}R_{\beta \gamma }^{}+g_{\beta \gamma }{}_{}{}^{(n)}R_{\delta \alpha }^{}g_{\alpha \gamma }{}_{}{}^{(n)}R_{\delta \beta }^{}g_{\beta \delta }{}_{}{}^{(n)}R_{\gamma \alpha }^{})`$ (10) $`+{\displaystyle \frac{1}{(n1)(n2)}}(g_{\alpha \gamma }g_{\beta \delta }g_{\alpha \delta }g_{\gamma \beta }){}_{}{}^{(n)}R.`$ (11) . Note that it is symmetric and traceless. Equation (4) is the modified Einstein field equation on the brane due to the bulk effects. Compared with the usual 4-dimensional Einstein field equations, there are two main modifications: the first one is the presence of $`S_{\mu \nu }`$ term which is quadratic in $`{}_{}{}^{(4)}T_{\mu \nu }^{}`$. Comparison of the second and third terms in the right hand side of (4) shows that it is important in high energy limit and dominates when the energy density $`\rho \lambda _b`$ and is negligible for $`\rho \lambda _b`$. The second correction is due to $`_{\mu \nu }`$, the projection of the bulk Weyl tensor on the brane. From the viewpoint of a brane-observer the former is local and the latter is non-local. One recovers the usual Einstein field equations, by taking the limit $`\kappa _{(5)}0`$ while keeping $`G`$ finite. Though there is a point to be made here, according to (6) it is impossible to define Newton’s gravitational constant $`G_N`$ without an unambiguous definition of $`\lambda _b`$. Using equation (4), the 4D contracted Bianchi identity $`𝒟^\nu {}_{}{}^{(4)}G_{\mu \nu }^{}=0`$ along with the conservation of $`{}_{}{}^{(4)}T_{\mu \nu }^{}`$, $`𝒟^\nu {}_{}{}^{(4)}T_{\mu \nu }^{}=0`$ lead to the following constraint;
$$𝒟^\mu _{\mu \nu }=\kappa _{(5)}^4𝒟^\mu S_{\mu \nu },$$
(12)
Where $`𝒟^\mu `$ is the covariant derivative with respect to $`q_{\mu \nu }`$. Being symmetric and traceless, $`_{\mu \nu }`$ could be decomposed irreducibly with respect to a chosen 4- velocity vector field $`u^\mu `$ as ;
$$_{\mu \nu }=\kappa ^4[\text{U}(u_\mu u_\nu +\frac{1}{3}h_{\mu \nu })+\text{P}_{\mu \nu }+2\text{Q}_{(\mu }u_{\nu )}]$$
(13)
where $`h_{\mu \nu }=q_{\mu \nu }+u_\mu u_\nu `$ is the projection tensor orthogonal to $`u^\mu `$, $`\kappa =\kappa _{(5)}/\kappa _{(4)}`$ and
U $`=`$ $`\kappa ^4_{\mu \nu }u^\mu u^\nu `$ (14)
$`\text{Q}_\mu `$ $`=`$ $`\kappa ^4h_\mu ^\alpha _{\alpha \beta }u^\beta `$ (15)
$`\text{P}_{\mu \nu }`$ $`=`$ $`k^4[h_{(\mu }^\alpha h_{\nu )}^\beta {\displaystyle \frac{1}{3}}h_{\mu \nu }h^{\alpha \beta }]_{\alpha \beta }`$ (16)
U is an effective nonlocal energy density on the brane, arising from the bulk free gravitational field. Note that this nonlocal energy density need not be positive (in fact U contributes to tidal acceleration on the brane in the off-brane direction ) and actually being negative is more consistent with the Newotonian picture in which gravitational field carries negative energy . Also $`\text{P}_{\mu \nu }`$ and $`\text{Q}_\mu `$ are respectively the effective nonlocal anisotropic stress and energy flux on the brane, both arising from the free gravitational field in the bulk.
## III Static Cylindrically Symmetric Solution
In this section we solve the 4-dimensional modified Einstein equation (4) in vacuum ($`T_{\mu \nu }=0`$) for a static, cylindrically symmetric case. In the vacuum $`T_{\mu \nu }`$ and consequently $`S_{\mu \nu }`$ vanish and the constraint (12) becomes
$$𝒟^\mu _{\mu \nu }=0$$
(17)
One can Choose $`\mathrm{\Lambda }=0`$ by taking $`\mathrm{\Lambda }_{(5)}=1/6\kappa _{(5)}^2\lambda _b^2`$ so that the effective Einstein equations on the brane reduce to
$${}_{}{}^{(4)}R_{\mu \nu }^{}=_{\mu \nu }$$
(18)
As pointed out in the set of equations (17-18) form a closed system of equations for static solutions on the brane. In particular Einstein-Maxwell solutions (with traceless energy-momentum tensor) in 4-dimensional general relativity could be properly interpreted as vacuum brane world solutions.
Representing nonlocal effects, one could specify $`_{\mu \nu }`$ in an inertial frame at any point on the brane where $`u^\mu =\delta _0^\mu `$. In static spacetimes this corresponds to comoving frames along the timelike Killing vector field for which
$$\text{Q}_\mu =0,h_\mu ^\nu =\mathrm{diag}(0,1,1,1)$$
(19)
and the constraint for $`_{\mu \nu }`$ takes the form
$$\frac{1}{3}\stackrel{~}{}_\mu \text{U}+\frac{4}{3}\text{U}A_\mu +\stackrel{~}{}^\nu \text{P}_{\mu \nu }+A^\nu \text{P}_{\mu \nu }=0$$
(20)
where $`\stackrel{~}{}_\mu =h_\mu ^\nu 𝒟_\nu `$ is the projected covariant derivative orthogonal to $`u^\mu `$ and $`A_\mu =u^\nu 𝒟_\nu u_\mu `$ is the 4-acceleration. In the static cylindrically symmetric case we may choose
$$A_\mu =A(\rho )\rho _\mu $$
(21)
and
$$\text{P}_{\mu \nu }=P(\rho )(\rho _\mu \rho _\nu \frac{1}{3}h_{\mu \nu }).$$
(22)
where $`A(\rho )`$ and $`P(\rho )`$ are scalar functions and $`\rho _\mu `$ is the radial unit vector. Now in the (inertial) comoving frame the projected bulk Weyl tensor takes the form;
$$_\nu ^\mu =\kappa ^4\mathrm{diag}(U,\frac{1}{3}(U+2P),\frac{1}{3}(UP),\frac{1}{3}(UP))$$
(23)
Choosing the following general form for a static cylindrically symmetric line element on the brane,
$$ds^2=e^{2f(\rho )}dt^2+e^{2f(\rho )}e^{2K(\rho )}(d\rho ^2+dz^2)+e^{2f(\rho )}W(\rho )^2d\varphi ^2$$
(24)
substituting it into the modified gravitational field equations (18) and using (23) we end up with,
$`{\displaystyle \frac{1}{W}}(2f^{\prime \prime }W2f^{}W^{}+f^2W+K^{\prime \prime }W+W^{\prime \prime })`$ $`=`$ $`\kappa ^4\text{U}e^{2f}e^{2K}`$ (25)
$`{\displaystyle \frac{1}{W}}(f^2WK^{}W^{})`$ $`=`$ $`{\displaystyle \frac{\kappa ^4}{3}}(\text{U}+2\text{P})e^{2f}e^{2K}`$ (26)
$`{\displaystyle \frac{1}{W}}(f^2WK^{}W^{}+W^{\prime \prime })`$ $`=`$ $`{\displaystyle \frac{\kappa ^4}{3}}(\text{U}\text{P})e^{2f}e^{2K}`$ (27)
$`(f^2+K^{\prime \prime })`$ $`=`$ $`{\displaystyle \frac{\kappa ^4}{3}}(\text{U}\text{P})e^{2f}e^{2K}`$ (28)
The equation (17) could be written explicitly as
$$\text{U}^{^{}}(4f^{^{}}2K^{^{}}\frac{2W^{^{}}}{W})\text{U}=0$$
(29)
From (26) and (27), we obtain
$$\frac{W^{\prime \prime }}{W}=\frac{\kappa ^4}{3}(2\text{U}+\text{P})e^{2f}e^{2K}.$$
(30)
To solve the above equations we choose the following equation of state
$$2\text{U}+\text{P}=0$$
(31)
There is no physical restriction that implies this choice of equation of state, but it is not an easy task to solve the above equations in their general form, i.e for arbitrary $`U(\rho )`$ and $`P(\rho )`$, even in the spherical case . On the other hand it is interesting to note that the above equation of state in terms of the components of the projected Weyl tensor reads;
$$_0^0+_3^3=0$$
(32)
This establishes, once more and now for cylindrical symmetry, the connection stated above between solutions of the Einstein-Maxwell equations in general relativity and those of the vacuum brane world. For, (32) is the equation satisfied by the energy momentum tensor of Maxwell fields leading to known cylindrically symmetric solutions to the Einstein-Maxwell equations . We will discuss this in more detail in the next section.
By the above choice of the equation of state and from (30) we find $`W^{\prime \prime }=0`$ or $`W=a\rho +b`$, which by appropriate scaling of coordinates means either $`W=\rho `$ or $`W=const`$. Choosing $`W=\rho `$, (27) and (28) lead to the following solutions for $`e^{2K}`$,
$$e^{2K}=\rho ^{2m^2}$$
(33)
or
$$e^{2K}=1$$
(34)
Using the first solution and solving equations (25-29) we end up with
$`e^{2f}={\displaystyle \frac{4m^2}{(c_1\rho ^m+c_2\rho ^m)^2}}`$ (35)
$`U(\rho )={\displaystyle \frac{16m^2c_1c_2}{\kappa ^4}}{\displaystyle \frac{\rho ^{22m^2}}{(c_1\rho ^m+c_2\rho ^m)^4}}`$ (36)
where $`m`$, $`c_1`$ and $`c_2`$ are constants and $`c_1c_2<0`$. This last condition on the constants $`c_1`$ and $`c_2`$ could be infered from the comparison of this solution with the corresponding Einstein-Maxwell solution (refer to the next section). For the latter choice of $`e^{2K}`$ we arrive at
$`e^{2f}={\displaystyle \frac{1}{(c_1\mathrm{ln}\rho +c_2)^2}}`$ (37)
$`U(\rho )={\displaystyle \frac{c_1^2}{\kappa ^4\rho ^2(c_1\mathrm{ln}\rho +c_2)^4}}`$ (38)
Now setting $`W=1`$ and solving Einstein field equations along with the constraint equation (29) we find
$`K`$ $`=`$ $`a\rho +b`$ (39)
$`f`$ $`=`$ $`\mathrm{ln}(\rho )`$ (40)
U $`=`$ $`\kappa ^4{\displaystyle \frac{e^{2(a\rho +b)}}{\rho ^4}}`$ (41)
## IV Relation to Einstein-Maxwell solutions
As pointed out previously, the brane world solutions found in the last section could be treated as the solutions to the Einstein-Maxwell equations in general relativity. So one could use this analogy to find expressions for the corresponding electromagnetic fields in terms of the dark energy $`U`$ (or dark pressure $`P`$). We note that the solutions (26-32) are static cylindrically symmetric electrovacs discussed in the exact solutions literature (see section 22.2 of and references therin). Assuming that the electromagnetic fields inherit the metric symmetry and are orthogonal to the orbits of the two dimentional orthogonally transitive group, the vector potential could be written as
$$A_\alpha dx^\alpha =𝒫(\rho )dt+𝒬(\rho )d\varphi .$$
(42)
This vector potential corresponds to a magnetic field along the z-direction and an electric field along the $`\rho `$-direction. The non-zero components of $`F_{\mu \nu }`$ are $`F_{01}=𝒫^{}`$ and $`F_{13}=𝒬^{}`$ so that the corresponding energy-momentum tensor components $`T_\mu ^\nu `$ are given by;
$`T_0^0`$ $`=`$ $`{\displaystyle \frac{1}{2}}e^{2K}𝒫^{}_{}{}^{}2{\displaystyle \frac{1}{2}}e^{4f}e^{2K}W^2𝒬^{}_{}{}^{}2`$ (43)
$`T_1^1`$ $`=`$ $`{\displaystyle \frac{1}{2}}e^{2K}𝒫^{}_{}{}^{}2+{\displaystyle \frac{1}{2}}e^{4f}e^{2K}W^2𝒬^{}_{}{}^{}2`$ (44)
$`T_2^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}e^{2K}𝒫^{}_{}{}^{}2{\displaystyle \frac{1}{2}}e^{4f}e^{2K}W^2𝒬^{}_{}{}^{}2`$ (45)
$`T_3^3`$ $`=`$ $`{\displaystyle \frac{1}{2}}e^{2K}𝒫^{}_{}{}^{}2+{\displaystyle \frac{1}{2}}e^{4f}e^{2K}W^2𝒬^{}_{}{}^{}2`$ (46)
It is clear from the above equations that
$$T_0^0+T_3^3=0$$
(47)
Drawing analogies by comparing this equation with (32) we equate the right hand sides of equations (23) and (43-46) to arrive at
$$𝒬=\mathrm{constant}$$
(48)
obviously the equation of state (31) is also satisfied. Therefore dark energy and dark pressure in terms of the vector potential are given by
$$\text{U}=\frac{P}{2}=\frac{1}{2}\kappa ^4e^{2K}𝒫^{}_{}{}^{}2$$
(49)
or converesely the solutions (26-32) are soultions of the Einstein-Maxwell equations with the electromagnetic field given by the vector potential (33) with
$$𝒫=\kappa ^2e^K\sqrt{P(\rho )}d\rho $$
(50)
and $`𝒬`$ a constant which could be taken to be zero. It is noted from (49) that the dark energy, as expected from equations (29), (31) and (32), is a negative quantity. It is shown, by calculating the tidal acceleration on the brane in the off-brane direction, that negative dark energy enhances the localization of the gravitational field near the brane .
## V Discussion and summary
we have investigated exact static cylindrically symmetric solutions of the modified Einstein field equations for the induced metric on the brane. It is shown that, as in the case of the Reissner-Nordestrom type solutions , for the special choice of the equation of state $`2U+P=0`$ for the dark energy and dark pressure, the solutions found are that of the Einstein-Maxwell eqautions in the usual 4-dimensional general relativity. Of course there are no electromagnetic fields present, instead, using the analogy with the solutions to the Einstein-Maxwell solutions, the fields present could be interpreted as tidal fields (50) arising from the imprints, on the brane, of the bulk free gravitational field. The dark energy corresponding to the bulk free gravitational field is shown to be negative in all the solutions obtained, indicating that it acts in favour of confining the gravitational field near the brane. Finding an exact bulk solution that reduces to the exact induced metric on the brane is not an easy task and simple embeddings of the 4-d solution into the 5-d equations either results in equations which could not be solved or change the nature of the 4-d solution. As an example it could be seen that the solution given by equations (28) and (29) when embedded into the Randall-Sundrum general form
$$ds_{bulk}^2=f(y)(dy^2+ds_{brane}^2)$$
satisfies the bulk equations only when either $`c_1=0`$ or $`c_2=0`$. But in either case the 4-d solution reduces to the Levi-Civita’s Ricci flat solution . This is already known to satisfy the bulk equations since the Minkowski metric in Randall-Sundrum AdS solution could be replaced by any Ricci flat solution .
## Acknowledgement
M. N-Z thanks University of Tehran for supporting this project under the grants provided by the research council. |
warning/0507/hep-lat0507033.html | ar5iv | text | # One-loop matching of Δ𝑆=2 four-quark operators with improved staggered fermions
## I Introduction
Lattice calculations have started extracting phenomenologically relevant results with high precision and controlled uncertainties Davies:2003ik . This requires simulations with dynamical quarks, lattice actions with small discretization errors and good control over perturbative corrections. Effective field theory methods are an important tool to achieve these goals. They enable a judicious design of lattice actions with small discretization errors. On the other hand, the complicated form of improved lattice actions puts a burden on perturbative calculations required to match the output of lattice simulations to continuum physics. For the matching to be meaningful, perturbative calculations have to be performed with exactly the same action and operators, with which the non-perturbative simulations are done.
The purpose of the present paper is to compute the matching coefficients of the four-fermion $`\mathrm{\Delta }S=2`$ operators relevant for $`K\overline{K}`$ mixing for the Asqtad lattice action, an improved staggered fermion action used in many of the current simulations with dynamical quarks. The strength of $`K\overline{K}`$ mixing is determined by the matrix element $`\overline{K}|(\overline{s}d)_{VA}(\overline{s}d)_{VA}|KB_K(\mu )`$, where $`\mu `$ is the renormalization scale. The determination of $`B_K(\mu )`$ has become an important goal of lattice simulations, since a precise value of this parameter translates into a stringent constraint on the unitarity triangle. The current constraint is dominated by theoretical uncertainties, the largest being the uncertainty in the value of $`B_K`$ Charles:2004jd . The experimental input comes from $`K\pi \pi `$ decays and has percent level precision Eidelman:2004wy .
Results of lattice calculations for $`B_K(2\mathrm{GeV})`$ have been summarized in Hashimoto:2004hn ; Wingate:2004xa ; a central value of $`B_K(2\mathrm{GeV})=0.58(4)`$ from quenched determinations was advocated in Hashimoto:2004hn . The uncertainty associated with this value does not include the systematic error induced by quenching, i.e. neglecting the effects of dynamical fermions. The errors from quenching are presumed to be of the order of 15%, but are very hard to estimate reliably Sharpe:1998hh . In order to go beyond this accuracy, the calculation needs to be performed with dynamical fermions. The virtue of the improved action for which we perform our calculation is that it allows for precise calculations with light, dynamical quarks. The action is used by the MILC collaboration Bernard:2001av and consists of the Asqtad discretization of staggered quarks Lepage:1998vj ; Orginos:1999cr coupled to a Symanzik improved gluon action Luscher:1985zq .
The nonperturbative calculation of the matrix elements of the lattice operators on the MILC dynamical configurations was discussed in Gamiz:2004qx , along with preliminary results for the bare $`B_K`$ parameter. Together with the matching coefficients given in this paper, these results determine the renormalized continuum value $`B_K(\mu )`$ CDGSW05 .
In Becher:2002if two of us suggested an approach to lattice perturbation theory that allows an efficient separation of energy scales that appear in matching calculations – the inverse lattice spacing, which plays the role of the ultraviolet cut-off, and the physical masses and momenta. As a result of this scale separation, we are able to treat the complicated integrals that appear in lattice perturbation theory by algebraic means. The approach allows a very high degree of automation, thereby circumventing the complexity of improved lattice actions to a certain degree. We have used this method in Becher:2003fu to compute the one-loop relation between the bare lattice mass and the pole mass for the Asqtad action. In the present paper we compute the matching coefficient for $`\mathrm{\Delta }S=2`$ four-fermion operators for the Asqtad action, using the algebraic technique Becher:2002if as well as numerically. The agreement of the results obtained using these two completely different techniques provides a strong check on the calculation.
The lattice-to-continuum matching coefficients for the $`\mathrm{\Delta }S=2`$ four-quark operators have been previously calculated for the standard, unimproved, discretization of staggered fermions Ishizuka:1993fs , as well as for a class of improved staggered actions Lee:2003sk . While none of these results is directly relevant for the Asqtad action, we have checked that we reproduce the known result for the unimproved action, once we switch off all improvement terms.
The rest of the paper is organized as follows. In the next Section we give the discretization of the four-fermion operators. In Section III we discuss how the matching is performed and give the result for each of the relevant Feynman diagrams. In Section IV we include the tadpole improvement terms. We present our results for the matching coefficients that relate bare lattice operators to continuum $`\overline{\mathrm{MS}}`$ operators in Section V. Finally, in the last Section we summarize our results and compare the size of the perturbative corrections found for different staggered actions.
## II Staggered and naive lattice operators
For energy scales below the charm quark mass, the continuum effective Hamiltonian relevant for $`K\overline{K}`$ mixing reads Buras:1990fn
$$H_{\mathrm{eff}}^{\mathrm{weak}}=\frac{G_F^2m_W^2}{16\pi ^2}C(\mu )Q(\mu ),$$
(1)
where
$$Q=(\overline{s}d)_{VA}(\overline{s}d)_{VA}.$$
(2)
The possibility to describe $`K\overline{K}`$ mixing by a single operator is a consequence of the $`VA`$ structure of weak interactions and the Fierz identities that facilitate expressing any relevant four-quark operator in canonical form (2).
Lattice simulations of the $`K\overline{K}`$ mixing require introducing the operator $`Q`$ on the lattice. It is convenient to perform the perturbative calculations with naive instead of staggered quarks and to carry out the spin-diagonalization and the reduction from the sixteen to four doublers (or “tastes”) afterwards. We use the standard discretization for the staggered four-quark operators and now show how these operators can be rewritten in terms of the naive fermion field.
The naive fermion field $`\psi `$ and the staggered field $`\chi `$ at lattice site $`n=(n_1,n_2,n_3,n_4)`$ are related by
$`\psi (n)`$ $`=\mathrm{\Gamma }_n\chi (n)`$ $`\overline{\psi }(n)`$ $`=\overline{\chi }(n)\mathrm{\Gamma }_n^{}`$ (3)
where
$$\mathrm{\Gamma }_n=(\gamma _1)^{n_1}(\gamma _2)^{n_2}(\gamma _3)^{n_3}(\gamma _4)^{n_4}.$$
(4)
We choose a Hermitian representation of the Euclidean Dirac algebra with
$`\{\gamma _i,\gamma _j\}`$ $`=2\delta _{ij}`$ and $`\gamma _5`$ $`=\gamma _1\gamma _2\gamma _3\gamma _4.`$ (5)
When written in terms of staggered fermions, the quark action becomes spin diagonal and three of the four components of the quark field $`\chi `$ can be dropped. In perturbation theory it is more convenient to keep all components, and replace $`n_fn_f/4`$ at the end of the calculations to obtain the result for the single component field.
The fields $`\chi (2N+A)`$ are collected into a set of Dirac fields $`q(2N)`$ that live on the even lattice sites and are spread over a unit hypercube Gliozzi:1982ib ; Kluberg-Stern:1983dg ; Sharpe:1993ur ($`A_\mu =0`$ or 1, $`\mu =1\mathrm{}4`$)
$$q(2N)_{\alpha ij}=\frac{1}{8}\underset{A}{}(\mathrm{\Gamma }_A)_{\alpha ,i}\chi _j(2N+A).$$
(6)
The index $`\alpha `$ is the Dirac index of the new field and $`i`$ the “taste” index. The second taste index $`j`$ runs over the components of the field $`\chi `$. Staggered fermions have only one component, while the four component of the naive field give rise to $`4\times 4`$ tastes. Bilinear quark operators with spin structure $`\gamma _S`$ and taste structure $`\xi _T=\mathrm{\Gamma }_T^{}`$ (with $`T`$ again on the unit hypercube) are
$$\begin{array}{c}\overline{q}(2N)(\gamma _S\xi _T)q(2N)\hfill \\ \hfill =\frac{1}{16}\underset{A,B}{}\overline{\chi }(2N+A)\chi (2N+B)\frac{1}{4}\mathrm{tr}(\mathrm{\Gamma }_A^{}\gamma _S\mathrm{\Gamma }_B\mathrm{\Gamma }_T^{}).\end{array}$$
(7)
Since the quark fields in the operators are at different lattice points, they need to be connected by Wilson lines in order to make the operators gauge invariant. We suppress these Wilson lines, but it is understood that gauge strings among all possible shortest paths connecting the quark fields are inserted and the operator is divided by the number of paths. Using (3) we rewrite the operator in terms of the naive fermion field:
$$\begin{array}{c}\overline{q}(2N)(\gamma _S\xi _T)q(2N)\hfill \\ \hfill =\frac{1}{16}\underset{A,B}{}\overline{\psi }(2N+A)\mathrm{\Gamma }_A\mathrm{\Gamma }_B^{}\psi (2N+B)\frac{1}{4}\mathrm{tr}(\mathrm{\Gamma }_A^{}\gamma _S\mathrm{\Gamma }_B\gamma _T^{}).\end{array}$$
(8)
Let us write out the vector and axial currents at $`N=0`$ with unit taste structure in terms of the naive quark field $`\psi (n)`$. We find
$$\begin{array}{c}\overline{q}(\gamma _\mu 1)q=\frac{1}{16}\underset{A}{}\delta _{A_\mu ,0}\hfill \\ \hfill \left[\overline{\psi }(A)\gamma _\mu \psi (A+\widehat{\mu })+\overline{\psi }(A+\widehat{\mu })\gamma _\mu \psi (A)\right],\end{array}$$
(9)
and ($`\overline{A}=(1,1,1,1)A`$)
$$\begin{array}{c}\overline{q}(\gamma _\mu \gamma _51)q=\frac{1}{16}\underset{A}{}\delta _{A_\mu ,0}\hfill \\ \hfill \left[\overline{\psi }(A)\gamma _\mu \gamma _5\psi (\overline{A+\widehat{\mu }})+\overline{\psi }(\overline{A})\gamma _\mu \gamma _5\psi (A+\widehat{\mu })\right].\end{array}$$
(10)
In the continuum limit, the massless staggered action has a $`SU(4)_L\times SU(4)_R`$ chiral symmetry (for a single staggered field) and one would naively expect to find fifteen Goldstone bosons after chiral symmetry breaking, corresponding to the fifteen traceless taste matrices. However, this symmetry is explicitly broken to an axial $`U(1)`$ symmetry by terms in the action proportional to the lattice spacing squared. The generator of the remaining axial symmetry is the matrix $`\gamma _S\xi _T=\gamma _5\gamma _5`$. In our case, this implies that only the kaon with taste structure $`\gamma _5`$ becomes massless in the chiral limit. For this reason, the simulation is done with currents having $`\xi _T=\gamma _5`$ taste structure and not the unit matrix as in (9) and (10). However, it has been shown that the result for diagrams in which the fermions in the operators are not contracted are identical for the Dirac and taste structures $`\mathrm{\Gamma }\mathrm{\Gamma }^{}`$ and $`\mathrm{\Gamma }\gamma _5\mathrm{\Gamma }^{}\gamma _5`$ Sharpe:1993ur .<sup>1</sup><sup>1</sup>1The reason is that in this case, one can view the four fermions as four different quark flavors and perform separate axial $`U(1)`$ transformations on each field. For our perturbative calculation we will thus use (10) instead of the vector current $`\gamma _\mu \gamma _5`$ and (9) in place of $`\gamma _\mu \gamma _5\gamma _5`$.
The operator $`Q`$ for which we want to perform the matching calculation is a product of two $`VA`$ currents. We will see that QCD corrections to bare lattice four-quark operators affect the vector and axial parts differently; as a consequence currents of the form $`V+A`$ are generated. A minimal set of lattice operators that matches to the continuum and closes under renormalization consists of 4 scalar and 2 pseudoscalar operators. Schematically these operators are
$`\begin{array}{cccc}\hfill Q_1& =(\overline{s}^ad^a)_V(\overline{s}^bd^b)_V,\hfill & & Q_2=(\overline{s}^ad^b)_V(\overline{s}^bd^a)_V,\hfill \\ \hfill Q_3& =(\overline{s}^ad^a)_A(\overline{s}^bd^b)_A,\hfill & & Q_4=(\overline{s}^ad^b)_A(\overline{s}^bd^a)_A,\hfill \end{array}`$
$`\begin{array}{cc}& Q_5=(\overline{s}^ad^a)_V(\overline{s}^bd^b)_A+(\overline{s}^ad^a)_A(\overline{s}^bd^b)_V,\hfill \\ & Q_6=(\overline{s}^ad^b)_V(\overline{s}^bd^a)_A+(\overline{s}^ad^b)_A(\overline{s}^bd^a)_V.\hfill \end{array}`$ (11)
$`V`$ and $`A`$ are the vector and axial currents with taste structure $`\xi _T=\gamma _5`$. The color indices $`a,b`$ indicate which fields are connected by Wilson lines. The pseudoscalar operators $`Q_{5,6}`$ only mix amongst themselves and their $`K\overline{K}`$ matrix element vanishes, because of parity. In the following we will therefore only consider the renormalization of the operators $`Q_{14}`$.
The above operators are used in the nonperturbative calculation of the matrix elements by the HPQCD collaboration Gamiz:2004qx ; CDGSW05 . In contrast to the lattice action, these operators are not Symanzik improved.
## III Matching calculation
To perform the matching we calculate a physical quantity in the continuum and on the lattice, expand around the continuum limit, and adjust the Wilson coefficients of the lattice operators in such a way that they reproduce the continuum result. Specifically, we determine the bare lattice Wilson coefficients from the quark-quark scattering amplitude
$`𝒜=Z^2C_i^{\mathrm{bare}}O_i^{\mathrm{bare}}=Z^2C_iZ_{ij}O_j=Z_{\mathrm{lat}}^2C_i^{\mathrm{lat}}O_i^{\mathrm{lat}},`$ (12)
where $`Z`$ and $`Z_{\mathrm{lat}}`$ are on-shell quark wave-function renormalizations and $`O`$ is the amputated four-quark Green’s function with an insertion of the operator $`O`$. The lattice-to-continuum matching coefficients are independent of the quantity chosen to perform the matching. At one-loop order, it is simplest to regulate infrared divergencies with a gluon mass and to calculate the scattering amplitude of massless quarks at zero external momentum.
At tree level in the continuum limit, the operators $`Q_i`$ are identical to their continuum counterparts; this implies that no tree level matching is required. At one loop the matching coefficients $`C_i^{\mathrm{lat}}`$ are obtained by calculating the difference between the lattice and continuum one-loop diagrams. In the difference, the infrared divergencies cancel. Below we present such a calculation for the set of lattice operators described above.
For unimproved staggered fermions the renormalization of the axial current operator with $`\gamma _S\xi _T=\gamma _\mu \gamma _5\gamma _5`$ is finite and identical in the continuum and on the lattice. This happens because it renormalizes in the same way as the vector current with unit taste structure $`\gamma _\mu 1`$ which has exactly the same form as the quark-gluon coupling in the Lagrangian. Because of the improvement terms, this is not true for the Asqtad action and we will need to evaluate the matching not only for the four-quark operators but also for the current.
The matching calculation requires computing matrix elements of certain operators in lattice perturbation theory; to do so we use the approach described in Becher:2002if ; Becher:2003fu . In those references an expansion of the lattice integrals around the continuum limit is introduced by utilizing the technique of asymptotic expansions. The asymptotic expansion around the continuum limit splits the lattice diagrams into two parts: (i) a soft part which is obtained by evaluating continuum loop integrals (in analytic regularization) and is independent of the discretization and (ii) a hard part which depends on the discretization but is independent of particle masses and momenta. The proximity of the soft part of the lattice loop integrals and the continuum integrals in dimensional regularization permits significant simplifications when computing the difference of the lattice and continuum Green’s functions required for matching calculations.
There are five diagrams (plus permutations) that contribute at $`𝒪(\alpha _s)`$ to the quark-quark scattering amplitude; they are shown in Fig. 1. To present the result of the calculations, we introduce the following notation for the difference of the lattice and continuum one-loop matrix elements of the four-fermion operators:
$$\delta Q^i_\alpha =\left(\frac{\alpha _s}{4\pi }\right)\left(\frac{\mu a}{2}\right)^{2ϵ}\underset{j=1}{\overset{4}{}}q_{\alpha ,j}^iQ^j,$$
(13)
where the label $`\alpha `$ refers to a particular diagram shown in Fig. 1. The spacetime dimension is $`d=42ϵ`$, $`\mu `$ is the renormalization scale in the $`\overline{\mathrm{MS}}`$ subtraction scheme and $`a`$ the lattice spacing. In what follows, we split the matrices $`q_i`$ into soft and hard pieces. We give the result for each diagram separately. We hope that the results for the individual diagrams will be useful for readers performing similar calculations in the future.
### III.1 Soft part
As a first step in the matching calculation, we compute the difference between the soft parts of lattice matrix elements and the result derived in dimensional regularization. To obtain the full result, we later add the hard part for a given lattice action. Technically the soft part is obtained by evaluating the continuum diagrams shown in Fig. 1, first in dimensional and then in analytic regularization. The calculation is simplified by noting that a difference between the two regulators can only arise from ultraviolet divergent loop integrals.
Below we give the results for the difference of the soft parts of the matrices $`q`$ and the corresponding continuum diagrams. Since the diagrams 4 and 5 in Fig. 1 do not have continuum analog, their contribution to the soft part is zero: $`q_4^{\mathrm{soft}}=q_5^{\mathrm{soft}}=0`$. For diagram $`1`$, we derive
$$q_1^{\mathrm{soft}}=\left(\begin{array}{cccc}C_Ff_{11}^{(1)}& 0& 0& 0\\ \frac{f_{11}^{(1)}}{2}& \frac{f_{11}^{(1)}}{2N}& 0& 0\\ 0& 0& C_Ff_{11}^{(1)}& 0\\ 0& 0& \frac{f_{11}^{(1)}}{2}& \frac{f_{11}^{(1)}}{2N}\end{array}\right).$$
The variable $`N=3`$ denotes the number of colors, $`C_F=(N^21)/2N`$ and
$`f_{11}^{(1)}=2(1+\xi )\left({\displaystyle \frac{1}{ϵ}}{\displaystyle \frac{1}{\delta }}\right)+12\xi .`$ (14)
The parameter $`\delta `$ is an intermediate analytic regulator, which will drop out in the sum of the hard and soft part Becher:2002if . Setting the gauge parameter $`\xi =0`$, one obtains the result in Feynman gauge while $`\xi =1`$ corresponds to Landau gauge. For details on our notation and method of calculation, we refer the reader to Becher:2003fu .
For diagrams 2 and 3 in Fig. 1 we obtain,
$$q_2^{\mathrm{soft}}=\left(\begin{array}{cccc}\frac{f_{11}^{(2)}}{2N}& \frac{f_{11}^{(2)}}{2}& \frac{f_{13}^{(2)}}{2N}& \frac{f_{13}^{(2)}}{2}\\ \frac{f_{11}^{(2)}}{2}& \frac{f_{11}^{(2)}}{2N}& \frac{f_{13}^{(2)}}{2}& \frac{f_{13}^{(2)}}{2N}\\ \frac{f_{13}^{(2)}}{2N}& \frac{f_{13}^{(2)}}{2}& \frac{f_{11}^{(2)}}{2N}& \frac{f_{11}^{(2)}}{2}\\ \frac{f_{13}^{(2)}}{2}& \frac{f_{13}^{(2)}}{2N}& \frac{f_{11}^{(2)}}{2}& \frac{f_{11}^{(2)}}{2N}\end{array}\right),$$
$$q_3^{\mathrm{soft}}=\left(\begin{array}{cccc}\frac{f_{11}^{(3)}}{2N}& \frac{f_{11}^{(3)}}{2}& \frac{f_{13}^{(3)}}{2N}& \frac{f_{13}^{(3)}}{2}\\ 0& C_Ff_{11}^{(3)}& 0& C_Ff_{13}^{(3)}\\ \frac{f_{13}^{(3)}}{2N}& \frac{f_{13}^{(3)}}{2}& \frac{f_{11}^{(3)}}{2N}& \frac{f_{11}^{(3)}}{2}\\ 0& C_Ff_{13}^{(3)}& 0& C_Ff_{11}^{(3)}\end{array}\right),$$
where
$`f_{11}^{(2)}=f_{11}^{(3)}+6=\left(5+2\xi \right)\left({\displaystyle \frac{1}{ϵ}}{\displaystyle \frac{1}{\delta }}\right)+{\displaystyle \frac{15}{2}}+2\xi ,`$
$`f_{13}^{(2)}=f_{13}^{(3)}=3\left({\displaystyle \frac{1}{ϵ}}{\displaystyle \frac{1}{\delta }}\right)+{\displaystyle \frac{5}{2}}.`$ (15)
To evaluate the continuum diagrams, a prescription for the treatment of $`\gamma _5`$ and the extension of the Fierz identities to $`d`$-dimensions Dugan:1990df ; Buras:1989xd ; Herrlich:1994kh has to be adopted. We use naive dimensional regularization for $`\gamma _5`$ and treat the evanescent operators exactly as in Buras:1990fn , to be compatible with the results for the Wilson coefficients given in these references. Note that with this prescription, the operators $`Q_1\gamma _\mu \gamma _\mu `$ and $`Q_3\gamma _\mu \gamma _5\gamma _\mu \gamma _5`$ also mix into operators with Dirac structure $`11`$, $`\gamma _5\gamma _5`$ and $`\sigma _{\mu \nu }\sigma _{\mu \nu }`$. We do not give the result for this mixing, since it drops out in the sums $`Q_1+Q_3`$ and $`Q_2+Q_4`$ which are relevant for our matching calculation. However, separating the vector and axial contributions also in the soft part allows us to read off the current renormalizations from the first diagram in Fig. 1 and provides additional consistency checks, such as cancellation of $`\delta `$-poles and gauge invariance of the sum of hard and soft parts.
In addition to the diagrams shown in Fig. 1, we have to account for the external leg corrections, given by the wave function renormalization of the massless quarks. This correction is universal in that it does not depend on the operator under consideration. We obtain
$$\delta Z^{(1)}O_i|_{\mathrm{soft}}=C_F\left(\frac{2}{ϵ}\frac{2}{\delta }+1\right)(1+\xi )O_i,$$
(16)
where $`Z_{\mathrm{lat}}^2Z^2`$=$`\frac{\alpha _s}{\pi }\delta Z^{(1)}+\mathrm{}`$ is the difference between the lattice (soft part only) and the on-shell quark wave function renormalization constant. This difference vanishes in Landau gauge, which is a consequence of the fact that in this gauge the continuum quark wave function renormalization is finite.
Note that the coefficients of the $`1/ϵ`$ and $`1/\delta `$ pieces in all the equations above are exactly opposite. The $`1/\delta `$-divergences will cancel against those of the hard parts. The fact that the above differences do not depend on any infrared scales (the gluon mass in our case) shows that the result of a matching calculation is process independent.
### III.2 Hard part
We now give the result of our calculation of the hard part. This amounts to presenting the matrices $`q_\alpha `$ for the five diagrams in Fig. 1. The symmetry factors for each of the diagrams are included. We begin with naive unimproved fermions coupled to unimproved glue.
The result for the hard part of the first three diagrams in Fig. 1, has the same structure as was found for the soft part:
$`q_1^{\mathrm{hard}}`$ $`=\left(\begin{array}{cccc}C_Fd_{11}^{(1)}& 0& 0& 0\\ \frac{d_{11}^{(1)}}{2}& \frac{d_{11}^{(1)}}{2N}& 0& 0\\ 0& 0& C_Fd_{33}^{(1)}& 0\\ 0& 0& \frac{d_{33}^{(1)}}{2}& \frac{d_{33}^{(1)}}{2N}\end{array}\right),`$
where
$`d_{11}^{(1)}`$ $`={\displaystyle \frac{2(1+\xi )}{\delta }}3.532.38\xi ,`$
$`d_{33}^{(1)}`$ $`={\displaystyle \frac{2(1+\xi )}{\delta }}3.62+0.695\xi ,`$ (17)
and
$`q_2^{\mathrm{hard}}`$ $`=\left(\begin{array}{cccc}\frac{d_{11}^{(2)}}{2N}& \frac{d_{11}^{(2)}}{2}& \frac{d_{13}^{(2)}}{2N}& \frac{d_{13}^{(2)}}{2}\\ \frac{d_{11}^{(2)}}{2}& \frac{d_{11}^{(2)}}{2N}& \frac{d_{13}^{(2)}}{2}& \frac{d_{13}^{(2)}}{2N}\\ \frac{d_{13}^{(2)}}{2N}& \frac{d_{13}^{(2)}}{2}& \frac{d_{11}^{(2)}}{2N}& \frac{d_{11}^{(2)}}{2}\\ \frac{d_{13}^{(2)}}{2}& \frac{d_{13}^{(2)}}{2N}& \frac{d_{11}^{(2)}}{2}& \frac{d_{11}^{(2)}}{2N}\end{array}\right),`$
$`q_3^{\mathrm{hard}}`$ $`=\left(\begin{array}{cccc}\frac{d_{11}^{(3)}}{2N}& \frac{d_{11}^{(3)}}{2}& \frac{d_{13}^{(3)}}{2N}& \frac{d_{13}^{(3)}}{2}\\ 0& C_Fd_{11}^{(3)}& 0& C_Fd_{13}^{(3)}\\ \frac{d_{13}^{(3)}}{2N}& \frac{d_{13}^{(3)}}{2}& \frac{d_{11}^{(3)}}{2N}& \frac{d_{11}^{(3)}}{2}\\ 0& C_Fd_{13}^{(3)}& 0& C_Fd_{11}^{(3)}\end{array}\right),`$
with
$`d_{11}^{(2)}`$ $`=d_{11}^{(3)}={\displaystyle \frac{(5+2\xi )}{\delta }}+0.157+0.689\xi ,`$
$`d_{13}^{(2)}`$ $`=d_{13}^{(3)}={\displaystyle \frac{3}{\delta }}0.341,`$ (18)
The remaining two diagrams have no continuum analogue. We obtain
$`q_{i=4,5}^{\mathrm{hard}}`$ $`=\left(\begin{array}{cccc}C_Fd_{11}^{(i)}& 0& 0& 0\\ d_{21}^{(i)}& C_Fd_{22}^{(i)}\frac{d_{21}^{(i)}}{N}& 0& 0\\ 0& 0& C_Fd_{33}^{(i)}& 0\\ 0& 0& d_{43}^{(i)}& C_Fd_{44}^{(i)}\frac{d_{43}^{(i)}}{N}\end{array}\right),`$
with entries
$`d_{11}^{(4)}`$ $`=16.06+18.38\xi ,`$
$`d_{21}^{(4)}`$ $`=1.18+1.69\xi ,`$
$`d_{22}^{(4)}`$ $`=13.70+15.00\xi ,`$ (19)
$`d_{33}^{(4)}`$ $`=12.23+12.23\xi ,`$
$`d_{43}^{(4)}`$ $`=0.731.38\xi ,`$
$`d_{44}^{(4)}`$ $`=13.70+15.00\xi ,`$
$`d_{11}^{(5)}`$ $`=73.409.19\xi ,`$
$`d_{21}^{(5)}`$ $`=2.200.85\xi ,`$
$`d_{22}^{(5)}`$ $`=48.97.50\xi ,`$ (20)
$`d_{33}^{(5)}`$ $`=24.476.12\xi ,`$
$`d_{43}^{(5)}`$ $`=1.57+0.69\xi ,`$
$`d_{44}^{(5)}`$ $`=48.97.50\xi .`$
We also give the result for the hard part fermion wave function renormalization on the lattice:
$`\delta Z^{(1)}O_i|_{\mathrm{hard}}`$ $`=C_F(10.61411.928\xi +{\displaystyle \frac{2(1+\xi )}{\delta }}`$ (21)
$`+[24.466+6.117\xi ])O_i.`$
We have separated out the tadpole contribution in (21) in square brackets. In Landau gauge this contribution exactly cancels against the tadpole improvement term.
We now give the results for the entries of the matrices for the case of Asqtad action. In this case, as explained in Becher:2003fu , we expand in the improvement terms, for both quark and gluon actions. We use half of the highest order terms as the numerical error estimate and multiply the last term of the series by $`2/3`$ to resum the higher order contributions that behave like an alternating geometric series with expansion parameter $`1/2`$ (see Becher:2003fu for details). We find
$`d_{11}^{(1)}`$ $`={\displaystyle \frac{2(1+\xi )}{\delta }}3.00(1)2.379\xi ,`$
$`d_{33}^{(1)}`$ $`={\displaystyle \frac{2(1+\xi )}{\delta }}+0.81(4)+0.695\xi `$
$`d_{11}^{(2)}`$ $`=d_{11}^{(3)}={\displaystyle \frac{(5+2\xi )}{\delta }}+2.79(10)+0.689\xi ,`$
$`d_{13}^{(2)}`$ $`=d_{13}^{(3)}={\displaystyle \frac{3}{\delta }}+2.10(9),`$
$`d_{11}^{(4)}`$ $`=18.54(6)+18.38\xi ,`$
$`d_{21}^{(4)}`$ $`=1.73(1)+1.69\xi ,`$
$`d_{22}^{(4)}`$ $`=15.09(3)+15.00\xi ,`$ (22)
$`d_{33}^{(4)}`$ $`=12.23+12.23\xi ,`$
$`d_{43}^{(4)}`$ $`=1.43(2)1.38\xi ,`$
$`d_{44}^{(4)}`$ $`=15.09(3)+15.00\xi ,`$
$`d_{11}^{(5)}`$ $`=58.0(5)9.19\xi ,`$
$`d_{21}^{(5)}`$ $`=2.060.84\xi ,`$
$`d_{21}^{(5)}`$ $`=39.1(3)7.50\xi ,`$
$`d_{33}^{(5)}`$ $`=20.3(1)6.12\xi ,`$
$`d_{43}^{(5)}`$ $`=1.56(1)+0.69\xi ,`$
$`d_{44}^{(5)}`$ $`=39.1(3)7.50\xi .`$
The gauge dependent part is the same as in the unimproved case because the improvement terms in the action are transverse Becher:2003fu .
For the Asqtad action, the wave function renormalization contribution is:
$`\delta Z^{(1)}O_i|_{\mathrm{hard}}`$ $`=C_F(14.3(2)11.28\xi +{\displaystyle \frac{2(1+\xi )}{\delta }}`$ (23)
$`+[30.1(6)+5.47\xi ])O_i.`$
The renormalization of the axial and vector current can be read off from the above expressions. In both cases it is given by one half of the corresponding entries in the matrices $`q_i`$ for diagrams 1, 4 and 5 in Fig. 1 and the contribution due to the wave function renormalization. As discussed earlier in this section, the discretized version of the axial current that we use in this paper renormalizes identically to the unimproved quark-gluon vertex. As a consequence, the lattice-to-continuum matching coefficient for the axial current should be equal to unity for the unimproved action, to all orders in $`\alpha _s`$. Using the results presented above, it is easy to verify this at the one-loop level.
### III.3 Matching coefficients from a different method
Most of the diagrams needed in the determination of the matching coefficients for four-fermion operators can be calculated from the renormalization coefficients for bilinears operators, as was first pointed out by Martinelli Martinelli84 . The method has been generalized to Landau-gauge operators in Sharpe:1993ur , general local operators in GBS96 and staggered gauge invariant operators in Lee:2003sk .
In Fig. 1, we did not show mirrored diagrams and did not display the color structure of the four-quark operators. Also, we have not indicated from which of the two currents the gluons in diagrams 4 and 5 are emitted. Listing all possibilities, one ends up with 16 diagrams of type 4 and 6 of type 5. Only 8 of the 16 diagrams of type 4 and 2 of the 6 diagrams of type 5 are genuine four-fermion contributions. All other diagrams factor into a product of a one-loop current renormalization diagram times a tree level current. Furthermore, only half of the remaining diagrams contribute to the mixing between operators whose bilinears have identical taste-structure, which is all that is relevant for matrix elements involving external states of definite taste. One cannot avoid calculating the genuine four-fermion diagrams, but the result for the other topologies can be obtained from the diagonal renormalization coefficients for bilinears operators using Fierz and charge conjugation transformations as described in Lee:2003sk . Apart from the 5 genuine four-fermion diagrams, the calculation reduces to the evaluation of the current renormalization diagrams in Fig. 2 and the universal wave function renormalization of the massless quarks. In Fig. 2, the crossed circle denotes the insertion of any bilinear with structure $`(\gamma _S\xi _T)`$. The calculation is further simplified by noting that the corrections are identical if the operators are multiplied by $`(\gamma _5\gamma _5)`$ due to the axial symmetry, as stressed above.
For diagonal four-fermion operators (with identical taste and spin structure in both bilinears) this procedure gives us a $`35\times 35`$ matrix from which we can extract the $`2\times 2`$ block describing the mixing between $`[\gamma _\mu \gamma _5][\gamma _\mu \gamma _5]`$ and $`[\gamma _\mu \gamma _5\gamma _5][\gamma _\mu \gamma _5\gamma _5]`$. This block is related to the $`4\times 4`$ matrix in (32) by including the appropriate color factors as explained in Lee:2003sk .
We have done the calculation of the matching coefficients for the $`\mathrm{\Delta }S=2`$ four-quark operators using the method described above and evaluating the integrals numerically instead of algebraically. In all cases we find full agreement with the direct evaluation of the diagrams. This independent calculation of the matching coefficients provides a non-trivial check of the validity of our results. In addition, we have also performed the calculation for the HYP actions used in Lee:2003sk and reproduce the result for the matching matrix given in that reference.
## IV Tadpole improvement
Quite generally one finds that perturbative corrections in lattice regularization are larger than those encountered in continuum calculations, e.g. in the $`\overline{\mathrm{MS}}`$ scheme. The large corrections typically arise from the additional diagrams that would be absent in the continuum, in our case the fourth and fifth diagram in Fig. 1. Since the size of the corrections is scheme dependent they are not by themselves meaningful. Instead of just working with the bare lattice parameters one can try to absorb the large perturbative corrections into a redefinition of the basic parameters of the theory, such as coupling constant, quark masses, Wilson coefficients of the weak Hamiltonian, etc. Such techniques are also used in the continuum, for example in the context of Heavy-Quark Effective theories: the perturbative results for inclusive heavy hadron decay rates get large perturbative corrections if they are expressed in terms of the heavy quark pole mass but the corrections are much smaller if one uses a more suitable mass definition Bigi:1997fj ; Czarnecki:1997sz ; Hoang:1998ng . Similar ideas are behind the so-called “physical” couplings and optimal scale setting prescriptions, introduced quite some time ago Brodsky:1982gc .
The large effects of the tadpole diagrams can be canceled by tadpole improvement Lepage:1992xa . Tadpole improvement is achieved by working with an action in which all gauge links are divided by an average link $`u_0`$. In the simulation the average link is determined numerically, while it is evaluated perturbatively in the matching calculation. We define
$$u_0=1\alpha _su_0^{(1)}+\mathrm{}.$$
(24)
The average link receives large perturbative corrections which then cancel out the large contributions from the tadpole diagrams. In this section we give the tadpole improvement terms relevant for our calculation. Some care is required because different variants of tadpole improvement are used in the literature. First of all, two different definitions of the the mean link $`u_0`$ are common: the mean link in Landau gauge
$$u_0=\frac{1}{3}\mathrm{Re}\mathrm{Tr}U_{n,\mu }_{\xi =1}$$
(25)
or the fourth root of the average plaquette
$$u_0^P=\frac{1}{3}\mathrm{Re}\mathrm{Tr}U_{n,\mu }U_{n+\mu ,\nu }U_{n+\mu +\nu ,\mu }^{}U_{n+\nu ,\nu }^{}^{1/4}.$$
(26)
The second definition is somewhat simpler since no gauge fixing is required and it is what MILC is using in their simulations with the Asqtad action. Furthermore, MILC divides the $`L`$-link terms in the action and operators not by $`u_0^L`$ as Lepage:1992xa but by $`u_0^{L1}`$ and absorbs a power of $`u_0`$ into the fermion mass. With their prescription no tadpole improvement factor is present in the fermion part of the naive unimproved action Orginos:1999cr .<sup>2</sup><sup>2</sup>2We thank C. Bernard for pointing this out to us. Since the fermion mass does not enter our calculation, the difference between the two prescriptions merely amount to a rescaling of the fermion field and leads to identical results for the matching coefficients we calculate.
If we divide the $`L`$-link terms by $`u_0^L`$, the tadpole improvement contribution from the current operators is
$$q^{(\mathrm{tad})}=4\pi u_0^{(1)}\left(\begin{array}{cccc}6& 0& 0& 0\\ 0& 4& 0& 0\\ 0& 0& 2& 0\\ 0& 0& 0& 4\end{array}\right).$$
(27)
In addition, there is a tadpole improvement correction to the square of the wave function renormalization
$`\delta Z_{\mathrm{unimp}.}^{\mathrm{tad}}Q_i`$ $`=8\pi u_0^{(1)}Q_i,`$
$`\delta Z_{\mathrm{Asqtad}}^{\mathrm{tad}}Q_i`$ $`=18\pi u_0^{(1)}Q_i.`$ (28)
The results given in the next section are obtained with $`u_0`$ defined as the mean link in Landau gauge (25). However, using the explicit form of the tadpole improvement terms given above the prescription can easily be changed.
## V Wilson coefficients of the lattice operators
The Wilson coefficients of the lattice operators can now be read off by imposing that the four-quark scattering amplitudes are equal in the continuum and on the lattice. Clearly, the solution to the matching equation (12) is in general not unique since there are many different discretizations of the same continuum operator. In the present case, the two $`(\overline{s}d)_{VA}(\overline{s}d)_{VA}`$ operators with different color structure become equivalent in the continuum limit where they are related by a Fierz transformation. The matching condition will thus not fix the Wilson coefficients of the operators with different color structure individually. To be specific, at tree level the matching conditions impose only three relations
$`C_1^{\mathrm{latt}}`$ $`=C_3^{\mathrm{latt}}+𝒪(\alpha _s),`$
$`C_2^{\mathrm{latt}}`$ $`=C_4^{\mathrm{latt}}+𝒪(\alpha _s),`$ (29)
$`C^{\mathrm{cont}}`$ $`=C_1^{\mathrm{latt}}+C_2^{\mathrm{latt}}+𝒪(\alpha _s),`$
among the four lattice operator Wilson coefficients. Here, $`C^{\mathrm{cont}}`$ is the Wilson coefficient of the continuum operator $`Q`$, while $`C_i^{\mathrm{latt}}`$ are the coefficients of $`Q_i`$.
In lattice simulations, the quarks in the operator $`Q`$ are treated as four distinct flavors Sharpe:1986xu . In this case, there is no ambiguity: the effective continuum Hamiltonian contains operators $`Q_i^{\mathrm{cont}}`$ that are the continuum limit of each $`Q_i`$. Also, instead of writing relations among coefficients, it is more convenient to rewrite the result of the matching in the form
$$Q_i^{\mathrm{cont}}\underset{j=1}{\overset{4}{}}C_{ij}Q_j^{\mathrm{lat}},$$
(30)
where “$``$” means equality at the level of matrix elements, see (12). The 4 by 4 matrix $`C_{ij}`$ has a perturbative expansion
$$C_{ij}=\delta _{ij}+\left(\frac{\alpha _s}{\pi }\right)C_{ij}^{(1)}+\mathrm{}$$
(31)
and is obtained from the results for the graphs in the previous section by
$`C_{ij}^{(1)}`$ $`={\displaystyle \frac{1}{4}}\left[\delta Z^{(1)}\delta _{ij}+\left({\displaystyle \frac{a\mu }{2}}\right)^{2ϵ}{\displaystyle \underset{\alpha }{}}q_{\alpha ,j}^i+\delta Z_{ij}^{(1)}\right].`$ (32)
Here $`\delta Z^{(1)}`$ stands for the wave function renormalization and $`Z_{ij}=\delta _{ij}+\frac{\alpha _s}{\pi }Z_{ij}^{(1)}`$ is the matrix of the $`\overline{\mathrm{MS}}`$ renormalization constants for the operators $`Q_i^{\mathrm{cont}}`$ in the continuum. For the purposes of computing $`C_{ij}^{(1)}`$, the role of this matrix is to remove any residual $`1/ϵ`$ dependence in (32). Using the results for the matrices $`q_\alpha `$ given in the previous section, it is straightforward to compute the matching coefficients. We split the matrix $`C_{ij}^{(1)}`$ into three parts
$$C^{(1)}=\mathrm{ln}\left(\frac{a\mu }{2}\right)\gamma ^{(1)}+C_C^{(1)}+C_U^{(1)}.$$
(33)
Here, $`\gamma ^{(1)}`$ is the universal one-loop anomalous dimensions matrix of the operators $`Q_i`$:
$$\gamma ^{(1)}=\left(\begin{array}{cccc}0& 0& \frac{1}{2}& \frac{3}{2}\\ \frac{3}{4}& \frac{9}{4}& \frac{3}{4}& \frac{7}{4}\\ \frac{1}{2}& \frac{3}{2}& 0& 0\\ \frac{3}{4}& \frac{7}{4}& \frac{3}{4}& \frac{9}{4}\end{array}\right).$$
(34)
The remainder $`C_C+C_U`$ depends on the lattice action and we give it for three cases: (a) Asqtad fermion and improved gauge action, (b) Asqtad fermion with plaquette gluon action, (c) standard action, no improvement. The reason for splitting the remainder into two pieces which are not separately gauge invariant, is as follows. Occasionally, simulations are performed with operators that are obtained by dropping the links that connect fermions on different sites. The simulations with these gauge non-invariant operators are performed in Landau gauge and their matching coefficients can be extracted from our calculation. To obtain them, we need to account for the contributions from graphs $`13`$ in Fig. 1, the wave function renormalization contributions as well as the tadpole improvement for the wave function given in (28). We denote the corresponding matching matrix by $`C_C`$. The remainder, denoted by $`C_U`$ is the contribution from graphs 4 and 5 which arises from the link fields in the current operator and the operator tadpole improvement term given in (27). We give the matrices $`C_{C,U}`$ in Landau gauge; their sum is gauge invariant.
For the Asqtad and improved gauge action, we find
$`C_C^{(1a)}`$ $`=\left(\begin{array}{cccc}2.83(14)& 0.75& 0.38(1)& 1.15(2)\\ 1.25(1)& 4.33(10)& 0.58(1)& 1.34(3)\\ 0.38(1)& 1.15(2)& 2.59(15)& 0.75\\ 0.58(1)& 1.34(3)& 1.34(2)& 4.36(10)\end{array}\right),`$ (39)
$`C_U^{(1a)}`$ $`=\left(\begin{array}{cccc}2.08(5)& 0& 0& 0\\ 0.29& 0.98(3)& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0.20(1)& 1.14(3)\end{array}\right).`$ (45)
These results are obtained after tadpole improvement with the “mean link in Landau gauge” definition of the improvement parameter, $`u_0^{(1)}=0.750(2)`$ for the improved gluon action and $`u_0^{(1)}=0.97432`$ for the unimproved action. For Asqtad with the standard plaquette gluon action, the result is
$`C_C^{(1b)}`$ $`=\left(\begin{array}{cccc}2.33& 0.75& 0.34& 1.03\\ 1.18& 3.61(1)& 0.51& 1.20\\ 0.34& 1.03& 2.09& 0.75\\ 0.51& 1.20& 1.27& 3.64(1)\end{array}\right),`$ (50)
$`C_U^{(1b)}`$ $`=\left(\begin{array}{cccc}2.97& 0& 0& 0\\ 0.32& 1.43& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0.20& 1.60\end{array}\right).`$ (56)
Notice that case (b) is not equivalent to what is considered in Lee:2003sk . These authors employ a different discretization of the operators using insertions of fattened links to make the operators gauge invariant while we are using thin links. Also, their fermion action is similar but not identical to the fermion part of the Asqtad action.
Without any improvement terms in the action, we get
$`C_C^{(1c)}`$ $`=\left(\begin{array}{cccc}0.81& 0.75& 0.18& 0.54\\ 0.85& 0.5& 0.27& 0.63\\ 0.18& 0.54& 0.25& 0.75\\ 0.27& 0.63& 0.46& 0.63\end{array}\right),`$ (61)
$`C_U^{(1c)}`$ $`=\left(\begin{array}{cccc}3.83& 0& 0& 0\\ 0.47& 1.86& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0.38& 2.14\end{array}\right).`$ (67)
The matrices for the unimproved action agree with the known results in the literature Ishizuka:1993fs ; Lee:2003sk .
We write the finite correction to the normalization of the axial-vector current in the form $`Z_A=1+\frac{\alpha _s}{\pi }\delta Z_A`$ and obtain
$`\delta Z_A^{(a)}`$ $`=1.17(7),`$ $`\delta Z_A^{(b)}`$ $`=0.921,`$ $`\delta Z_A^{(c)}`$ $`=0.`$ (68)
The current renormalization happens to be identical for the gauge invariant and non-invariant operators in Landau gauge. This occurs because in this gauge, the contribution of the fourth graph vanishes and the contribution from fifth graph gets canceled by tadpole improvement.
## VI Summary and discussion
We have calculated the matching of the Wilson coefficients in the $`\mathrm{\Delta }S=2`$ flavor changing effective Hamiltonian from the continuum to the lattice for the Asqtad lattice action. Together with non-perturbative calculations of the $`K\overline{K}`$ matrix element of the lattice operators, our results determine the amount of indirect CP violation in the kaon system. The matrix elements are currently being evaluated by the HPQCD and UKQCD collaborations Gamiz:2004qx ; CDGSW05 . The calculation uses the configurations with $`n_f=2+1`$ dynamical flavors generated by the MILC collaboration.
Previous calculations of the quantity $`B_K`$, which determines the strength of $`K\overline{K}`$ mixing, were performed in the quenched approximation. This induces an essentially unknown and irreducible systematic error into the result. Precise simulations with dynamical fermions are necessary in order to be able to make full use of the experimental data on $`K\overline{K}`$ mixing to constrain the CKM matrix elements entering the effective Hamiltonian. Among the currently available fermion actions, only staggered actions allow for simulations with low statistical errors at small quark masses. However, the unimproved staggered action suffers from large taste-changing interactions. The Asqtad action, for which we have performed the matching, has been designed to reduce these as well as other cut-off effects.
We have performed our calculation with two independent methods. First, we have evaluated the diagrams in two different ways: by directly calculating the various diagrams and by first separating off the part which can be inferred from the renormalization of the current operators. Second, we have evaluated the lattice integrals both algebraically and numerically. For the algebraic evaluation, we have expanded the diagrams around the continuum limit. This produces a set of lattice tadpole integrals which we then reduced to a minimal set of master integrals using computer algebra. The agreement between these two rather different methods provides a strong check on our results.
We find that the size of one-loop corrections to the tree level matching for the Asqtad action is very similar to what is found with unimproved staggered and other improved staggered actions, such as the HYP action Lee:2004qv . In quenched calculations at a lattice spacing $`1/a=1.6\mathrm{GeV}`$, the shift in the value of the bare $`B_K`$ due to $`𝒪(\alpha _s)`$ corrections is around 10% for all these actions CDGSW05 . In dynamical simulations with the Asqtad action the shift turns out to be over 15% at the same lattice spacing, simply because the value of the strong coupling is much larger for $`n_f=3`$ than for $`n_f=0`$ dynamical flavors. Also, it turns out that the values for the matching coefficients are not very sensitive to the improvement of the gluon action. Switching off the gluon improvement terms in the Asqtad action changes the entries of the one-loop matching matrix by less than 20% and the changes are even smaller for the larger elements on the diagonal. The motivation for using an improved staggered action in the calculation of $`B_K`$ is not to reduce the size of the one-loop corrections, which are already small in the unimproved case, but to correct the large $`𝒪(a^2)`$ scaling violations that highly affect the unimproved case; see Gamiz:2004qx ; CDGSW05 for further discussion.
The situation is different for the calculation of other weak matrix elements relevant in the study of CP-violating effects in the kaon system. For the unimproved staggered action, Sharpe and Patel Sharpe:1993ur found very large perturbative corrections for operators with scalar and tensor Dirac structure. The contributions are large even after tadpole improvement, which casts doubt on the perturbative expansion for the unimproved action. The corrections are reduced to an acceptable level with the HYP action Hasenfratz:2001hp and operators which are made gauge invariant by the insertion of fattened links LS02 ; Lee:2003sk . In order to check whether these large corrections are present for the Asqtad action with standard four-quark operators, we have analyzed the mixing of scalar operators. We find that the corrections are small with the improved action. The suppression arises because the improvement terms reduce taste-changing interactions. In Appendix A we present the results for the mixing of scalar operators and discuss the suppression of the large contributions in more detail.
The absence of anomalously large one-loop corrections suggests that the two-loop corrections to our results should be reasonably small. However, in order to reduce the uncertainty in dynamical calculations to the level of a few per-cent, a two-loop matching calculation will be necessary.
Acknowledgments We are grateful to Christine Davies and Andreas Kronfeld for discussions and comments on the manuscript. We thank the KITP in Santa Barbara and the INT in Seattle where part of this work was carried out for their hospitality and financial support. This research was supported in part by the U.S. Department of Energy contracts DE-AC02-76CH03000, DE-FG03-94ER-40833 and the Outstanding Junior Investigator Award DE-FG03-94ER-40833, and by the Alfred P. Sloan Foundation. E.G. is indebted to the European Commission for a Marie-Curie Grant No. MEIF-CT-2003-501309. Fermilab is operated by Universities Research Association Inc., under contract with the U.S. Department of Energy.
## Appendix A Mixing of scalar operators
In this appendix we evaluate the matching for the operators
$`Q_1^S`$ $`=(\overline{\psi }_1^a\psi _2^a)_S(\overline{\psi }_3^b\psi _4^b)_S,`$ (69)
$`Q_2^S`$ $`=(\overline{\psi }_1^a\psi _2^b)_S(\overline{\psi }_3^b\psi _4^a)_S,`$
where the bilinears $`(\overline{\psi }_i^a\psi _j^a)_S`$ have unit Dirac and taste structure, $`\gamma _S\xi _T=11`$ (or equivalently $`\gamma _5\gamma _5`$). These operators are not needed for the evaluation of $`B_K`$, but scalar and pseudoscalar operators are present in the $`\mathrm{\Delta }S=1`$ effective Lagrangian used to determine $`\epsilon ^{}`$. The scalar operators mix into tensor operators under renormalization, but the form of this mixing is not important for our discussion.
For gauge non-invariant operators, the above matching has been studied by Sharpe and Patel Sharpe:1993ur , who find very large perturbative corrections for these operators even after tadpole improvement. Such large corrections are absent for the HYP action Hasenfratz:2001hp with operators which are made gauge invariant by the insertion of fattened links LS02 ; Lee:2003sk . As we show below, this is true also for the Asqtad action with standard operators.
Using the same notations and conventions as in Section V, the $`2\times 2`$ matrix for the anomalous dimension of the operators in (69) is
$$\gamma ^{(1)}=\left(\begin{array}{cc}4& 0\\ \frac{3}{2}& \frac{1}{2}\end{array}\right).$$
(70)
After tadpole improvement, the corresponding matching matrices for the unimproved action are
$`C_C^{(1c)}`$ $`=\left(\begin{array}{cccc}17.26& 0& & \\ 6.28& 1.59& & \end{array}\right),`$ (73)
$`C_U^{(1c)}`$ $`=\left(\begin{array}{cccc}0& 0& & \\ 2.45& 2.83& & \end{array}\right).`$ (77)
The one-loop correction to the matching is indeed very large, in particular for the element that describes the mixing of $`Q_1^S`$ into itself. The large correction is entirely due to the first diagram in Fig. 1. It is twice the renormalization of the scalar current. Note that the operator $`Q_1^S`$ is completely local. In terms of the naive field, it is simply
$$\overline{q}(11)q=\frac{1}{16}\underset{A}{}\overline{\psi }(A)\psi (A).$$
(78)
This implies that for the unimproved action large values are obtained, independently of the prescription adopted to make the operators gauge invariant.
The values obtained for the Asqtad action are
$`C_C^{(1a)}`$ $`=\left(\begin{array}{cccc}1.8(1)& 0& & \\ 0.6(1)& 3.6(1)& & \end{array}\right),`$ (81)
$`C_U^{(1a)}`$ $`=\left(\begin{array}{cccc}0& 0& & \\ 1.81(1)& 1.68(3)& & \end{array}\right).`$ (85)
It is comforting to see that the corrections are smaller with the improved action. The dramatic reduction is mostly due to the fat-link terms in the fermion action.
The origin of the large corrections for the unimproved action, such as those in (73), has been analyzed by Golterman Golterman:1998jj . He splits the integration region for the loop momenta $`\pi <k_\mu \pi `$ into a region around zero $`\pi /2<k_\mu \pi /2`$ and a remainder which contains the corners of the Brillouin zone. For Wilson fermions, he finds that after tadpole improvement the main contribution arises from the integration region around zero. For staggered fermions, on the other hand, large corrections arise from the corners of the Brillouin zone. In this region, the gluon propagator is off-shell, but the staggered fermion propagator has poles.<sup>3</sup><sup>3</sup>3Note that the contribution we are talking about is from the region of hard loop momentum; there are no singularities from propagating doublers. Such contributions only arise at higher order in the expansion in the lattice spacing. Because the gluon propagator is far off-shell, the first graph in Fig. 2 can be viewed as a fermion tadpole in the corner region. This is the interpretation put forward by Golterman.
The reason for the reduction achieved with the improved action is that the improved quark gluon coupling is designed to suppress taste-changing interactions. The quark-gluon interaction switches off if the in- and outgoing tastes differ. This mechanism suppresses the contribution from the corners of the Brillouin zone to the first diagram in Fig. 1, which corresponds precisely to the situation where the taste of the internal quark line is different from the taste of the external line.
To conclude, we find that the anomalously large perturbative corrections present in some cases for the unimproved staggered action are suppressed in the Asqtad results due to the smaller taste-symmetry breaking of this action. We thus expect that the two-loop corrections to our matching calculation will be reasonably small. |
warning/0507/cond-mat0507265.html | ar5iv | text | # Magnetic reordering in the vicinity of a ferromagnetic/antiferromagnetic interface
## 1 Introduction
Magnetic reordering in the vicinity of an interface has for a long time attracted the interest of researchers. In fact, when two magnetically ordered systems are in atomic contact with each other, it is quite natural to expect that in the vicinity of the interface a novel magnetic arrangement, different from the bulk one, will set in. This phenomenon is usually referred to as the magnetic proximity effect (MPE). To the best of our knowledge this effect was first investigated to treat a ferromagnet in contact with a paramagnet Zuc73 . Since then a vast literature on the subject has been published, of which we mention just a few examples KiZ73 .
The interest in the MPE has revived lately in relation to the exchange bias effect NoS99 . It occurs when a thin ferromagnetic (FM) film is deposited on an antiferromagnetic (AFM) material, resulting in a shift of the hysteresis loop from its normal (symmetric) position. If the AFM has a compensated interface (’in-plane AFM’), i.e., if the number of bonds between parallel and antiparallel spin pairs across the interface is the same, the AFM often assumes an almost orthogonal magnetization with respect to the FM magnetic direction, while the spins of the AFM interface layer adopt a canted configuration. This magnetic arrangement of the AFM is usually called spin-flop-phase, in analogy to an AFM system in an external magnetic field Nee67 . The occurrence of the exchange bias effect is, in most likelihood, related to a certain amount of interface disorder MNL98 .
When considering coupled FM-AFM systems we realized that results for fully ordered structures are scarce. In previous studies the FM film is usually treated as a system with uniform layers, i.e., the spins within a given FM layer remain strictly parallel to each other Koo97 ; HiM86 ; CrC01 ; JeD02 . Whereas different magnetization directions for different FM layers are considered, each layer rotates solidly. We stress that in the case of a compensated FM-AFM interface an MPE may be present also for the FM layers close to the interface. Thus the magnetic structure of each FM layer is represented, in perfect analogy to the AFM layers, by two interpenetrating sublattices with different magnetization directions. The consideration of a nonuniform intralayer magnetic structure in the FM subsystem leads to new features, which in turn are strongly dependent on the underlying lattice symmetry.
Results concerning the spin reorientation of full magnetic layers have been obtained previously for various magnetic systems but, to the best of our knowledge, caused by magnetic anisotropies HiM86 ; CrC01 ; JeD02 . It is important to stress that although we also incorporate anisotropy, the spin rotation in ordered FM-AFM films is mainly caused by the isotropic exchange interactions. Moreover, these systems exhibit a rotation of the magnetic sublattices, and not a net spin reorientation of the full layers. These properties constitute an essential difference of the present treatment with respect to previous studies.
In order to derive a number of general results, while keeping the analysis as straightforward as possible, we examine at first the magnetic arrangement of a perfectly ordered bilayer consisting of a single FM layer that is coupled to a single AFM layer. Using a mean field approximation, this particular structure yields results which can be written in an analytical form. In addition, we present a number of results for more realistic systems having thicker FM and AFM films. In particular, we investigate the effect of the interface coupling on the characteristics and magnitude of the MPE at zero and finite temperatures. Of special concern is whether, and to which degree, magnetic order is induced by the subsystem with the higher (bare) ordering (Néel or Curie) temperature into the one with the lower ordering temperature. The resulting magnetic arrangements for various cases of the bilayer system, for films with several atomic layers, and for the corresponding ordering temperatures are determined. In fact, we show that, depending on the lattice structure, the proximity effect is not always present, and that under certain circumstances two different critical temperatures can occur.
This paper is organized as follows. In Section 2 we define our physical model. In Section 3 the magnetic properties of the bilayer system at zero temperature are discussed, which exhibits already a number of general features. Results obtained for finite temperatures are presented in Section 4. Thicker films with several FM and AFM layers are considered in Section 5. Conclusions are drawn in the last Section.
## 2 Theory
To model the magnetic arrangement and ordering temperatures of a coupled FM-AFM system we use an XYZ- Heisenberg Hamiltonian with localized quantum spins $`𝐒_i`$ and spin number $`S`$,
$$=\frac{1}{2}\underset{i,j}{}\left(J_{ij}𝐒_i𝐒_j+D_{ij}^xS_i^xS_j^x+D_{ij}^yS_i^yS_j^y\right).$$
(1)
We take into account the isotropic exchange interaction $`J_{ij}`$ between spins located on nearest-neighbor lattice sites $`i`$ and $`j`$. In addition in-plane easy-axis exchange anisotropies $`D_{ij}^x`$ and $`D_{ij}^y`$ are considered, which for a particular layer are directed either along the $`x`$\- or along the $`y`$-direction. Note that for two-dimensional (2D) magnets a long-range magnetic order at finite temperatures exists only in presence of such anisotropies MeW66 . A perfectly ordered layered structure in the $`xy`$-plane is assumed, consisting of an FM film with $`n_{\mathrm{FM}}`$ layers and an AFM film with $`n_{\mathrm{AFM}}`$ layers. Each layer is represented by two interpenetrating sublattices, applying otherwise periodic lateral boundary conditions. The lattice symmetry, which is assumed to be the same for both FM and AFM films, is characterized by the numbers of nearest neighbors $`z_0`$ and $`z_1`$ within a layer and between adjacent layers, respectively. The latter value also refers to the number of bonds with which an FM spin is coupled across the interface to neighboring spins in the AFM layer. In this study the sc(001) and fcc(001) lattices are taken as representative and extremal examples corresponding to $`z_1=1`$ and $`z_1=4`$, respectively, and $`z_0=4`$ for both symmetries hexagonal . As will become apparent in the next Sections, the magnetic properties of these two types of coupled FM-AFM films differ markedly.
The FM and AFM subsystems are characterized by the exchange couplings $`J_{\mathrm{FM}}>0`$ and $`J_{\mathrm{AFM}}<0`$, and by the usually much weaker exchange anisotropies $`D_{\mathrm{FM}}>0`$ along the $`x`$-axis and $`D_{\mathrm{AFM}}<0`$ along the $`y`$-axis. Due to shape anisotropy the magnetizations of both subsystems are confined to the film plane, besides this demagnetizing effect the magnetic dipole interaction is not considered explicitely dipol . Furthermore, the FM and AFM films are coupled across the interface by the interlayer exchange coupling $`J_{\mathrm{int}}`$, where we consider $`J_{\mathrm{int}}>0`$ without loss of generality, and $`D_{\mathrm{int}}=0`$. The (unperturbed) ground state for a small interface coupling $`J_{\mathrm{int}}0`$ is defined by a mutually perpendicular arrangement of the FM and AFM magnetic directions. The choice of the anisotropy easy axes support this perpendicular magnetic arrangement. A net magnetic binding results only if the spins of at least one of the subsystems are allowed to deviate from the unperturbed state. Hence, the magnetic moments cannot be represented by Ising-like spins.
In this study we apply a single-spin mean field approximation (MFA). Within this method the site-dependent magnetizations $`𝐒_i=𝐌_i(T)`$ with components $`M_i^x(T)`$ and $`M_i^y(T)`$ are calculated, yielding the magnitudes $`M_i(T)=|𝐌_i(T)|`$ and in-plane angles $`\mathrm{tan}\varphi _i(T)=\pm M_i^y(T)/M_i^x(T)`$. Furthermore, the ordering temperatures are determined. For decoupled monolayers ($`J_{\mathrm{int}}=0`$, $`n_{\mathrm{FM}}=n_{\mathrm{AFM}}=1`$) the bare Curie temperature $`T_C^0`$ of the FM and the analogous Néel temperature $`T_N^0`$ of the AFM are given by
$`T_C^0`$ $`=`$ $`{\displaystyle \frac{S(S+1)}{3}}z_0(J_{\mathrm{FM}}+D_{\mathrm{FM}}),`$
$`T_N^0`$ $`=`$ $`{\displaystyle \frac{S(S+1)}{3}}z_0|J_{\mathrm{AFM}}+D_{\mathrm{AFM}}|,`$ (2)
where the Boltzmann constant $`k_B`$ is set equal to unity. For thicker films the ordering temperature is determined by the largest eigenvalue of a particular matrix. We will investigate the MPE for a number of different cases, i.e., whether and to which degree a magnetic order propagates from the subsystem with the larger bare ordering temperature into the other one. The corresponding magnetic structure is characterized by the magnetization vectors $`𝐌_i(T)`$. As mentioned, at first we will consider the particularly simple bilayer system ($`n_{\mathrm{FM}}=n_{\mathrm{AFM}}=1`$) which allows to draw a number of general results and analytical expressions. Later on we take into account coupled FM-AFM systems with thicker films. Since the explicit consideration of anisotropies is not decisive within MFA, for simplicity we include them into the exchange couplings: $`J_{\mathrm{FM}}+D_{\mathrm{FM}}J_{\mathrm{FM}}`$ and $`J_{\mathrm{AFM}}+D_{\mathrm{AFM}}J_{\mathrm{AFM}}`$. For the spin quantum number we use $`S=1`$ throughout.
## 3 Bilayers: Zero temperature
### 3.1 sc(001) – bilayer
For this lattice type both the FM and AFM layers are disturbed from their ground state, thus also the FM layer ‘dimerizes’ and exhibits a noncollinear magnetization. The undisturbed magnetic arrangement of an sc(001) bilayer is depicted in Figure 1a. For $`J_{\mathrm{int}}>0`$ both FM and AFM layers assume a canted magnetic arrangement, as sketched in Figure 1b. The canting angles $`\varphi _{\mathrm{FM}}`$ and $`\varphi _{\mathrm{AFM}}`$ represent the deviations from the decoupled bilayer.
The energy of such an arrangement is given by
$`E_{\mathrm{sc}(001)}(\varphi _{\mathrm{FM}},\varphi _{\mathrm{AFM}})`$ $`=`$ (3)
$`{\displaystyle \frac{z_0}{2}}J_{\mathrm{FM}}\mathrm{cos}(2\varphi _{\mathrm{FM}}){\displaystyle \frac{z_0}{2}}|J_{\mathrm{AFM}}|\mathrm{cos}(2\varphi _{\mathrm{AFM}})`$
$`z_1J_{\mathrm{int}}\mathrm{cos}(\pi /2\varphi _{\mathrm{FM}}\varphi _{\mathrm{AFM}}).`$
Differentiation of $`E_{\mathrm{sc}(001)}(\varphi _{\mathrm{FM}},\varphi _{\mathrm{AFM}})`$ with respect to $`\varphi _{\mathrm{FM}}`$ and $`\varphi _{\mathrm{AFM}}`$ yields the conditions for the equilibrium angles $`\varphi _{\mathrm{FM}}^0`$ and $`\varphi _{\mathrm{AFM}}^0`$,
$`z_0J_{\mathrm{FM}}\mathrm{sin}(2\varphi _{\mathrm{FM}}^0)`$ $`=`$ $`z_0|J_{\mathrm{AFM}}|\mathrm{sin}(2\varphi _{\mathrm{AFM}}^0)`$ (4)
$`=z_1J_{\mathrm{int}}\mathrm{cos}(\varphi _{\mathrm{FM}}^0+\varphi _{\mathrm{AFM}}^0).`$
We emphasize that this behavior refers to a magnetic rotation of the two sublattices, with angles $`\varphi _i^0`$ and $`\pi \varphi _i^0`$, and not to a net spin reorientation of layer $`i`$.
First we consider $`|J_{\mathrm{AFM}}|<J_{\mathrm{FM}}`$. In Figure 2 the angles $`\varphi _{\mathrm{FM}}^0`$ and $`\varphi _{\mathrm{AFM}}^0`$ are shown as functions of the interlayer coupling $`J_{\mathrm{int}}`$ for different values of $`|J_{\mathrm{AFM}}|`$. The following properties are quite apparent:
* For a small $`|J_{\mathrm{AFM}}|`$ the AFM spins quickly turn into the direction of the FM as $`J_{\mathrm{int}}`$ increases. A parallel orientation of the AFM spins with respect to the FM, i.e., $`\varphi _{\mathrm{AFM}}^0=90^{}`$ and $`\varphi _{\mathrm{FM}}^0=0^{}`$, is reached at the particular strength $`J_{\mathrm{int}}^{}`$ of the interlayer coupling, given by
$$J_{\mathrm{int}}^{}=\frac{z_0}{z_1}\frac{2J_{\mathrm{FM}}|J_{\mathrm{AFM}}|}{J_{\mathrm{FM}}|J_{\mathrm{AFM}}|}.$$
(5)
The larger $`|J_{\mathrm{AFM}}|`$ the larger is the value of $`J_{\mathrm{int}}^{}`$ required to reach that limit.
* With increasing $`J_{\mathrm{int}}`$ the FM angle $`\varphi _{\mathrm{FM}}^0`$ increases and exhibits a maximum at
$$J_{\mathrm{int}}^{\mathrm{max}}=\frac{z_0}{z_1}|J_{\mathrm{AFM}}|\sqrt{\frac{2J_{\mathrm{FM}}}{J_{\mathrm{FM}}|J_{\mathrm{AFM}}|}},$$
(6)
assuming the value $`\mathrm{sin}(2\varphi _{\mathrm{FM}}^{\mathrm{max}})=|J_{\mathrm{AFM}}|/J_{\mathrm{FM}}`$ and coinciding with $`\varphi _{\mathrm{AFM}}^0=45^{}`$. Notice that $`\varphi _{\mathrm{FM}}^0(J_{\mathrm{int}})`$ and $`\varphi _{\mathrm{AFM}}^0(J_{\mathrm{int}})`$ in general are not symmetric with respect to $`J_{\mathrm{int}}^{\mathrm{max}}`$.
* For the limiting case $`|J_{\mathrm{AFM}}|=J_{\mathrm{FM}}`$, no maximum of $`\varphi _{\mathrm{FM}}^0`$ is obtained. Instead one has $`\mathrm{tan}(2\varphi _{\mathrm{FM}}^0)=\mathrm{tan}(2\varphi _{\mathrm{AFM}}^0)=(z_1J_{\mathrm{int}})/(z_0J_{\mathrm{FM}})`$ For $`J_{\mathrm{int}}\mathrm{}`$ one obtains $`\varphi _{\mathrm{FM}}^0=\varphi _{\mathrm{AFM}}^0=45^{}`$.
* For $`J_{\mathrm{int}}<0`$ the same results emerge, if one performs the transformations $`\varphi _{\mathrm{FM}}^0\varphi _{\mathrm{FM}}^0`$ and $`\varphi _{\mathrm{AFM}}^0\varphi _{\mathrm{AFM}}^0`$.
* The sc(001) bilayer is characterized by an apparent symmetry between the FM and AFM layers as determined within MFA. For $`J_{\mathrm{FM}}<|J_{\mathrm{AFM}}|`$ the behavior of the FM and AFM layers, in particular the equilibrium angles $`\varphi _{\mathrm{FM}}^0`$ and $`\varphi _{\mathrm{AFM}}^0`$, becomes interchanged, as can be seen from the symmetry of equations (3) and (4). If one exchanges $`J_{\mathrm{FM}}`$ and $`|J_{\mathrm{AFM}}|`$ in the preceding deduction, Figure 2a is valid for $`\varphi _{\mathrm{FM}}^0`$ and Figure 2b for $`\varphi _{\mathrm{AFM}}^0`$. Thus, for $`J_{\mathrm{int}}>J_{\mathrm{int}}^{}`$, antiferromagnetic order of the FM layer on top of the undisturbed AFM layer sets in. Another system exhibiting this behavior is the bcc(110) bilayer.
### 3.2 fcc(001) – bilayer
The fcc(001) bilayer is characterized by the fact that for an undisturbed AFM the sum of the coupling energies to a given FM spin vanishes, as can be seen from the undisturbed arrangement illustrated in Figure 3a. By setting up an equation similar to equation (3) one can show that $`\varphi _{\mathrm{FM}}^0=0`$, hence in this case the spin structure of the FM always remains strictly collinear. The resulting magnetic structure of a coupled fcc(001) bilayer is shown in Figure 3b. Thus, the symmetry between the FM and AFM subsystems of the sc(001) bilayer is no longer present for the fcc(001) one. This is a consequence of the fact that for the sc(001) bilayer each FM spin couples to a single AFM sublattice, while for the fcc(001) interface each FM spin couples identically to both AFM sublattices. A similar behavior holds for bcc(001) films.
The corresponding energy expression $`E_{\mathrm{fcc}(001)}`$ is obtained from equation (3) by setting $`\varphi _{\mathrm{FM}}=0`$. Differentiation with respect to $`\varphi _{\mathrm{AFM}}`$ yields the equilibrium angle $`\varphi _{\mathrm{AFM}}^0`$ of the disturbed AFM spin arrangement,
$$\mathrm{sin}(\varphi _{\mathrm{AFM}}^0)=\frac{z_1J_{\mathrm{int}}}{2z_0|J_{\mathrm{AFM}}|},$$
(7)
which is shown in Figure 4 as function of $`J_{\mathrm{int}}`$. For $`z_1J_{\mathrm{int}}>2z_0|J_{\mathrm{AFM}}|`$ one obtains $`\varphi _{\mathrm{AFM}}^0=90^{}`$, i.e., the spins of the AFM layer order parallel to the FM ones. The case $`J_{\mathrm{int}}<0`$ is recovered by replacing $`\varphi _{\mathrm{AFM}}^0\varphi _{\mathrm{AFM}}^0`$.
## 4 Bilayer: Finite temperatures
We now turn our attention to the magnetic arrangement of the coupled FM-AFM bilayer at finite temperatures. Like in the previous Section we distinguish between an sc(001) and an fcc(001) symmetry. Furthermore, we treat the cases $`T_N^0<T_C^0`$, $`T_N^0>T_C^0`$, and $`T_N^0=T_C^0`$ separately.
Let us at first present the ordering temperature $`T_C`$ for a coupled magnetic bilayer with a collinear magnetization. Its two layers are characterized by the exchange couplings $`J_1`$ and $`J_2`$, which can be of either sign. Within MFA one obtains
$`T_C`$ $`=`$ $`{\displaystyle \frac{S(S+1)}{6}}[z_0(J_1+J_2)`$ (8)
$`+\sqrt{z_0^2(J_1J_2)^2+4(z_1J_{\mathrm{int}})^2}].`$
Except for the cases that will be mentioned below, $`T_C`$ of the coupled bilayer is always larger than the largest bare ordering temperature ($`T_N^0`$ or $`T_C^0`$) of the decoupled monolayers, regardless of the relative magnitude of $`J_1`$ and $`J_2`$, and of the sign of $`J_{\mathrm{int}}`$. For unequal layers ($`J_1J_2`$) and a small coupling $`J_{\mathrm{int}}`$ one obtains an increase of $`T_C`$ given approximately by
$$\mathrm{\Delta }T_C(J_{\mathrm{int}})\frac{S(S+1)}{3}\frac{(z_1J_{\mathrm{int}})^2}{z_0|J_1J_2|}.$$
(9)
From the denominator of equation (9) one observes that the increase of $`T_C`$ for an FM bilayer ($`J_1,J_2>0`$) will be larger than the one for a corresponding FM-AFM bilayer ($`J_1>0`$, $`J_2<0`$). Within MFA the results for $`J_{\mathrm{int}}<0`$ are identical to the corresponding ones for $`J_{\mathrm{int}}>0`$, if the signs of $`\varphi _{\mathrm{FM}}^0(T)`$ and $`\varphi _{\mathrm{AFM}}^0(T)`$ are adapted appropriately.
### 4.1 sc(001) – bilayer
a) $`T_N^0<T_C^0`$. For the AFM coupling we choose $`|J_{\mathrm{AFM}}|/J_{\mathrm{FM}}`$ $`=0.5`$. In Figure 5 we display the magnetizations $`M_{\mathrm{FM}}(T)`$ and $`M_{\mathrm{AFM}}(T)`$, and the corresponding equilibrium angles $`\varphi _{\mathrm{FM}}^0(T)`$ and $`\varphi _{\mathrm{AFM}}^0(T)`$, as functions of the temperature $`T`$. Different values of the interlayer coupling $`J_{\mathrm{int}}>0`$ are used as indicated. At low temperatures both subsystems deviate from the undisturbed magnetic arrangement. With increasing temperature the equilibrium angle $`\varphi _{\mathrm{FM}}^0(T)`$ of the FM layer decreases, whereas $`\varphi _{\mathrm{AFM}}^0(T)`$ of the AFM layer increases. Approaching the sublattice reorientation temperature $`T_R`$, given implicitly by the relation
$`z_1J_{\mathrm{int}}\left[J_{\mathrm{FM}}M_{\mathrm{FM}}^2(T_R)|J_{\mathrm{AFM}}|M_{\mathrm{AFM}}^2(T_R)\right]`$ (10)
$`=`$ $`2z_0J_{\mathrm{FM}}|J_{\mathrm{AFM}}|M_{\mathrm{FM}}(T_R)M_{\mathrm{AFM}}(T_R),`$
the AFM spins turn into the direction of the FM spins, and one obtains $`\varphi _{\mathrm{FM}}^0(T)0^{}`$, $`\varphi _{\mathrm{AFM}}^0(T)90^{}`$. Thus, for $`T>T_R`$ the AFM layer adopts ferromagnetic order. $`M_{\mathrm{AFM}}(T)`$ exhibits a sharp kink at $`T_R`$, whereas $`M_{\mathrm{FM}}(T)`$ shows no particular features. The FM-AFM bilayer becomes paramagnetic above the ordering temperature $`T_C`$ given by equation (8).
b) $`T_N^0>T_C^0`$. Here the strength of the AFM coupling is stronger than the FM one. We adopt $`|J_{\mathrm{AFM}}|/J_{\mathrm{FM}}=2`$ for comparison with the previous case. Then the results for $`T_N^0>T_C^0`$ are fully symmetric with the ones derived for $`T_N^0<T_C^0`$ shown in Figure 5, if one interchanges $`M_{\mathrm{FM}}(T)M_{\mathrm{AFM}}(T)`$, $`\varphi _{\mathrm{FM}}^0(T)\varphi _{\mathrm{AFM}}^0(T)`$, and $`T_N^0T_C^0`$. Thus new figures are not required. Notice that for $`T>T_R`$ the FM layer assumes an antiferromagnetic structure. This finding demonstrates the symmetry of the FM and AFM layers of the sc(001) lattice within MFA, which also holds for finite temperatures.
c) $`T_N^0=T_C^0`$. For the particular case $`J_{\mathrm{AFM}}=J_{\mathrm{FM}}`$ the angles $`\varphi _{\mathrm{FM}}^0(T)`$ and $`\varphi _{\mathrm{AFM}}^0(T)`$ are independent of the temperature and are given by $`\mathrm{tan}(2\varphi _{\mathrm{FM}}^0)=\mathrm{tan}(2\varphi _{\mathrm{AFM}}^0)=(z_1J_{\mathrm{int}})/(z_0J_{\mathrm{FM}})`$. The magnetizations $`M_{\mathrm{FM}}(T)`$ and $`M_{\mathrm{AFM}}(T)`$ are identical and vanish at the ordering temperature, cf. equation (8),
$$T_C=\frac{S(S+1)}{3}z_0J_{\mathrm{FM}}\sqrt{1+j_{\mathrm{int}}^2}.$$
(11)
### 4.2 fcc(001) – bilayer
As mentioned in Section 3, the behavior of the fcc(001) bilayer system is not symmetric, which also holds for finite temperatures. Several cases have to be distinguished, for this purpose we define the crossover interlayer coupling
$$J_{\mathrm{int}}^{}=\frac{z_0}{z_1}\sqrt{2|J_{\mathrm{AFM}}|\left(|J_{\mathrm{AFM}}|J_{\mathrm{FM}}\right)}.$$
(12)
a) $`T_N^0<T_C^0`$. This case is similar to the analogous sc(001) one. However, unlike that system, for the fcc(001) bilayer the spins in the FM layer remain always collinear, i.e., $`\varphi _{\mathrm{FM}}^0(T)=0`$. With increasing temperature the equilibrium angle $`\varphi _{\mathrm{AFM}}^0(T)`$ of the AFM layer increases and approaches $`90^{}`$ for the temperature $`T_R`$ given by
$$z_1J_{\mathrm{int}}M_{\mathrm{FM}}(T_R)=2z_0|J_{\mathrm{AFM}}|M_{\mathrm{AFM}}(T_R).$$
(13)
For $`T>T_R`$ the AFM spins remain in a ferromagnetic structure up to the ordering temperature given by equation (8). This behavior is depicted in Figure 6 for different values of $`J_{\mathrm{int}}`$. Notice that due to the larger number $`z_1`$ of interlayer bonds the influence of the interlayer coupling for the fcc(001) bilayer is more pronounced as compared to the sc(001) system.
b) $`T_N^0>T_C^0`$ and $`J_{\mathrm{int}}>J_{\mathrm{int}}^{}`$. In effect this case is similar to the preceding one, i.e., with increasing temperature the AFM spins rotate into the direction of the FM. However, although the FM exchange is weaker than the AFM exchange in this case, due to the strong interlayer coupling the FM layer dominates the behavior of the AFM, and results in an ordering temperature, cf. equation (8), even larger than $`T_N^0`$. The lack of a similar mechanism for the AFM layer emphasizes the asymmetry of the two subsystems. Results are illustrated in Figure 7 for different values of $`J_{\mathrm{int}}`$.
c) $`T_N^0>T_C^0`$ and $`J_{\mathrm{int}}<J_{\mathrm{int}}^{}`$. The asymmetric behavior of the FM and AFM layers for the fcc(001) bilayer becomes even stronger. As before, the FM layer remains collinear. However, in this case the disturbance of the AFM layer and the angle $`\varphi _{\mathrm{AFM}}^0(T)`$ decrease with increasing temperature, as shown in Figure 8 for different values of $`J_{\mathrm{int}}`$. At the critical (Curie-) temperature $`T_C^{}>T_C^0`$, given by
$$T_C^{}=\frac{S(S+1)}{3}\left[z_0J_{\mathrm{FM}}+\frac{(z_1J_{\mathrm{int}})^2}{2z_0|J_{\mathrm{AFM}}|}\right],$$
(14)
the FM layer becomes paramagnetic, although in principle coupled to a still ordered AFM layer. However, no magnetization is induced in the FM for $`T>T_C^{}`$, since for $`\varphi _{\mathrm{AFM}}^0(T)=0`$ the couplings of an FM spin across the interface to the two AFM sublattices cancel exactly, and since the scalar product of the interlayer exchange coupling, cf. equation (1), vanishes for perpendicularly oriented FM and AFM layers. The AFM layer becomes disordered at $`T_C=T_N^0`$, thus the bilayer ordering temperature is not given by equation (8). Evidently, in this case the interlayer exchange coupling $`J_{\mathrm{int}}`$ is not strong enough to allow the FM layer to dominate the AFM, like in the previous case. Hence, the coupled magnetic system has two critical temperatures. This behavior is present as long as $`T_C^{}`$ is smaller than $`T_N^0`$. Equating $`T_C^{}=T_N^0`$ yields the relation for the crossover interlayer coupling $`J_{\mathrm{int}}^{}`$ given by equation (12).
## 5 Thicker Films
In this Section we will present a number of results for coupled FM-AFM systems, where the individual FM and AFM films are thicker than just a monolayer. Evidently, the magnetizations $`M_{\mathrm{FMi}}(T)`$ and $`M_{\mathrm{AFMi}}(T)`$, and the sublattice canting angles $`\varphi _{\mathrm{FMi}}^0(T)`$ and $`\varphi _{\mathrm{AFMi}}^0(T)`$ will depend on the layer $`i`$. The deviation from the undisturbed magnetic arrangement, cf. Figures 1a and 3a, is expected to be particularly pronounced for the layers close to the interface, whereas will vanish rapidly with increasing distance from the interface. In Figure 9 the equilibrium angles are shown for an sc(001) lattice symmetry at $`T=0`$ as function of the AFM film thickness. For the FM film one and two layers are considered. The angles $`\varphi _\mathrm{i}^0`$, particularly those close to the interface, saturate within two AFM layers, while thicker AFM films exhibit a weak oscillatory behavior of decreasing amplitude which cannot be observed on the scale of the figure. An alternating sign of $`\varphi _{\mathrm{AFMi}}^0`$ is obtained for neighboring AFM layers. For distances from the interface larger than approximately three layers the AFM remains virtually undisturbed. A corresponding behavior is obtained by varying the FM film thickness. Similar results have been reported for instance in Fin04 .
Moreover, we also investigate sc(001) FM-AFM systems with thicker AFM films at finite temperatures. As for the bilayer, and also for thicker films and for $`T_N^0<T_C^0`$, the AFM spins exhibit a rotation of the sublattice magnetization. With increasing temperature they turn into the direction of the FM film and become collinear above the sublattice reorientation temperature $`T_R`$, cf. Figure 10. The AFM magnetic arrangement for $`T>T_R`$ represents a ‘layered AFM structure’ consisting of ferromagnetic layers with an alternating orientation for neigboring layers. All AFM layers become collinear at the same temperature, the variation of $`\varphi _{\mathrm{AFMi}}^0(T)`$ is the steeper the larger the distance of layer $`i`$ from the interface. A similar behavior is also obtained for FM films thicker than a monolayer. In addition, for $`T_N^0>T_C^0`$ the behaviors of the FM and AFM subsystems are interchanged. Thus, the mentioned symmetry between FM and AFM films for the sc(001) symmetry, as calculated within MFA, is also present for thicker films.
The discussion of the corresponding behavior of fcc(001) FM-AFM films requires some introductory remarks. Unlike FM films, and unlike sc(001) AFM films, as calculated within MFA the Néel temperature $`T_N(n_{\mathrm{AFM}})`$ of an fcc(001) AFM film with an in-plane AFM order and with nearest neighbor exchange interactions only does not increase with increasing thickness $`n_{\mathrm{AFM}}`$. Merely, a constant $`T_N(n_{\mathrm{AFM}})`$ given by the one of the monolayer ($`n_{\mathrm{AFM}}=1`$) results. Consequently, the same magnetizations $`M_i(T)`$, independent of the individual layer $`i`$, are obtained. Hence, the expression for the ordering temperature, cf. equation (8), is not valid for an fcc(001) AFM bilayer. The reason is that for such a system with a collinear magnetization each layer is virtually decoupled from the rest. Only in the case of noncollinear magnetic order, as is present e.g. close to the FM-AFM interface, a net coupling between neighboring AFM layers results.
Keeping these features in mind we now discuss the finite-temperature properties of an FM-AFM system with an fcc(001) symmetry and for $`n_{\mathrm{AFM}}>1`$. As for the bilayer, all FM spins remain strictly collinear for all temperatures. In Figure 11 the magnetizations $`M_i(T)`$ and angles $`\varphi _i^0(T)`$ for FM-AFM films with $`n_{\mathrm{AFM}}=2`$ and $`n_{\mathrm{AFM}}=3`$ close to their critical temperatures $`T_C`$ is presented. The FM film thickness $`n_{\mathrm{FM}}=1`$ and the coupling constants are the same for both cases and are chosen in such a way that $`T_N^0>T_C^0`$. The case $`n_{\mathrm{AFM}}=2`$ corresponds to the situation shown in Figure 8. The FM layer becomes paramagnetic above the critical temperature $`T_C^{}`$ in the presence of a still ordered AFM film. Thus, two critical temperatures can occur also for thicker FM-AFM films. The magnetizations $`M_i(T)`$ of the two AFM layers are identical to each other over the whole temperature range and vanish at $`T_N^0(n_{\mathrm{AFM}}=1)`$, cf. Figure 11a. In contrast, the equilibrium angles $`\varphi _{\mathrm{AFMi}}^0(T)`$ are different for both AFM layers, and approach $`\varphi _{\mathrm{AFMi}}^0(T)0`$ for $`TT_C^{}`$ (not shown).
For the applied coupling constants this behavior changes drastically if three AFM layers are considered. Although still $`T_C^0<T_N^0`$, the FM layer now dominates and causes a similar behavior as shown in Figure 7 for an FM-AFM bilayer with a strong interlayer coupling. As can be seen from Figure 11c, the angles $`\varphi _{\mathrm{AFMi}}^0(T)`$ of the AFM layers increase with increasing temperature. The AFM spins eventually become collinear with respect to the FM, with an alternating magnetic orientation for neighboring AFM layers. In contrast to the sc(001) film shown in Figure 10, the sublattice reorientation temperature $`T_{R,i}`$ is now layer dependent and increases as the distance of the layer $`i`$ from the interface becomes larger. Moreover, as long as the AFM layers maintain a noncollinear structure, the magnetizations $`M_i(T)`$ are identical and independent of the layer index. Only for temperatures $`T>T_{R,i}`$ the $`M_i(T)`$ differ from each other, and vanish together with the magnetization of the FM film at the common ordering temperature $`T_C`$ of the total FM-AFM system, cf. Figure 11b.
The different behavior of the AFM magnetizations in coupled sc(001) and fcc(001) FM-AFM films can be understood as follows. For the former symmetry the magnetic structures of all FM and AFM layers are disturbed for $`T=0`$, and become collinear at the same temperature. On the other hand, for fcc(001) films the FM layers always remain collinear ($`\varphi _{\mathrm{FMi}}^0(T)=0`$). Consider the situation depicted in Figures 11b,c. If, e.g., the spins of the AFM interface layer (AFM1) turn into the direction of the FM, the remaining AFM layers virtually experience an ordered FM film with an increased thickness. As before, the remaining AFM layers can maintain a noncollinear magnetic arrangement, and there is no need for all layers to become simultaneously collinear. In addition, we note that the AFM films above the sublattice reorientation temperatures exhibit, for both AFM thicknesses shown in Figure 11, a collinear structure. Nevertheless they behave differently since for $`n_{\mathrm{AFM}}=2`$, Figure 11a, the magnetic structure refers to an ‘in-plane AFM’ for $`T>T_C^{}`$, and for $`n_{\mathrm{AFM}}=3`$ and $`T>T_{R,i}`$, Figure 11b,c, to a ‘layered AFM structure’. In the latter case the AFM magnetizations $`M_i(T)`$ are layer dependent, and the corresponding ordering temperature depends on the AFM film thickness.
## 6 Conclusion
In this study we investigated how the magnetic structure rearranges in the vicinity of the interface between a ferromagnet and an antiferromagnet. Thin film systems with sc(001) and fcc(001) symmetries have been solved for both zero and finite temperatures within the framework of a mean field approximation. A variety of configurations was obtained, and the underlying physics has been discussed. In contrast with previous work Koo97 ; HiM86 ; CrC01 ; JeD02 , these properties are mainly determined by the isotropic exchange interactions. The consideration of a particularly simple bilayer system, and the application of an MFA at finite temperatures, allows us to derive analytical expressions for various quantities. These serve as estimates of the magnetic behavior for more realistic coupled FM-AFM systems having thicker FM and AFM films.
We emphasize the different behavior of the sc(001) and fcc(001) lattice symmetries. In particular, a canting of the sublattice magnetizations of both FM and AFM layers is obtained for the former case, whereas for the latter only the AFM layer is disturbed. Moreover, if the bare Curie temperature $`T_C^0`$ of the FM film is larger than the bare Néel temperature $`T_N^0`$ of the AFM film, the AFM spins become collinear with respect to the FM system above the sublattice reorientation temperature $`T_R`$ for both investigated symmetries. For $`T>T_R`$ the AFM film assumes a ‘layered AFM structure’. For an sc(001) lattice this reorientation happens simultaneously for all layers at the same temperature. For $`T_N^0>T_C^0`$ a corresponding behavior with an interchanged role of the FM and AFM films results, which within MFA is perfectly symmetric to the case $`T_N^0<T_C^0`$. In contrast, such a symmetry between FM and AFM is not present for the fcc(001) lattice. Merely, the FM spins always remain strictly collinear. The different AFM layers turn into the direction of the FM at different sublattice reorientation temperatures.
Moreover, the possibility of two critical temperatures is pointed out, as derived for fcc(001) FM-AFM bilayers for $`T_N^0>T_C^0`$ and $`J_{\mathrm{int}}<J_{\mathrm{int}}^{}`$. In this case the FM film becomes paramagnetic at temperatures $`T_C^{}`$ where the AFM film is still magnetically ordered. The presence of two different $`T_C`$’s in magnetic systems is well known, for instance, for two coupled semi-infinite ferromagnets. Similarly, if a magnetic film with a strong exchange interaction is deposited on a bulk ferromagnet, two different ordering temperatures may exist KTS83 . In contrast, to our knowledge two critical temperatures for coupled magnetic films with finite thicknesses have not been reported previously.
However, the existence of two $`T_C`$’s is expected to be fragile. In fact, small deviations from the fcc(001) symmetry, for example in presence of disorder near the interface, could destroy the lower one. The reason is that in this case the couplings of the two AFM sublattices across the interface do not cancel exactly, and a magnetization will be induced in the FM for $`T>T_C^{}`$. In general, we note that in real FM-AFM interfaces disorder is always present, like step, vacancies, interdiffusion, etc. In this case the lateral periodicity of the magnetic structure as sketched in Figures 1b and 3b will vanish with increasing degree of disorder. The presented results are obtained for fully ordered interfaces and thus will serve as starting points to investigate the role of disorder at FM-AFM interfaces. For example, the resulting magnetic arrangement can be a mixture of the two extremal cases represented by the sc(001) and fcc(001) stackings. For a strong disorder compensated and noncompensated interfaces can no longer be distinguished Fin04 . Moreover, as mentioned in the Introduction, the consideration of disorder seems to be essential to explain the exchange bias effect MNL98 .
As noted in Section 2, we have chosen anisotropy easy axes of the FM and AFM films which support a perdendicular magnetic arrangement of both subsystems. Anisotropies with different symmetries and arbitrary directions of the easy axes can be considered as well. In that case the magnetic structure and the (sublattice) spin rotation will also depend on the anisotropies. In presence of disorder the anisotropy easy axes will be site-dependent which, if the anisotropy is sufficiently strong, can disturb the lateral periodicity of the magnetic arrangements sketched in Figures 1b and 3b.
In this connection we like to point out an important difference with our prior work JeD02 , which also dealt with coupled FM-AFM films. There the magnetization of the FM undergoes a full spin reorientation transition (SRT) as a function of temperature, i.e., the net magnetization of each layer changes its direction whereas its magnitude stays approximately constant. To exhibit such an SRT a significant anisotropy in the FM must be present, eventually competing with the interlayer exchange. In contrast, in the present study both sublattices in every layer exhibit a magnetic reorientation, with opposite sense of the rotations. The directions of the net layer magnetizations remain constant and do not show an SRT, whereas their magnitudes vary considerably. These differences should become apparent in possible experimental realizations, e.g., within an element-specific X-ray magnetic linear or circular dichroism (XMLD, XMCD) measurement SLA04 . Whether a full SRT like in JeD02 , or whether the magnetic arrangement as described in the present study dominates, depends on the actual FM-AFM system under consideration.
Finally, we like to discuss the influence of collective magnetic excitations (spin waves). As is well known, for 2D magnetic systems these excitations play a very important role, which however are neglected in the MFA used in this study. It is therefore important to apply improved methods which take into account collective excitations, for example, within a many-body Green’s function theory (GFT) Tya67 . FM-AFM bilayer and multilayer systems have been investigated previously by this method, considering a collinear magnetization Die89 . In SuS98 the collective excitations were discussed to be a possible source for the exchange bias effect. The GFT has recently been generalized JKW03 to take into account several nonvanishing components of the magnetization, hence allowing the investigation of noncollinear magnetic strucures. To avoid the catastrophe of the Mermin-Wagner-theorem MeW66 magnetic anisotropies must be incorporated explicitly. Analytical results, which can be drawn from the much simpler MFA, may not be obtained from such improved theoretical approaches. Also, it has been shown that MFA yields at least qualitatively correct results for anisotropic magnetic thin films, although quantitatively it strongly overestimates the ordering temperatures. Preliminary results calculated with GFT show that the main properties obtained in the present study are supported. In particular, this is valid for the sublattice magnetic reorientation, and the distortion of both FM and AFM layers in case of sc(001) FM-AFM films. On the other hand, the exact symmetry between the FM and AFM layers for the sc(001) system turns out to be an artifact of the MFA. The reason is that the spin wave dispersion relations for an FM and an AFM differ qualitatively, as do the respective ordering temperatures even for the same strengths of the exchange couplings. However, MFA incorrectly yields the same ordering temperatures. Further investigations using GFT are underway.
PJJ acknowledges support by the Deutsche Forschungsgemeinschaft through SFB 290, TP A1. HD and MK acknowledge support by ECOS. MK was also supported by Fundación Andes and by FONDECYT, Chile, under grant No. 1030957. |
warning/0507/quant-ph0507217.html | ar5iv | text | # On quantum control limited by quantum decoherence
## I Introduction
Generally people can utilize an external field to manipulate the time evolution of a quantum system from an arbitrary initial state to reach any wanted target state. If the external field is classical and can be artificially controlled to be time-dependent, then we refer this kind of manipulation as a classical control q-control . In quantum computations q-inform , the quantum logic gate operations can be regarded as classical controls in most cases where the controller is essentially classical and the control can be turned on or off classically at certain instants.
In this paper we consider the quantum control, in which the controller is quantized and obeys the laws of quantum mechanics. It is shown that the back action of the controlled system should be considered, which may have a negative side-effect on the controllability. There are two motivations for our investigations.
Firstly, it is exciting to explore the finiteness of human being’s abilities to control the nature, and a “down-to-earth” starting point for this exploration in physics should be a concrete model even though it is oversimplified. With some reasonable models one could demonstrate how the fundamental laws of physics impose the limits on the controllability in principle. These refer to some basic issues in physics, such as the energy bound, the basic precision of measurement (or standard quantum limit (SQL)SQL ). It is emphasized that the quantum decoherence may result from *the control itself* when the controller is essentially considered as a quantum subsystem.
Secondly, though the physical implementation of quantum computation seems to be difficult, the huge power of quantum computation has been demonstrated by some quantum algorithms in principle. The limit of quantum control can bring a physical limit to quantum computation architecture since it is based on complete quantum blocks including the controller. Lloyd discussed how the physical constants impose a limit on the power and the memory in the quantum computer Lloyd2000 , while Ozawa and Gea-Banacloche Ozawa2002 ; Julio2002b considered the conservation law and the minimum energy requirement for quantum computation respectively. Our present study can also be regarded as a part of the growing body of the explorations in this direction.
In Sec. II, we start with a model with a single mode field as a controller and a two-level system (qubit) as the controlled system. We found that it is possible to implement some phase gate controls without inducing decoherence to the controlled system. However, the single mode example is far from practical cases, and thus we further study the quantum control in a more general case in Sec. III. In Sec. IV the control induced decoherence is explained as a phase uncertainty by associating it with the SQL. In Sec. V the obtained results is highlighted as the complementarity of the controllability and the control induced decoherence. An inequality similar to the Heisenberg uncertainty relation is presented as the accurate bound of quantum gates under the quantum control.
## II An exactly soluble model for the quantum control
To have a clear picture about the quantum control, let us first start with a simple model. The total system that we concern is closed, which consists of the controller $`C`$ with the Hamiltonian $`H_c`$ and the controlled system $`Q`$ with the Hamiltonian $`H_q`$. The system is in the initial states $`|\psi _c(0,R)`$ and $`|\psi _q(0)=_nc_n|n`$ respectively, where $`R`$ represents the controlling parameters. For a given target state $`|\psi _t`$ of $`Q`$, the quantum control is described as a factorized evolution
$$|\psi _q(0)|\psi _c(0,R)|\psi _q(T)|\psi _c(T)$$
(1)
of the total system driven by the interaction Hamiltonian $`H_{qc}`$ within the time duration $`(0,T)`$. If one could choose an appropriate initial state and the corresponding parameters $`R`$ such that the partial wave function $`|\psi _q(T)U_q(T)|\psi _q(0)`$ is just the target one $`|\psi _t`$, where a global phase difference is allowed, then we could say that an ideal quantum control is realized. Usually $`U_q(T)`$ defines a quantum logic operation in quantum computation.
We now consider an exactly soluble example, where the controlled system is a qubit with two basis states $`|0`$ and $`|1`$ and the controller is a single mode boson field with free Hamiltonian $`H_c=\mathrm{}\omega a^{}a`$, here $`a^{}`$ ($`a`$) is the creation (annihilation) operator. The interaction
$$H_{qc}=|11|V11|(ga+g^{}a^{})$$
(2)
between them is of non-demolitionSQL , i.e., $`[H_{qc},H_c]0`$ and $`[H_{qc},H_q]=0`$. Since $`H_q`$ is conserved during the evolution we take $`H_q=0`$ without loss of the generality. In the interaction picture the time-dependent potential
$$V_I(t)=gae^{i\omega t}+h.c$$
(3)
acts only on the state $`|1`$, but not on $`|0`$. This Hamiltonian originates from the atom-field system in the large detuning limit, but the problem is greatly simplified for convenience q-opt .
Now we explore the possibility of automatically creating a phase gate operation
$$|\psi _q(0)=c_0|0+c_1|1|\psi _q(t)=c_0|0+c_1e^{i\varphi }|1$$
(4)
driven by $`H_{qc}`$. Essentially, the phase gate operation is supposed to generate a relative phase $`\varphi `$ between $`|0`$ and $`|1`$ and the total system experiences a factorized evolution
$$(c_0|0+c_1|1)|\psi _c(0)(c_0|0+c_1e^{i\varphi }|1)|\psi _c(T).$$
(5)
We will show that, only a class of phase gates with special phases depending on the global parameters, such as the coupling coefficients $`g`$ and the gate operation time $`T`$, can be implemented precisely, while the other phase gates definitely result in a decoherence in the qubit system, and can only be implemented in an inaccurate way.
Obviously the Hamiltonian $`H=H_{qc}+H_c`$ describes a typical conditional dynamics sun95 . Let the total system be initially in a superposition of
$$|\mathrm{\Psi }(0)=(c_0|0+c_1|1)|\alpha ,$$
(6)
where the boson field is in a coherent state $`|\alpha `$. The total system will evolve into an entangled state
$$|\mathrm{\Psi }(t)=c_0|0|\alpha +c_1|1e^{i\stackrel{}{\mathrm{\Phi }}}|\alpha ,$$
(7)
where
$$e^{i\stackrel{}{\mathrm{\Phi }}}=\widehat{T}\mathrm{exp}(i_0^tV_I(t^{})𝑑t^{})$$
(8)
is a time ordered integral. A formal phase operator can be explicitly calculated as
$$\stackrel{}{\mathrm{\Phi }}(t)=\eta (t)a+h.c+\phi (t)+i\xi (t),$$
(9)
where time-dependent coefficients
$`\eta (t)`$ $`=`$ $`i{\displaystyle \frac{g}{\omega }}(1e^{i\omega t}),`$
$`\phi (t)`$ $`=`$ $`{\displaystyle \frac{|g|^2}{\omega ^2}}(\omega t\mathrm{sin}\omega t),`$ (10)
$`\xi (t)`$ $`=`$ $`{\displaystyle \frac{|g|^2}{\omega ^2}}(1\mathrm{cos}\omega t),`$
are obtained through the Wei-Norman algebraic technique sun91 . Then we can write down the total wave function as an entangled state
$$|\mathrm{\Psi }(t)=c_0|0|\alpha +e^{i\phi (t)\xi (t)}c_1|1|\alpha +\eta (t).$$
(11)
Obviously, at the special instants
$$t=T=2k\pi /\omega ,$$
(12)
where $`k𝖹`$, both the decay factor $`\xi (t)`$ and the displacement $`\eta (t)`$ in the coherent state $`|\alpha +\eta (t)`$ vanish. And a real phase
$$\phi (T)=\varphi _s=\frac{|g|^2}{\omega }T.$$
(13)
occurs in the above entanglement state. Thus we realize a phase gate operation Eq.(4) of certain phase $`\varphi _s`$, which is induced by the factorized evolution
$$|\mathrm{\Psi }(0)|\mathrm{\Psi }(T)=(c_0|0+c_1e^{i\varphi _s}|1)|\alpha .$$
(14)
It defines a reduced density matrix of a pure state
$$\rho _q=|c_0|^2|00|+|c_1|^2|11|+c_1c_0^{}e^{i\varphi _s}|10|+h.c$$
(15)
for the qubit system.
If the evolution time is not just at the instant $`t=T`$, the reduced density matrix
$$\rho _r=|c_0|^2|00|+|c_1|^2|11|+c_1c_0^{}D(t)|10|+h.c$$
(16)
is not of a pure state duo to the decoherence factor
$$D(t)=\alpha |W(t)|\alpha =e^{i\varphi (t)\xi (t)}$$
(17)
where
$$\varphi (t)=2\mathrm{𝐈𝐦}[\frac{(1e^{i\omega t})g\alpha }{\omega }]+\frac{(\omega t\mathrm{sin}\omega t)|g|^2}{\omega ^2},$$
(18)
and
$$\xi (t)=\frac{|g|^2}{\omega ^2}(1\mathrm{cos}\omega t).$$
The difference between $`\rho _q`$ and $`\rho _r`$ can be characterized by the control fidelity $`F(t)=Tr(\rho _q\rho _r)`$, which is defined as the overlap of the target state $`\rho _q`$ and the final state $`\rho _r`$. By a straightforward calculation, we have
$`F(t)`$ $`=`$ $`12|c_0|^2|c_1|^2[1\mathrm{𝐑𝐞}(D(t)e^{i\varphi _s})]`$
$`=`$ $`12|c_0|^2|c_1|^2[1e^{\xi (t)}\mathrm{cos}(\varphi (t)\varphi _s)].`$
In Fig.1 we plot the curve $`F(t)`$, where $`g=0.1`$, $`\omega =1`$ and $`\alpha =1.5`$. For convenience, we have take $`|c_0|^2=|c_1|^2=1/2`$ and then
$$F(t)=1\frac{1}{2}[1e^{\xi (t)}\mathrm{cos}(\varphi (t)\varphi _s)].$$
(19)
It can be seen that $`F(t)`$ is a periodic function with unity as the maximum value. As a functional, the period $`T=f[F]`$ is a function of the function $`F(t)`$. $`f[F]`$ is determined by the system parameters $`g`$ and $`\omega `$. When $`t=f[F(t)]`$, the control fidelity takes its maximum $`F(t)=1`$ and then we realized an ideal phase gate operation with the phase $`\varphi _s=|g|^2T/\omega `$.
In order to realize a real control we require that the effective interaction $`V_I(t)`$ could be automatically switched on and off at time $`0`$ and $`T`$, i.e., the controllable condition ($`CABC`$)
$$V_I(0)=V_I(T)=0,$$
(20)
is satisfied for $`V_I(t)=\psi _c(t)|V_I(t)|\psi _c(t)`$. For the above example, this requirement means
$`\mathrm{𝐑𝐞}(g\alpha )`$ $`=`$ $`0,`$ (21)
$`\mathrm{𝐈𝐦}[g\alpha ]\mathrm{sin}\omega T`$ $`=`$ $`\omega \xi (T),`$ (22)
for $`\mathrm{sin}\omega T0`$. When there is no loss of qubit coherence at the instance $`t=T`$ ($`\xi (T)=0`$), the requirement Eq.(21) and Eq.(22) for an ideal quantum control is just $`g\alpha =0`$. It is absurd and impracticable. However, there exist the situations ($`\mathrm{sin}\omega T=0`$) satisfying the requirement for quantum control: $`\mathrm{𝐑𝐞}[g\alpha ]=0`$ and $`\omega T=2k\pi `$, $`k𝖹`$, which is reasonable in principle since a pure imaginary number $`g\alpha =i|g\alpha |`$ does not vanish even though it have a vanishing real part. Therefore some target states are obtained as the superpositions state of $`|0`$ and $`|1`$ with specific relative phases that can be implemented perfectly by the quantum control.
However, the above phase gate control could only generate particular phases $`\varphi _s`$ on the qubit state $`|1`$, which is completely determined by the coupling factor $`g`$ and the controller field frequency $`\omega `$. In this sense we can not achieve a quantum control of implementing universal phase gates for a given total system with fixed $`g`$ and the controller field frequency $`\omega `$. To overcome this problem the local parameters of the initial states of the controller should be used in the quantum control rather than the fixed global parameters of the total system. We will explore this possibility in Sec. V where the quantum decoherence will be considered based on the uncertainty relation that relates to a multi-mode coherent field.
## III Quantum control by general controller
Staring with an idealized model, the above investigations provide us some insights into the quantum control problem. In order to consider the more practical cases, we will analyze the quantum controllability in this section. To focus on the central idea we do not consider the influence of the environment yet. The entire system that we consider is an isolated system including the controller $`C`$ with the Hamiltonian $`H_c`$ and the controlled system $`Q`$ with the Hamiltonian $`H_q`$ . To bring out more clearly the physical picture of such a quantum control, the minimal assumption is that the Hamiltonian includes only two items: $`H_{qc}`$ and $`H_c`$. Matching this assumption, there exists a practical case that the non-demolition control satisfies $`[H_q,H_{qc}]=0`$ and then the free evolution of the controlled system is eliminated.
Conveniently we work in the interaction picture with the Schr$`\ddot{o}`$dinger equation
$$i\mathrm{}\frac{d}{dt}|\mathrm{\Psi }^I(t)=H_{qc}^I(t)|\mathrm{\Psi }^I(t).$$
(23)
Formally, the quantum control requires that the interaction Hamiltonian
$$H_{qc}^I(t)=e^{iH_ct/\mathrm{}}H_{qc}e^{iH_ct/\mathrm{}}$$
(24)
can be automatically turned on and off at certain instants $`t=0`$ and $`t=T`$ during the evolution of the controller system. Under the quantum control a quantum gate operation is accomplished by the controlled system. Besides, it is also required that the controlling parameters depend on the initial state of the controller system. By applying them to quantum computing, the quantum computer implements the operations programmed by the controller.
Without loss of the generality, we still take the controlled system as a qubit with two basis states $`|0`$ and $`|1`$. An ideal quantum control with $`U_q(T)`$ exerting on the qubit can be described as a factorized evolution
$$U_I(T)=e^{\frac{i}{\mathrm{}}(H_{qc}+H_c)T}=U_q(T)U_c(T)$$
(25)
of the total system. So that a controlled evolution of the qubit system is implemented as $`|\psi _q(T)=U_q(T)|\psi _q(0)`$, while $`|\psi _c(T)=U_c(T)|\psi _c(0)`$ defines the final state of the controller. Here, $`|\psi _q(0)=c_0|0+c_1|1`$ and $`|\psi _c(0)`$ are the initial states of the qubit and the controller respectively. We note that, because the Hermitian operators $`H_{qc}`$ and $`H_c`$ do not commute with each other, thus there is not simply $`\mathrm{exp}(iH_{qc}T)=U_q(T)`$ in practice. We emphasize that, due to the limitation resulting from the Heisenberg uncertainty principle, the realistic control can not be carried out in such a perfect way as a completely factorized evolution.
Generally, the Hermitian operators $`H_{qc}`$ and $`H_c`$ do not commute with each other and there exists an uncertainty relation:
$$\mathrm{\Delta }H_{qc}^I\mathrm{\Delta }H_c\frac{1}{2}|[H_{qc}^I,H_c]|,$$
(26)
where the variations $`\mathrm{\Delta }A=\sqrt{A^2A^2}`$, $`A=H_{qc}`$ and $`H_c`$. In a consistent approach for quantum measurement zls , this uncertainty relation is also responsible for the decoherences induced by the detector as well as those induced by the quantum control. Roughly speaking, the variation $`\mathrm{\Delta }H_{qc}^I(t)`$ is relavent to the induced decoherence in the qubit system, while the term $`|[H_{qc}^I(t),H_c]|`$ indicates the influence of quantum control, and $`\mathrm{\Delta }H_c`$ is associated with the power or the average energy of the controller. The conservation laws throw some limits on such implementation of quantum gates Ozawa2002 . For example, a quantum control to complete a CNOT gate usually concerns the transfer of some conservation quantities between qubits. To focus on the problems in the following, we will only consider the quantum control itself, which does not involve the transfer of any known conservation quantities.
Now we assume a non-demolition controlling interaction $`H_{qc}=|11|V`$ with a potential $`V`$ that acts only on the qubit state $`|1`$. It does not play any role at the beginning and the end of the gate operation, but we require that it is generated by the controller, and a nontrivial phase is left on the qubit state $`|1`$. Actually, as for the quantum controls in quantum information processing, it is expected that a quantum computer could work like electronic computers: when the programs are designed then stored in it initially, the quantum computer should be able to carry out computations without any other assistants. The basic requirement for the quantum control is that the interaction can be switched on and off automatically at certain instants, e.g., at $`t=0`$ and $`t=T`$,
$`V_I(0)`$ $`=`$ $`\psi _c(0)|V_I(0)|\psi _c(0)=0,`$
$`V_I(T)`$ $`=`$ $`\psi _c(T)|V_I(T)|\psi _c(T)=0,`$ (27)
where $`V_I(t)\mathrm{exp}(iH_ct/\mathrm{})V(iH_ct/\mathrm{})`$. The above Eq.(27)is the general controllable condition. The sandwich $`V`$ is defined as the average of the operator $`V`$ over the controller state. This means that the effective interaction is obtained by taking the average of $`V_I(t)`$ over the instantaneous controller states $`|\psi _c(t)`$.
Generally, the controller in physical implementations of the quantum control are various fields that are supposed to be classical. For example, the microwave electromagnetic fields are used to manipulate the nuclear spin-qubits in NMR, the laser fields are applied to control the atomic qubits and the classical magnetic flux and voltage are utilized to adjust the Josephson-Junction based qubits. However, the controlling fields are essentially of quantization and are usually described by coherent states or some quantum mechanical mixture.
Starting from the initial state where the qubit is in $`|\psi _q=c_0|0+c_1|1`$, the total system evolves according to the entangled state
$$|\mathrm{\Psi }(t)=c_0|0|\psi _c(0)+c_1|1e^{i𝚽(t)}|\psi _c(0),$$
(28)
where we have defined the time-order integral
$$e^{i𝚽}=\widehat{T}\mathrm{exp}(\frac{i}{\mathrm{}}_0^tV_I(\tau )𝑑\tau ).$$
(29)
The decoherence factor sun93 is an expectation of the unitary operator
$$D(T)=\psi _c(0)|e^{i𝚽}|\psi _c(0),$$
(30)
which can be used to characterize the quantum controllability.
Now we need to consider that in what cases the above entangled state $`|\mathrm{\Psi }(t)`$ can become a factorized state Eq.(5) at ceratain instant $`t=`$ $`T`$ so that the ideal quantum control is realized by choosing the initial state $`|\psi _c(0)`$ of the controlling system. A simplest illustration is that $`V_I(t)=V`$ is a static potential and thus
$$e^{i𝚽}=\mathrm{exp}(iTV/\mathrm{}).$$
(31)
If we choose $`|\psi _c(0)=|\varphi `$ with the eigenvalue $`\varphi `$, then $`\mathrm{exp}(i𝚽)`$ becomes a $`cnumber`$ phase factor $`\phi `$, and the time evolution automatically generate a phase gate operation with the $`cnumber`$ phase:
$$|\mathrm{\Psi }(T)=(c_0|0+c_1e^{i\phi }|1)|\psi _c(0).$$
(32)
Indeed, the phase $`\varphi `$ multiplied to the qubit state $`|1`$ is well defined and can be generated with arbitrary precision at a suitable instant $`T`$ by choosing the initial state $`|\psi _c(0)=|\varphi `$ of the controller. This is what we want: the qubit system is controlled by the parameters of the initial state as well as the evolution time. It seems that no fundamental restrictions exists for $`|\psi _c(0)`$ and $`T`$.
However, the above idealized situation is far from the realistic cases in practical quantum controls. Firstly, the precision of the quantum control is guaranteed by the stability of potential, i.e., $`[V,H_c]=0=[H_{qc},H_c]`$. However, this means that the free Hamiltonian evolution of controller has no influence on the effective interaction by $`H_{qc}`$ and thus the $`CABC`$ can not be satisfied automatically. Therefore, we infer that, in order to realize a quantum control with the “switched on and off”, the potential $`V_I(t)`$ could not be a static one. In this case the $`cnumber`$ phase is not well defined by the initial state of the controller and thus there exists a phase fluctuation $`\mathrm{\Delta }\mathrm{\Phi }`$ in the implementation of the quantum control.
To explore the possibility of assorting with the $`CABC`$ and the precision of the quantum control, we distinguish two cases by whether the potential $`V_I(t)`$ generated by the controller is commutative or not at different instants, i.e.,
$`Case.1`$ $`:`$ $`[V_I(t),V_I(t^{})]=0,`$ (33)
$`Case.2`$ $`:`$ $`[V_I(t),V_I(t^{})]0.`$ (34)
In the first case a phase factor operator can be simply defined as
$$𝚽=\frac{1}{\mathrm{}}_0^TV_I(\tau )𝑑\tau .$$
(35)
Under the small variation $`\mathrm{\Delta }𝚽1`$, the decoherence factor can be calculated as
$$D(T)e^{i𝚽\frac{1}{2}(\mathrm{\Delta }𝚽)^2}e^{i𝚽}d(T).$$
(36)
Similar to the arguments about the exactly solvable model in Sec.II, an observation is that the ideal quantum control can be characterized by whether or not the decoherence factor $`|D(T)|=|\mathrm{exp}(i𝚽)|`$ can reach unity. Actually the phase multiplied to the qubit state $`|1`$ is the real part of the expectation value of the phase factor operator $`𝚽`$ plus a decay factor from its quantum fluctuation $`(\mathrm{\Delta }𝚽)^2`$Aha .Thus the quantum controllability is destroyed by the phase fluctuation $`(\mathrm{\Delta }𝚽)^2`$ in general.
In the following, we will show that the phase fluctuation $`(\mathrm{\Delta }𝚽)^2`$ will result in a loss of quantum coherence or quantum dephasing. To this end we calculate
$$(\mathrm{\Delta }𝚽)^2\frac{1}{\mathrm{}^2}_0^T𝑑t_0^t\mathrm{\Delta }V_I(t)\mathrm{\Delta }V_I(\tau )𝑑\tau ,$$
(37)
which shows that phase fluctuation $`(\mathrm{\Delta }𝚽)^2`$ is just the correlated fluctuation of Heisenberg interaction. Thus $`d(T)=\mathrm{exp}[(\mathrm{\Delta }𝚽)^2/2]`$ is a decaying factor in $`D(T)`$ accompanying the off-diagonal terms of the reduced density matrix of the qubit system. To quantitatively describe that to what extent the target sate
$$|\psi _t=c_0|0+c_1e^{i𝚽}|1$$
can be reached by the controlled time evolution $`|\mathrm{\Psi }(t)`$, the control fidelity
$`F(t)`$ $`=`$ $`Tr[|\mathrm{\Psi }(t)\mathrm{\Psi }(t)|(1|\psi _t\psi _t|)]`$ (38)
$`=`$ $`Tr_c(\psi _t|\mathrm{\Psi }(t)\mathrm{\Psi }(t)|\psi _t)=Tr(\rho _t\rho _r)`$
is defined in terms of the reduced density matrix $`\rho _t`$ and the reduced density matrix $`\rho _r=Tr_c[|\mathrm{\Psi }(t)\mathrm{\Psi }(t)|]`$, where $`Tr_c`$ indicates tracing over the variables of the controller. In this case the result is obtained as
$$F(t)=12|c_0|^2|c_1|^2(1e^{\frac{1}{2}(\mathrm{\Delta }𝚽)^2}).$$
Thus the corresponding error measure
$$\epsilon =1F(t)=2|c_0|^2|c_1|^2[1d(t)]$$
(39)
describes the failure probability of the quantum control.
For the second case, due to the non-vanishing commutator between $`V_I(t)`$ at different instants we can not generally define a phase factor operator $`𝚽`$, but we can still formally write $`D(T)=\mathrm{exp}(i𝚽)`$ or
$$D(T)=e^{i\mathrm{\Phi }\xi }\mathrm{exp}(i𝚽\frac{1}{2}(\mathrm{\Delta }𝚽)^2).$$
(40)
This can give all similar results as the case 1. The exactly solvable model in Sec. II belongs to the second case. This result is exact for the above example presented in the last section where
$$\frac{1}{2}(\mathrm{\Delta }\mathrm{\Phi })^2=\xi (t),𝚽=\varphi (t).$$
(41)
As discussed in the above, the decoherence induced limit to the quantum control has been explained based on the phase uncertainty. In fact, this understanding reveals once again the inherence of the quantum decoherence in the generalized two-slit experiment about $`|0`$ and $`|1`$, whose interference fringe vanishes when one determines which slit the particle comes from. According to Heisenberg, this is due to the randomness of relative phases heisenberg from the quantum control. Furthermore, we can conclude from the above exact solution that the large random phase change just originates from Heisenberg’s position-momentum uncertainty relation $`\mathrm{\Delta }x_k\mathrm{\Delta }p_k=1/2`$. This observation will help us to discover a bound on the quantum control.
## IV Phase uncertainty due to standard quantum limit
Based on our previous explorations on the relation between the two explanations for quantum decoherence zls , using the position-momentum uncertainty relation, we now can associate the physical limit of quantum control with the standard quantum limit (SQL) in quantum measurement context SQL through a concrete example as follows.
This is a more practical example that the qubit is controlled by a multi-mode electromagnetic field
$$E=\underset{k}{}(u_k(x)a_ke^{i\omega _kt}+h.c)$$
(42)
with the mode functions $`u_k(x)`$. The controlling Hamiltonian $`H_{qc}(t)=|11|V_I(t)`$ in the interaction picture reads as
$`H_{qc}(t)`$ $`=`$ $`|11|{\displaystyle \underset{k}{}}H_k`$ (43)
$``$ $`|11|{\displaystyle \underset{k}{}}\mathrm{}(g_ka_ke^{i\omega _kt}+h.c),`$
where $`\omega _k`$ are the mode frequencies, $`a_k`$ and $`a_k^{}`$ the creation and annihilation operators respectively, and $`g_k`$ the mode couplings constants between the qubit and the field modes. We suppose that the electromagnetic field is initially prepared in a multi-mode coherent state
$$|\psi _c(0)=|\alpha \underset{k}{}|\alpha _k$$
(44)
as a direct product of the coherent state $`|\alpha _k`$ of $`k`$th mode. In such a initial state, the observable is the the average of the field operator
$$\alpha |E|\alpha =\underset{k}{}[u_k(x)\alpha _ke^{i\omega _kt}+h.c],$$
(45)
which is a wave packet, the superposition of many plane waves. This means that, to realize a more realistic quantum control, we need a wave packet rather than a single mode or a plane wave.
The free Hamiltonian of the qubit system has been omitted without loss of generality. The potential $`V_I(t)`$ exerts on the qubit state $`|1,`$ but not on the qubit state $`|0`$. Then the evolution can be obtained as
$$U(t)=|00|\mathrm{𝟏}+|11|e^{i𝚽},$$
where
$$e^{i𝚽}=\underset{k}{}e^{i\mathrm{\Phi }_k}=\underset{k}{}U_k\underset{k}{}\widehat{T}\mathrm{exp}(\frac{i}{\mathrm{}}_0^tH_k(\tau )𝑑\tau ).$$
(46)
We can explicitly calculate the phase operator $`\mathrm{\Phi }=_k\mathrm{\Phi }_k`$ defined above by the method similarly to that used for the example about the single mode field in Sec. II. It is obtained by
$$\mathrm{\Phi }_k=\eta _k(t)a_k+h.c+\phi _k(t)+i\xi _k(t),$$
with three time-dependent parameters
$`\eta _k(t)`$ $`=`$ $`i{\displaystyle \frac{g_k}{\omega _k}}(1e^{i\omega _kt}),`$
$`\phi _k(t)`$ $`=`$ $`{\displaystyle \frac{|g_k|^2}{\omega _k^2}}(\omega _kt\mathrm{sin}\omega _kt),`$ (47)
$`\xi _k(t)`$ $`=`$ $`{\displaystyle \frac{|g_k|^2}{\omega _k^2}}(1\mathrm{cos}\omega _kt).`$
The phase operator can be re-written as $`𝚽=\mathrm{\Omega }(t)+𝚽_a`$ in terms of the constant phase $`\mathrm{\Omega }(t)=_k\phi _k(t)`$ plus the operator
$$𝚽_a=\underset{k}{}𝚽_{ak}(t)=\underset{k}{}(\eta _k(t)a_k+h.c).$$
(48)
The decoherence factor can be calculated similarly as
$$D(t)=\alpha |e^{i𝚽}|\alpha =e^{i\varphi (t)\xi (t)},$$
(49)
where $`\xi (t)=_k\xi _k(t)`$ and
$`\varphi (t)`$ $`=`$ $`2{\displaystyle \underset{k}{}}\{\mathrm{𝐈𝐦}[{\displaystyle \frac{g_k\alpha _k}{\omega _k}}(1e^{i\omega _kt})]`$ (50)
$`+{\displaystyle \underset{k}{}}{\displaystyle \frac{|g_k|^2}{\omega _k^2}}(\omega _kt\mathrm{sin}\omega _kt).`$
It is easy to check that the phase generated by the quantum control is just the average value of the phase operator
$`\alpha |𝚽|\alpha `$ $`=`$ $`\alpha |𝚽_a|\alpha +\mathrm{\Omega }(t)`$ (51)
$`=`$ $`{\displaystyle \underset{k}{}}(\eta _k(t)\alpha _k+h.c)+{\displaystyle \underset{k}{}}\phi (t).`$
The analytical expression of the phase fluctuation is
$`(\mathrm{\Delta }\mathrm{\Phi })^2`$ $`=`$ $`(\mathrm{\Delta }\mathrm{\Phi }_a)^2={\displaystyle \underset{k=1}{\overset{N}{}}}(\mathrm{\Delta }\mathrm{\Phi }_{ak})^2`$
$`=`$ $`{\displaystyle \underset{k=1}{\overset{N}{}}}|\eta _k|^2=2{\displaystyle \underset{k=1}{\overset{N}{}}}\xi _k(t)=2\xi (t),`$
where we have considered each uncertain phase change as an independent stochastic variable. Namely, the relation $`\xi (t)=(\mathrm{\Delta }\mathrm{\Phi })^2/2`$ or the exact expression $`D(T)=\mathrm{exp}(i𝚽(\mathrm{\Delta }𝚽)^2/2)`$ still holds for the multi-mode case with the specialized initial state. Correspondingly, the error measure is estimated as
$$\epsilon =1F(t)=\lambda (\mathrm{\Delta }𝚽_a)^2=2\lambda \xi (t),$$
(53)
where $`\lambda =|c_0|^2|c_1|^2`$. Different from the single mode case, it is hard to find a proper instant $`T`$ such that $`\epsilon =\xi (T)=0`$ in general. Namely, it is hard to achieve an ideal quantum control without any error.
In the above discussions, the realization of quantum control boils down to the appearance of the $`c`$-number phase $`\varphi (t)`$ that contains the controllable part depending on the initial state of the controller. An ideal quantum control means the vanishing error $`\lambda (\mathrm{\Delta }𝚽_a)^2`$. But it is almost impossible because of the intrinsic decoherence due to quantum control itself. In fact, if the electromagnetic field could carry out a completely efficient control on the controlled system, then the interaction Hamiltonian should not commute with that of the controller. These facts are responsible for the inaccuracy of the phase gate or decoherence in the controlled system under the quantum control. We have to point out that the conclusion drawn above seems to depend on the choice of the initial state, but now we can argue that this is not the case with the above consideration. So we need to consider the universality of the conclusions.
Physically, every variable of the controller can independently exert a different impact on the different components of controller state. Since every uncertain phase is an independent stochastic variable, we have
$$(\mathrm{\Delta }𝚽_𝐚)^2=\underset{k=1}{\overset{N}{}}(\mathrm{\Delta }𝚽_{ak})^2N\mathrm{min}\{(\mathrm{\Delta }𝚽_{ak})^2|k=1,2..N\}$$
for a general initial state of the controller. We note that the phase uncertainty $`(\mathrm{\Delta }𝚽_𝐚)^2`$ caused by the controller variables can be amplified to a number much larger than unity when $`N\mathrm{}`$, i.e., the system states acquire a very large random phase factor. The decay factor
$$|D(t)|=e^{\frac{1}{2}(\mathrm{\Delta }𝚽)^2}\mathrm{exp}(\frac{N}{2}\mathrm{min}\{(\mathrm{\Delta }𝚽_{ak})^2|).$$
So $`|D(t)|0`$ when $`N\mathrm{}`$, i.e., the macroscopic controller can wash out the quantum coherence of the controlled system.
To be more concrete we assume that, in the initial state $`|0=_{k=1}^N|\psi _k`$ of the controller, each component $`|\psi _k`$ is a wave packet, symmetric with respect to both the “canonical coordinate” $`x_k=(a_k+a_k^{})/\sqrt{2}`$ and the corresponding “canonical momentum” $`p_k=i(a_ka_k^{})/\sqrt{2}`$. So $`x_k\psi _k|x_k|\psi _k=0`$ and $`p_k=0`$. We do not need the concrete form of the initial state. For convience we assume it to be of Gaussian type with the variance $`\sigma _k`$ $`=\mathrm{\Delta }x_k`$ in $`x_kspace`$. Physically, once $`\mathrm{\Delta }x_k`$ is given, the variance of $`p_k`$ cannot be arbitrary since there is a Heisenberg’s position-momentum uncertainty relation $`\mathrm{\Delta }x_k\mathrm{\Delta }p_k1/2`$. In the following we will show that the uncertainty relation will give a low bound to the variance of $`\mathrm{\Delta }𝚽_a`$. In the above reasoning about $`|D(t)|0`$ when $`N\mathrm{}`$, we have considered that there exists a finite minimum value of $`(\mathrm{\Delta }𝚽_{ak})^2`$. In the quantum measurement theory, the finite minimum value of $`(\mathrm{\Delta }𝚽_{ak})^2`$ is implied by the so called standard quantum limit (SQL) on the continuous measurement of phase operator.
To see this we rewrite the phase operator
$$𝚽_a=\underset{k}{}𝚽_{ak}=\underset{k}{}[\varkappa _k(t)x_k+\mu _k(t)p_k],$$
(54)
in terms of the “canonical coordinate” and the corresponding “canonical momentum”, and the coefficients are
$`\varkappa _k(t)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\eta _k(t)+\eta _k^{}(t)),`$
$`\mu _k(t)`$ $`=`$ $`{\displaystyle \frac{i}{\sqrt{2}}}(\eta _k(t)\eta _k^{}(t)).`$
The existence of SQL is guaranteed by the Heisenberg’s position-momentum uncertainty relation. Because each $`𝚽_{ak}=\varkappa _k(t)x_k+\mu _k(t)p_k`$ is a linear combination of $`x_k`$ and $`p_k`$ with a property $`x_kp_k+p_kx_k=0`$ for the average over the real initial state. The phase fluctuation $`\mathrm{\Delta }𝚽_{ak}`$ can be derived as
$`\mathrm{\Delta }𝚽_{ak}`$ $`=`$ $`\sqrt{|\varkappa _k(t)|^2(\mathrm{\Delta }x_k)^2+|\mu _k(t)|^2(\mathrm{\Delta }p_k)^2}`$
$``$ $`\sqrt{|\varkappa _k(t)\mu _k(t)|},`$
or
$$(\mathrm{\Delta }𝚽_{ak})^28\frac{g_k^2}{\omega _k^2}|\mathrm{sin}^3\frac{\omega _kt}{2}\mathrm{cos}\frac{\omega _kt}{2}|.$$
(55)
Here, we considere the variance $`\mathrm{\Delta }(\xi x)=|\xi |(\mathrm{\Delta }x)`$ for a stochastic variable $`x`$ and a real number $`\xi `$, and suppose $`g_k/\omega _k`$ being a real number.
In the above arguments, $`x_k`$ and $`p_k`$ are not only regarded as a pair of uncorrelated stochastic variables in the terminology of classical stochastic process, the uncertainty relation $`\mathrm{\Delta }x_k\mathrm{\Delta }p_k1/2`$ of them is also taken into account. This constraint just reflects the uncertainty of phase change in the quantum control process. Therefore, we have a time dependent minimum value of phase uncertainty with a low bound
$$(\mathrm{\Delta }𝚽_𝐚)^2N\mathrm{min}\{|\varkappa _k(t)\mu _k(t)||k=1,2..N\}.$$
This result qualitatively illustrates the many-particle amplification effect of uncertain phase change due to quantum control itself. The large random phase variance $`(\mathrm{\Delta }𝚽_𝐚)^2`$ implies that it is hard to satisfy the exact condition $`(\mathrm{\Delta }𝚽_𝐚)^2=0`$ in principle, and thus one can only optimize both the system parameters and the initial state of the controller to approach what we want.
To see the above observation analytically, we calculate $`𝚽`$ in comparison with $`\mathrm{\Delta }𝚽`$ in the decoherence factor $`D(T)=\mathrm{exp}(i𝚽(\mathrm{\Delta }𝚽)^2/2)`$. The most simple, but somewhat trivial case is that all modes are degenerate, i.e., $`g_k=g`$ and $`\omega _k=\omega `$, then
$$\mathrm{\Delta }𝚽=\sqrt{8N}\frac{|g|}{\omega }|\mathrm{sin}\frac{\omega t}{2}|,$$
(56)
while the phase we wanted is
$`\varphi (t)`$ $`=`$ $`2N\{\mathrm{𝐈𝐦}[{\displaystyle \frac{g\alpha }{\omega }}(1e^{i\omega t})]`$ (57)
$`+2N{\displaystyle \frac{|g|^2}{\omega ^2}}(\omega t\mathrm{sin}\omega t).`$
Obviously, for very large $`N`$, the phase fluctuation $`\mathrm{\Delta }𝚽`$ can be neglected since $`\mathrm{\Delta }𝚽/|\varphi (t)|1/\sqrt{N}0`$. In general, we need to consider divergence of the phase fluctuation
$`(\mathrm{\Delta }𝚽)^2`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{N}{}}}16{\displaystyle \frac{g_k^2}{\omega _k^2}}\mathrm{sin}^2{\displaystyle \frac{\omega _kt}{2}}`$ (58)
$`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}16{\displaystyle \frac{g_k^2}{\omega _k^2}}\rho (\omega _k)\mathrm{sin}^2{\displaystyle \frac{\omega _kt}{2}}d\omega _k`$
for various spectrum distributions of the controller, where an unspecific spectrum distribution $`\rho (\omega _k)`$ is used to discuss the case with continuous spectrum. For example, when $`\rho (\omega _k)=\gamma /g_k^2`$, the decoherence factor is exponentially decaying since the above integral converges to a number $`8\pi \gamma t/9`$ proportional to time $`t`$. Another example is the Ohmic distribution $`\rho (\omega _k)=2\eta \omega _k^2/(\pi g_k^2)`$, which results in a diverging phase fluctuation for $`t0`$.
## V Low bound of the control induced decoherence and quantum computation
In this section we will show that, it is the back-action of the controller on the controlled system, implied by Heisenberg’s position-momentum uncertainty relation, that disturb the phases of states of the controlled system and then induce a quantum decoherence, which is relevant to the SQL. In order to quantitatively characterized such limit to the quantum controllability, we now return to the discussion about the quantum control with multi-mode field initially prepared in a coherent state.
The commutation relation of the number operator $`𝐍`$ and the phase operator $`𝚽_a`$ defines an operator
$$𝚯=i\underset{k}{}(\eta _k(t)a_k+\eta _k^{}(t)a_k^{})$$
(59)
dual to the phase operator $`𝚽_a,`$that is
$$𝚯=i[𝐍,𝚽_a].$$
(60)
To see the meaning of the defined $`𝚯,`$ we calculate the commutation relation of $`𝐍`$ and $`𝚽_a`$ to find a close algebra by
$`[𝐍,𝚯]`$ $`=`$ $`i𝚽_a,`$
$`[𝚽_a,𝚯]`$ $`=`$ $`iF(t)`$ (61)
where $`F(t)=2_k|\eta _k(t)|^2`$ is a time dependent constant. This means that $`𝐏=𝚯/F(t)`$ is a conjugate variable with respect to $`𝚽_a`$ since we have the canonical commutation relation $`[𝚽_a,𝐏]=1.`$ In this sense we call $`𝚯`$ a dual phase operator (DPO). A constant uncertainty relation can be found for $`𝚽_a`$ and $`𝐏`$ , which can be minimized by the corresponding coherent state $`|\alpha =_k|\alpha _k`$.
The above arguments about minimization of the uncertainty by $`[𝚽_a,𝚯]`$ can enlighten us to find a low bound for the control induced decoherence. To this end we consider the uncertainty relation
$$𝐍\mathrm{\Delta }𝚽_a=\mathrm{\Delta }𝐍\mathrm{\Delta }𝚯\frac{1}{2}|[𝚯,𝐍]|=\frac{1}{2}|𝚽_a|$$
(62)
about DPO and the photon number operator $`𝐍=_ka_k^{}a_k`$.
To derive the above uncertainty relation (62), we have considered
$`\mathrm{\Delta }𝐍`$ $`=`$ $`𝐍,`$
$`(\mathrm{\Delta }𝚽_a)^2`$ $`=`$ $`(\mathrm{\Delta }𝚯)^2`$ (63)
for the average $``$ over the coherent state $`|\alpha `$. We check the above results (63) by the straightforward calculations
$`(\mathrm{\Delta }\mathrm{\Phi }_a)^2`$ $`=`$ $`\alpha |\mathrm{\Phi }_a^2\mathrm{\Phi }_a^2|\alpha ={\displaystyle \underset{k}{}}|\eta _k(t)|^2,`$
$`(\mathrm{\Delta }𝚯)^2`$ $`=`$ $`\alpha |𝚯^2𝚯^2|\alpha ={\displaystyle \underset{k}{}}|\eta _k(t)|^2.`$
The novel uncertainty relation (62) defines a low bound for the phase variation $`\mathrm{\Delta }\mathrm{\Phi }_a`$ for a given phase $`𝚽_a`$ to be achieved by the quantum control, i.e.
$$\mathrm{\Delta }𝚽_a\frac{|𝚽_a|}{2𝐍}.$$
(64)
Eq.(64) clearly implies that we need much larger energy to reduce the low bound of the phase fluctuation. Actually, we can formally write down the expectation of the photon energy of the controller
$$E=\mathrm{}\underset{k}{}\omega _ka_k^{}a_k\mathrm{}𝐍\omega $$
(65)
in terms of the average photon number $`𝐍=_k|\alpha _k|^2`$ and the average frequency of photons
$$\omega =\frac{_k\omega _k|\alpha _k|^2}{_k|\alpha _k|^2}.$$
(66)
Then Eq.(64) becomes
$$\mathrm{\Delta }\mathrm{\Phi }_a\frac{\mathrm{}\omega }{2E}|𝚽_a|.$$
(67)
The small low bound requires that a large quantum controller (implied by large $`𝐍`$ or large energy $`E`$) possesses a very small average frequency. In this sense Eq.(62) defines a necessary condition for the quantum control that can manipulate the qubit system reaching the target state accurately. This requirement is very similar to that the apparatus should be sufficiently so “large” that to be ”classical” in the quantum measurement in the so-called “Copenhagen interpretation”. Since the quantum control relies on the ability to preserve quantum coherence of the qubit system during controlling it, the controller should be much “larger” than the controlled system. In this sense, the bac-action of the qubit system on the controller can be neglected.
Next we consider the controllable condition (27) that the controller field is switched on and off over a time duration $`T`$, which can be roughly realized as a periodical phenomenon with the average period $`T2\pi /\omega `$. Since the average frequency of the field can be approximated by $`\omega 2\pi /T`$, there is a low bound
$$\epsilon \frac{\lambda h^2}{4E^2T^2}|𝚽_a|^2\frac{\lambda h^2}{4S^2}|𝚽_a|^2$$
(68)
for the error measure estimation of $`\epsilon \lambda (\mathrm{\Delta }𝚽_a)^2`$ of the quantum control. So the larger action $`S=ET`$ from the controller is brought on the qubit system, the less quantum decoherence characterized by the control induced error $`\epsilon `$ becomes; the more one wants to change by the phase $`𝚽_a`$ of the qubit system, the larger quantum decoherence is induced by the quantum control. Therefore Eq.(68) imposes a fundamental limit on the accuracy of quantum control. In the following we can consider this physical limit for quantum computing
It is well known that the controllability of qubits is a basic requirement for universal quantum computations, but according to the above arguments a well-armed controls in quantum computing would cause the extra decoherence in the qubits system. Thus in competition with the induced decoherence, the controllability for quantum computation is limited.
In the last section a low bound of decoherence from the quantum control is obtained. It throws an accuracy limit to the quantum controls in quantum computation. According to Eq.(68) this limit is about $`10^{20}`$ for the typical setting $`\mathrm{\Phi }_a=\pi `$, $`E=10^9J`$ and $`T=1\mu s`$ in an ion trap schemes. This is such a small limit that it is negligible in comparison with other errors, such as the environments induced errors in the current experiments of implementing quantum computation. However, in principle, Eq.(68) do throw a fundamental limit to the accuracy of the quantum control and thus on quantum computations. There are some numerical estimates in Fig.3, which demonstrate the similar limit to the power of quantum computer. It is known that for an algorithm consisting of $`L`$ operations on the qubits system, an up bound of error $`ϵ1/L`$ is required in each operation for a faithful result of the entire computation. The inequality Eq.(68) tells us that the minimum amount of time needed for a single gate is
$$t_{\mathrm{min}}=\frac{h\sqrt{\lambda L}}{2E}|\mathrm{\Phi }_a|$$
(69)
and so the total time needed to carry out a particular algorithm consisting of $`L`$ elementary gates is about $`Lt_{\mathrm{min}}`$.
For a general algorithm as an arbitrary unitary operation on $`nqubit`$, the amount of elementary gates $`L`$ needed is about $`O(n^24^n)`$ Barenco ; for the Grover algorithm on $`nqubit`$, the amount is about $`O(\sqrt{2^n})`$; for the Shor large number factorization, the amount is about $`O(n^3log^32)`$. The time duration needed for a general algorithm, the Grover algorithm and the Shor’s algorithm are estimated with the optimistic assumptions, in which the only restriction is from the quantum control. Thus, in Fig.3 it could be found that the practise of quantum computation heavily depends on the sophisticated quantum algorithms and arbitrary quantum operations on about several tens qubits is already inaccessible even in principle. This handicap on the quantum computation stands when the quantum computation is carried out by tandem elementary gates under the quantum control.
## VI Conclusion
In this paper we present a universal description for the quantum control based on the quantized controller. We discovered the complementarity about the competition between the controllability and the control induced quantum decoherence in the view of quantum measurement. Starting with an exactly-soluble example, a general model of quantum control is proposed to describe this novel complementarity or a new type of uncertainty relation. Our investigations show that it is possible to realize the decoherence free quantum controls only with some special phases at the finite energy scale and in finite time. If the parameters of phase is to be determined by the initial state of the controller, then there exists a low bound for the systematic errors resulted from the decoherence cause by the quantum control itself.
The above arguments also show that the decoherences from the quantum control are different from those induced by the environments through the unwanted interactions. This is because the negative influence of the controller happens in the quantum control process itself. If one eliminates this influence out and out, the positive role of the quantum control would perish together. Therefore, for quantum computing, these kinds of errors induced by the control itself can not be overcame totally by the conventional error management protocols ECCs . At least, it has not been proved that the control induced decoherences can also be conquered efficiently by the well-estabilished error management protocols. The better method to solve this problem is to optimize the control operations when the target of control is given. Without doubt, this is an open question to challenge for the physical implementation of quantum computing as well as the other protocols of quantum information processing.
This work was supported by CNSF (Grant Nos. 90203018, 90303023, 10474104 and 60433050. It is also funded by the National Fundamental Research Program of China with Nos. 2001CB309310 and 2005CB724508. |
warning/0507/gr-qc0507114.html | ar5iv | text | # Probing anisotropies of gravitational-wave backgrounds with a space-based interferometer II: perturbative reconstruction of a low-frequency skymap
## I Introduction
The gravitational-wave background, incoherent superposition of gravitational waves coming from many unresolved point-sources and/or diffuse-sources would be a cosmological gold mine to probe the dark side of the Universe. While these signals are generally random and act as confusion noises with respect to a periodic gravitational-wave signal, the statistical properties of gravitational-wave backgrounds (GWBs) carry valuable information such as physical processes of gravitational-wave emission, source distribution and populations. Moreover, the extremely early universe beyond the last-scattering surface of cosmic microwave background (CMB) can be directly explored using GWB. Therefore, GWBs may be regarded as an ultimate cosmological tool alternative to the CMB.
Currently, several grand-based gravitational-wave detectors are now under scientific operation, and the search for gravitational waves enters a new era. There are several kinds of target GW sources in these detectors, such as, coalescence of neutron star binaries and core-collapsed supernova. On the other hand, planned space-based detector, LISA (Laser Interferometer Space Antenna) and the next-generation detectors, e.g., DECIGO Seto et al. (2001) and BBO BBO (2003) aim at detecting gravitational waves in low-frequency band $`0.1`$mHz –$`0.1`$Hz, in which many detectable candidates for GWBs exist. Among them, the GWB originated from the inflationary epoch may be detected directly in this band (e.g., Ungarelli and Vecchio (2001a); Smith et al. (2005)). Therefore, for future application to cosmology, various implications for these GWBs should deserve consideration both from the theoretical and the observational viewpoint.
One fundamental issue to access a new subject of cosmology is to make a skymap of GWB as a first step. Similar to the case of the CMB, the intensity map of the GWBs is of particular interest and it provides valuable information, which plays a key role to clarify the origin of GWBs. Importantly, the low-frequency GWBs observed by LISA are expected to be anisotropic due to the contribution of Galactic binaries to the confusion noise (Hils et al. (1990); Bender and Hils (1997); Nelemans et al. (2001), see also recent numerical simulations Benacquista et al. (2004); Edlund et al. (2005); Timpano et al. (2005)). Thus, the GWB skymap is potentially useful to discriminate the individual backgrounds from many superposed GWBs, as well as to identify the underlying physical processes. The basic idea to create a skymap of GWB is to use the directional information obtained through the time-modulation of the correlation signals, which is caused by the motion of gravitational-wave detector. As detector’s antenna pattern sweeps out the sky, the amplitude of the gravitational-wave signal would gradually change in time if the intensity distribution of GWB is anisotropic. Using this, a method to explore anisotropies of GWB has been proposed Giampieri and Polnarev (1997a, b); Allen and Ottewill (1997). Later, the methodology was applied to study the anisotropic GWB observed by space interferometer, LISA Ungarelli and Vecchio (2001b); Cornish (2001, 2002a); Seto (2004); Seto and Cooray (2004).
In the previous paper Kudoh and Taruya (2005a), which is referred to as Paper I in the present paper, we have investigated the directional sensitivity of space interferometer to the anisotropy of GWB. Particularly focusing on the geometric properties of antenna pattern functions and their angular power, we found that the angular sensitivity to the anisotropic GWBs is severely restricted by the data combination and the symmetry of detector configuration. As a result, in the case of the single LISA detector, detectable multipole moments are limited to $`\mathrm{}8`$$`10`$ with the effective strain sensitivity $`h10^{20}`$ Hz<sup>-1/2</sup>. This is marked contrast to the angular sensitivity to the chirp signals emitted from point sources, in which the angular resolution can reach at a level of a square degree or even better than that Cutler (1998); Moore and Hellings (2002); Peterseim et al. (1997); Takahashi and Nakamura (2003).
Despite the poor resolution of LISA detector with respect to GWBs, making a skymap of GWBs is still an important general issue and thus needs to be investigated. In the present paper, we continue to investigate the map-making problem and consider how the intensity map of the GWBs is reconstructed from the time-modulation signals observed via space interferometer. In particular, we are interested in the low-frequency GWB, the wavelength of which is longer than the arm-length of the detector. In such case, the frequency dependence of the detector response becomes simpler and a perturbative scheme based on the low-frequency expansion of the antenna pattern functions can be applied. Owing to the least-squares approximation, we present a robust reconstruction method. The methodology is quite general and is also applicable to the map-making problem in the case of the ground detectors. We demonstrate how the present reconstruction method works well in a specific example of GWB source, i.e., Galactic confusion-noise background. With a sufficient high signal-to-noise ratio for each anisotropic components of GWB, we show that the space interferometer LISA can create the low-frequency skymap of Galactic GWB with angular resolution $`\mathrm{}5`$. Since the resultant skymap is obtained in a non-parametric way without any assumption of source distributions, we hope that despite the poor angular resolution, the present methodology will be helpful to give a tight constraint on the luminosity distribution of GWBs if we combine it with the other observational techniques.
The organization of the paper is as follows. In Sec. II, a brief discussion on the detection and the signal processing of anisotropic GWBs is presented, together with the detector response of space interferometer. Sec.III describes the details of the reconstruction of a GWB skymap in low-frequency regime. Owing to the least-squares approximation and the low-frequency expansion, a perturbative reconstruction scheme is developed. Related to this (in the case of LISA), and we give an important remark on the degeneracy between some multipole moments (Appendix B). In Sec.IV, the reconstruction method is demonstrated in the case of the Galactic GWB. The signal-to-noise ratios for anisotropy of GWB are evaluated and the feasibility to make a skymap of GWB is discussed. Finally, Section V is devoted to a summary and conclusion.
## II Basic formalism
### II.1 Correlation analysis
Let us begin with briefly reviewing the signal processing of gravitational-wave backgrounds based on the correlation analysis Kudoh and Taruya (2005a). Stochastic gravitational-wave backgrounds are described by incoherent superposition of plane gravitational waves $`𝐡=h_{ij}`$ given by
$`𝐡(t,𝐱)={\displaystyle \underset{A=+,\times }{}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑f{\displaystyle 𝑑𝛀e^{i\mathrm{\hspace{0.17em}\hspace{0.17em}2}\pi f(t𝛀𝐱)}\stackrel{~}{h}_A(f,𝛀)𝐞^A(𝛀)}.`$ (1)
The Fourier amplitude $`\stackrel{~}{h}_A(f,𝛀)`$ of the gravitational waves for the two polarization modes $`𝐞^A`$ ($`A=+,\times `$) is assumed to be characterized by a stationary random process with zero mean $`\stackrel{~}{h}_A=0`$. The power spectral density $`S_h`$ is then defined by
$`\stackrel{~}{h}_A^{}(f,𝛀)\stackrel{~}{h}_A^{}(f^{},𝛀^{})`$ $`=`$ $`{\displaystyle \frac{1}{2}}\delta (ff^{}){\displaystyle \frac{\delta ^2(𝛀𝛀^{})}{4\pi }}\delta _{AA^{}}S_h(|f|,𝛀),`$ (2)
where $`𝛀`$ is the direction of a propagating plane wave. Note that the statistical independence between two different directions in the sky is implicitly assumed in equation (2), which might not be generally correct. Actually, CMB skymap exhibits a large-angle correlation between the different skies, which reflects the primordial density fluctuations. For the GWB of our interest, the wavelength of the tensor fluctuations detected via space interferometers is much shorter than the cosmological scales and thereby the assumption put in equation (2) is practically valid.
The detection of a gravitational-wave background is achieved through the correlation analysis of two data streams. The planned space interferometer, LISA and also the next generation detectors DECIGO/BBO constitute several spacecrafts, each of which exchanges laser beams with the others. Combining these laser pulses, it is possible to synthesize the various output streams which are sensitive (or insensitive) to the gravitational-wave signal. The output stream for a specific combination $`I`$ denoted by $`s_I(t)`$ is generally described by a sum of the gravitational-wave signal $`h_I(t)`$ and the instrumental noise $`n_I(t)`$ by
$$s_I(t)=h_I(t)+n_I(t).$$
We assume that the noise $`n_I(t)`$ is treated as a Gaussian random process with spectral density $`S_\mathrm{n}(f)`$ and zero mean $`n_I=0`$. The gravitational-wave signal $`h_I`$ is obtained by contracting $`𝐡`$ with detector’s response function and/or detector tensor $`𝐃:=D_{ij}`$:
$$h_I(t)=\underset{A=+,\times }{}_{\mathrm{}}^+\mathrm{}𝑑f𝑑𝛀e^{i2\pi f(t𝛀𝐱_I)}\text{D}_I(𝛀,f;t)\text{:}𝐞^A(𝛀)\stackrel{~}{h}_A(f,𝛀).$$
(3)
Note that the response function explicitly depends on time. The time variation of response function is caused by the orbital motion of the detector and it plays a key role in reconstructing a skymap of the GWBs. Typically, the orbital frequency of the detector motion is much lower than the observed frequency and within the case, the expression (3) is validated.
Provided the two output data sets, the correlation analysis is examined depending on the strategy of data analysis, i.e., self-correlation analysis using the single data stream or cross-correlation analysis using the two independent data stream. Defining $`S_{IJ}(t)s_I(t)s_J(t)`$, the Fourier counterpart of it $`\stackrel{~}{S}_{IJ}(t,f)`$, which is related with $`S_{IJ}(t)`$ by $`S_{IJ}(t)=𝑑f\stackrel{~}{S}_{IJ}(t,f)`$, becomes <sup>1</sup><sup>1</sup>1 In Paper I, we had used $`C_{IJ}`$ to denote the output data $`s_I(t)s_J(t)`$ itself. This coincides with the present definition if we neglect the instrumental noises.
$`\stackrel{~}{S}_{IJ}(t,f)`$ $`=`$ $`\stackrel{~}{C}_{IJ}(t,f)+\delta _{IJ}S_n(|f|),`$ (4)
where we define
$`\stackrel{~}{C}_{IJ}(t,f)={\displaystyle \frac{d𝛀}{4\pi }S_h(|f|,𝛀)_{IJ}^E(f,𝛀;t)}.`$ (5)
Here, the function $`_{IJ}^E`$ is the so-called antenna pattern function defined in an ecliptic coordinate, which is expressed in terms of detector’s response function:
$`_{IJ}^E(f,𝛀;t)=e^{i\mathrm{\hspace{0.17em}2}\pi f𝛀(\text{x}_I\text{x}_J)}{\displaystyle \underset{A=+,\times }{}}F_I^A(𝛀,f;t)F_J^A(𝛀,f;t)`$ (6)
$`F_I^A(𝛀,f;t)=𝐃_I(𝛀,f;t)\text{:}\text{e}^A(𝛀)`$ (7)
Apart from the second term, equation (4) implies that the luminosity distribution of GWBs $`S_h(f,𝛀)`$ can be obtained by deconvolving the all-sky integral of antenna pattern function from the time-series data $`S_{IJ}(t)`$. To see this more explicitly, we focus on equation (5) and decompose the antenna pattern function and the luminosity distribution into spherical harmonics in an ecliptic coordinate, i.e., sky-fixed frame:
$`S_h(|f|,𝛀)={\displaystyle \underset{\mathrm{},m}{}}[p_\mathrm{}m(f)]^{}Y_\mathrm{}m^{}(𝛀),`$ (8)
$`_{IJ}^E(f,𝛀;t)={\displaystyle \underset{\mathrm{},m}{}}a_\mathrm{}m^E(f,t)Y_\mathrm{}m(𝛀).`$ (9)
Note that the properties of spherical harmonics yield $`p_\mathrm{}m^{}=(1)^mp_{\mathrm{},m}`$ and $`a_\mathrm{}m^{}=(1)^\mathrm{}ma_{\mathrm{},m}`$, where the latter property comes from $`_{IJ}^{}(f,\mathrm{\Omega };t)=_{IJ}(f,\mathrm{\Omega };t)`$ Kudoh and Taruya (2005a). Substituting (9) into (5) becomes
$`\stackrel{~}{C}_{IJ}(t,f)={\displaystyle \frac{1}{4\pi }}{\displaystyle \underset{\mathrm{}m}{}}\left[p_\mathrm{}m(f)\right]^{}a_\mathrm{}m^E(f,t),`$ (10)
where we have dropped the contribution from the detector noise. The expression (10) still involves the time-dependence of antenna pattern functions. To eliminate this, one may further rewrite equation (10) by employing the harmonic expansion in detector’s rest frame. We denote the multipole coefficients of the antenna pattern in detector’s rest frame by $`a_\mathrm{}m`$. The transformation between the detector rest frame and the sky-fixed frame is described by a rotation matrix by the Euler angles $`(\psi ,\vartheta ,\phi )`$, whose explicit relation is expressed in terms of the Wigner $`D`$ matrices <sup>2</sup><sup>2</sup>2 We are specifically concerned with the time dependence of directional sensitivity of antenna pattern functions, not the real orbital motion. Hence, only the time evolution of directional dependence is considered in the expression (11). Allen and Ottewill (1997); Cornish (2002a); Edmonds (1957):
$`a_\mathrm{}m^E(f,t)={\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}e^{in\psi }d_{nm}^{\mathrm{}}(\vartheta )e^{im\phi }a_\mathrm{}n(f).`$ (11)
Here the Euler rotation is defined to perform a sequence of rotation, starting with a rotation by $`\psi `$ about the original $`z`$ axis, followed by rotation by $`\vartheta `$ about the original $`y`$ axis, and ending with a rotation by $`\phi `$ about the original $`z`$ axis. The Wigner $`D`$ matrices for $`nm`$ is
$`d_{nm}^{\mathrm{}}(\theta )=(1)^\mathrm{}n\sqrt{{\displaystyle \frac{(\mathrm{}+n)!(\mathrm{}n)!}{(\mathrm{}+m)!(\mathrm{}m)!}}}\left(\mathrm{cos}{\displaystyle \frac{\theta }{2}}\right)^{n+m}\left(\mathrm{sin}{\displaystyle \frac{\theta }{2}}\right)^{nm}P_\mathrm{}n^{(n+m,nm)}(\mathrm{cos}\theta )`$ (12)
with $`P_n^{(a,b)}`$ being Jacobi polynomial. For $`n<m`$, we have $`d_{nm}^{\mathrm{}}=(1)^{nm}d_{mn}^{\mathrm{}}`$.
Let us now focus on the orbital motion of the LISA constellation (see Fig.1). The LISA orbital motion can be expressed by $`\psi =\omega t`$, $`\vartheta =\pi /3`$, $`\phi =\omega t`$, where $`\omega =2\pi /T_0`$ is the orbital frequency of LISA ($`T_0=1`$ sidereal year) <sup>3</sup><sup>3</sup>3 The relation $`\psi =\phi `$ does not necessarily hold for orbital motion of LISA and there may be some possibilities to impose a constant phase difference, i.e., $`\psi =\phi +c`$. However, for the sake of simplicity, we put $`c=0`$. Since the antenna pattern function is periodic in time due to the orbital motion, the expected signals also vary in time with the same period $`T_0`$. It is therefore convenient to express the output signals by
$$\stackrel{~}{C}_{IJ}(t,f)=\underset{k=\mathrm{}}{\overset{+\mathrm{}}{}}\stackrel{~}{C}_{IJ,k}(f)e^{ik\omega t}.$$
(13)
Using the relation (11) with specific parameter set, the Fourier component $`\stackrel{~}{C}_{IJ,k}(f)`$ is then given by
$`\stackrel{~}{C}_k(f)`$ $``$ $`{\displaystyle \frac{1}{T_0}}{\displaystyle _0^{T_0}}𝑑te^{ik\omega t}\stackrel{~}{C}_{IJ}(t,f)`$ (14)
$`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle \underset{\mathrm{}=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}k}{}}}[p_\mathrm{}m(f)]^{}d_{(m+k),m}^{\mathrm{}}\left({\displaystyle \frac{\pi }{3}}\right)a_{\mathrm{},(m+k)}(f)`$
for $`k0`$. As for $`k<0`$, the lower and the upper limit of the sum over $`m`$ are changed to $`m=\mathrm{}k`$ and $`m=\mathrm{}`$, respectively.
Equation (14) as well as (5) is the theoretical basis to reconstruct the skymap of GWBs. Given the output data $`\stackrel{~}{C}_k(f)`$ (or $`\stackrel{~}{C}_{IJ}(t,f)`$) experimentally, the task is to solve the linear system (14) with respect to $`p_\mathrm{}m(f)`$ for given antenna pattern functions. One important remark deduced from equation (14) is that the accessible multipole coefficients $`p_\mathrm{}m`$ are severely restricted by the angular sensitivity of antenna pattern functions. The important properties of the antenna pattern functions are summarized in next subsection. Another important message is that the above linear systems are generally either over-constrained or under-determined. For the expression (14), if one truncates the multipole expansion with $`\mathrm{}=\mathrm{}_{\mathrm{max}}`$, the system consists of $`\frac{1}{2}(\mathrm{}_{max}+1)(2\mathrm{}_{max}+1)`$ unknowns and $`N(2\mathrm{}_{max}+1)`$ equations (for $`k0`$), where $`N`$ is the number of available modes of the antenna pattern functions. Thus this deconvolution problem is, in principle, over-determined for a relatively small truncation multipole $`\mathrm{}_{\text{max}}`$, while it becomes under-determined for a larger value of $`\mathrm{}_{\text{max}}`$. We will later discuss the over-determined case in Sec.IV.3.1, where $`\mathrm{}_{\text{max}}=5`$ and $`N5`$.
### II.2 Detector response and antenna pattern functions
The output signals of space interferometer sensitive to the gravitational waves are constructed by time-delayed combination of laser pulses. In the case of LISA, the technique to synthesize data streams canceling the laser frequency noise is known as time-delay interferometry (TDI), which is crucial for our subsequent analysis. In the present paper, we use the optimal set of TDI variables independently found by Prince et al. Prince et al. (2002) and Nayak et al.Nayak et al. (2003), which are free from the noise correlation <sup>4</sup><sup>4</sup>4 The optimal TDI variables adopted here may not be the best TDI combinations for the present purpose. There might be a better choice of the signal combinations, although a dramatic improvement of the angular sensitivity would not be expected. . A simple realization of such data set is obtained from a combination of Sagnac signals, which are the six-link observables using all six LISA oriented arms. For example, the Sagnac signal S<sub>1</sub> measures the phase difference accumulated by two laser beams received at space craft $`1`$, each of which travels around the LISA array in clockwise or counter-clockwise direction. The explicit expression of detector tensor for such signal, which we denote by $`𝐃_{\mathrm{S}_1}`$, is given in equation (21) of Paper I (see also Cornish (2002b); Armstrong et al. (1999)). The analytic expression for other detector tensors $`𝐃_{\mathrm{S}_2}`$ and $`𝐃_{\mathrm{S}_3}`$ are also obtained by the cyclic permutation of the unit vectors $`𝐚,𝐛`$ and $`𝐜`$.
Combining the three Sagnac signals, a set of optimal data combinations can be constructed Prince et al. (2002); Nayak et al. (2003) (see also Krolak et al. (2004)):
$`𝐃_\mathrm{A}={\displaystyle \frac{1}{\sqrt{2}}}(𝐃_{\mathrm{S}_3}𝐃_{\mathrm{S}_1}),`$
$`𝐃_\mathrm{E}={\displaystyle \frac{1}{\sqrt{6}}}(𝐃_{\mathrm{S}_1}2𝐃_{\mathrm{S}_2}+𝐃_{\mathrm{S}_3}),`$
$`𝐃_\mathrm{T}={\displaystyle \frac{1}{\sqrt{3}}}(𝐃_{\mathrm{S}_1}+𝐃_{\mathrm{S}_2}+𝐃_{\mathrm{S}_3}).`$ (15)
These three detector tensors are referred to as $`A,E,T`$-variables, which generate six kinds of antenna pattern functions. Notice that in the equal arm-length limit, frequency dependence of these functions is simply expressed in terms of the normalized frequency:
$$\widehat{f}\frac{f}{f_{}},$$
(16)
where the characteristic frequency is given by $`f_{}=c/(2\pi L)`$ with $`L`$ being the arm-length of detector. With the arm-length $`L=5\times 10^6`$ km, the characteristic frequency of LISA becomes $`f_{}9.54`$ mHz. In what follows, we use the analytic expressions for equal arm-length limit to demonstrate the reconstruction of GWB skymap. This is sufficient for the present purpose, because we are concerned with a fundamental theoretical basis to map-making capability of GWBs. The idea provided in the present paper would allow us to examine a more realistic situation and the same strategy can be applied to an extended analysis in the same manner.
Finally, we note that the antenna pattern functions for the optimal combinations of TDI have several distinctive features in angular sensitivity. In the low frequency limit ($`\widehat{f}1`$), the antenna pattern functions for the self-correlations and the cross-correlations are expanded as
$`_{AA}(f,𝛀)=`$ $`_{AA}^{(2)}(𝛀)\widehat{f}^2`$ $`+`$ $`_{AA}^{(4)}(𝛀)\widehat{f}^4`$ $`+O(\widehat{f}^6)`$ (17)
$`_{TT}(f,𝛀)=`$ $`_{TT}^{(4)}(𝛀)\widehat{f}^4`$ $`+O(\widehat{f}^6)`$
$`_{AE}(f,𝛀)=`$ $`_{AE}^{(2)}(𝛀)\widehat{f}^2`$ $`+`$ $`_{AE}^{(3)}(𝛀)\widehat{f}^3`$ $`+`$ $`_{AE}^{(4)}(𝛀)\widehat{f}^4`$ $`+O(f^5)`$
$`_{AT}(f,𝛀)=`$ $`_{AT}^{(3)}(𝛀)\widehat{f}^3`$ $`+`$ $`_{AT}^{(4)}(𝛀)\widehat{f}^4`$ $`+O(\widehat{f}^5)`$
The frequency dependence of $`_{EE}`$ and $`_{ET}`$ are the same as for $`_{AA}`$ and $`_{AT}`$, respectively. Since the leading term of the TT-correlation is $`O(\widehat{f}^4)`$, the $`TT`$-correlation becomes insensitive to the gravitational waves in the low-frequency regime. We will use all correlated data set except for the $`TT`$-correlation. In Table 1, we summarize the important properties for the multipole moments of antenna pattern functions. Also, in Appendix A, employing the perturbative approach based on the low-frequency approximation $`\widehat{f}1`$, the spherical harmonic expansion for antenna pattern functions are analytically computed, which will be useful in subsequent analysis.
## III Perturbative reconstruction method for GWB skymap
We are in position to discuss the methodology to reconstruct a skymap of GWB based on the expression (14) (or (5)). Since we are specifically concerned with low-frequency sources observed via LISA, it would be helpful to employ a perturbative approach using the low-frequency expansion of antenna pattern function. In Sec.III.1, owing to the expression (14), a perturbative reconstruction method is presented. Sec.III.2 discusses alternative reconstruction method based on the time-series representation (5).
### III.1 General scheme
Hereafter, we focus on the GWBs in the low-frequency band of the detector, the wavelength of which is typically longer than the arm-length of the gravitational detector, i.e., $`\widehat{f}1`$. Without loss of generality, we restrict our attention to the reconstruction of a GWB skymap in a certain narrow frequency range, $`ff+\mathrm{\Delta }f`$, within which a separable form of the GWB spectrum is a good assumption:
$$S_h(f,𝛀)=H(f)P(𝛀).$$
(18)
Further, for our interest of the narrow bandwidth, it is reasonable to assume that the spectral density $`H(f)`$ is described by a power-law form as $`H(f)=𝒩f^\alpha `$. In the reconstruction analysis discussed below, the spectral index $`\alpha `$ is assumed to be determined beforehand from theoretical prediction and/or experimental constraint <sup>5</sup><sup>5</sup>5In our general scheme, we do not assume a priori the normalization factor $`𝒩`$ in the function $`H(f)`$, which should be simultaneously determined with the reconstruction of an intensity distribution $`P(𝛀)`$. In the low-frequency approximation, $`\widehat{f}1`$, however, there exists a degeneracy between the monopole and the quadrupole components and one cannot correctly determine the normalization $`𝒩`$. (See Appendix.B.) .
In the low frequency approximation up to the order $`𝒪(\widehat{f}^3)`$, we have five output signals which respond to the GWBs, i.e., $`AA`$-, $`EE`$-, $`AE`$-, $`AT`$\- and $`ET`$-correlations, each output of which is represented by equation (14). Collecting these linear equations and arraying them appropriately, the linear algebraic equations can be symbolically written in a matrix form as
$`𝐜(f)=𝐀(f)𝐩.`$ (19)
In the above expression, while the vector $`𝐩`$ represents the unknowns consisting of the multipole coefficients of GWBs $`p_\mathrm{}m`$, the vector $`𝐜(f)`$ contains the correlation signals $`\stackrel{~}{C}_k(f)`$. Here, the frequency dependence of the function $`H(f)`$ has been already factorized and thereby the vector $`𝐩`$ only contains the information about the angular distribution. The matrix $`𝐀(f)`$ is the known quantity consisting of the multipole coefficients of each antenna pattern and the Wigner $`D`$ matrices (see Appendix C for an explicit example).
As we have already mentioned, the linear systems (19) become either over-determined or under-determined system. In most of our treatment in the low-frequency regime, the linear systems (19) tend to become over-determined, but this is not always correct depending on the amplitude of GWB spectrum $`H(f)`$ (see Sec.IV.2). In any case, the matrix $`𝐀`$ would not be a square matrix and the number of components of the vector $`𝐜`$ does not coincide with the one in the vector $`𝐩`$. In the over-determined case, while there is a hope to get a unique solution $`𝐩`$ to produce the gravitational-wave signal $`𝐜`$, it seems practically difficult due to the errors associated with the instrumental noise and/or numerical analysis. Hence, instead of pursuit of a rigorous solution, it would be better to focus on the issue how to get an approximate solution by a simple and systematic method.
In the case of our linear system, the approximate solution for the multipole coefficient $`𝐩_{\mathrm{approx}}`$ can be obtained from the least-squares method in the following form:
$$𝐩_{\mathrm{approx}}=𝐀^+𝐜.$$
(20)
The matrix $`𝐀^+`$ is called the pseudo-inverse matrix of Moore-Penrose type, whose explicit expression can be uniquely determined from the singular value decomposition (SVD) Press et al. (2002). According to the theorem of linear algebra, the matrix $`𝐀`$, whose number of rows is greater than or equal to its number of columns, can be generally written as $`𝐀=U^T\text{diag}[w_i]V`$. Here, the matrices $`U`$ and $`V`$ are orthonormal matrices which satisfy $`U^{}U=V^{}V=\mathrm{𝟏}`$, where the quantity with subscript represents the Hermite conjugate variable. The quantity $`\text{diag}[w_i]`$ represents an diagonal matrix with singular values $`w_i`$ with respect to the matrix $`𝐀`$. Then the pseudo-inverse matrix becomes
$$𝐀^+=V^T\text{diag}[w_i^1]U.$$
(21)
It should be stressed that the explicit form of the pseudo-inverse matrix $`𝐀^+`$ is characterized only by the angular dependence of antenna pattern functions.
In principle, the least-squares method by SVD can work well and all the accessible multipole moments of GWB would be obtained as long as the antenna pattern functions have the corresponding sensitivity to each detectable multipole moment. As we mentioned in Sec.II.2, however, the angular power of antenna pattern function depends sensitively on the frequency. In the low-frequency regime, the frequency dependence of the non-vanishing multipole moments appears at $`𝒪(\widehat{f}^2)`$ for $`\mathrm{}=0,\mathrm{\hspace{0.17em}\hspace{0.17em}2}`$ and $`4`$, and $`𝒪(\widehat{f}^3)`$ for $`\mathrm{}=1,\mathrm{\hspace{0.17em}\hspace{0.17em}3}`$ and $`5`$ (see Table 1). In this respect, by a naive application of the least-squares method, it is difficult to extract the information about $`\mathrm{}=`$odd modes of GWBs because the singular values of the matrix $`𝐀`$ are dominated by the lowest-order contribution of the antenna pattern functions.
For a practical and a reliable estimate of the odd multipoles in low-frequency regime, the least-squares method should be applied combining with the perturbative scheme described below. Let us recall from the expression (17) that the matrix $`𝐀`$ can be expanded as
$`𝐀=\widehat{f}^2𝐀^{(2)}+\widehat{f}^4𝐀^{(4)}+\mathrm{}`$ (22)
for a matrix consisting of the self-correlation signals $`_{AA}`$ and $`_{EE}`$,
$`𝐀=\widehat{f}^2𝐀^{(2)}+\widehat{f}^3𝐀^{(3)}+\widehat{f}^4𝐀^{(4)}+\mathrm{}`$ (23)
for a matrix consisting of the cross-correlation signal $`_{AE}`$, and
$`𝐀=\widehat{f}^3𝐀^{(3)}+\widehat{f}^4𝐀^{(4)}+\mathrm{}`$ (24)
for a matrix consisting of the cross-correlation signals $`_{AT}`$ and $`_{ET}`$. Then the resultant matrices $`𝐀^{(i)}`$ become independent of the frequency. The above perturbative expansion implies that the output signal $`𝐜`$ is also expanded in powers of $`\widehat{f}`$. Since the frequency dependence of the function $`H(f)`$ has been already factorized in equation (19), we have
$`𝐜(f)=\{\begin{array}{ccc}& \widehat{f}^2𝐜^{(2)}+\widehat{f}^4𝐜^{(4)}+\mathrm{},\hfill & \text{for AA-, EE-correlations}\hfill \\ & & \\ & \widehat{f}^2𝐜^{(2)}+\widehat{f}^3𝐜^{(3)}+\widehat{f}^4𝐜^{(4)}+\mathrm{},\hfill & \text{for AE-correlation}\hfill \\ & & \\ & \widehat{f}^3𝐜^{(3)}+\widehat{f}^4𝐜^{(4)}+\mathrm{},\hfill & \text{for AT-,ET-correlations}\hfill \end{array}`$ (30)
Substituting the terms (22)-(24) and (30) into the expression (19) and collecting the terms of each order of $`\widehat{f}`$, we have
$$𝐜^{(i)}=𝐀^{(i)}𝐩^{(i)},(i=2,\mathrm{\hspace{0.17em}\hspace{0.17em}3},\mathrm{\hspace{0.17em}\hspace{0.17em}4},\mathrm{})$$
(31)
where the subscript <sup>(i)</sup> means the quantity consisting of the order $`𝒪(\widehat{f}^i)`$ terms. The vector $`𝐩^{(i)}`$ represents the accessible multipole moments $`p_\mathrm{}m`$ to which the antenna pattern functions become sensitive in this order. For example, the vector $`𝐩^{(2)}`$ contains the $`\mathrm{}=0,\mathrm{\hspace{0.17em}2}`$ and $`4`$ modes, while the vector $`𝐩^{(3)}`$ have the multipole moments with $`\mathrm{}=1,\mathrm{\hspace{0.17em}3}`$ and $`5`$. Since the expression (31) has no explicit frequency dependence, we can reliably apply the least-squares solution by the SVD to reconstruct the odd modes of GWBs, as well as the even modes:
$$𝐩_{\mathrm{approx}}^{(i)}=[𝐀^{(i)}]^+𝐜^{(i)}.$$
(32)
Each step of the perturbative reconstruction scheme is summarized in Fig. 2 (see Appendix C in more details).
Finally, we note that the low-frequency reconstruction method presented here assumes the perturbative expansion form of the output signal $`𝐜`$. To determine the coefficients $`𝐜^{(i)}`$ in equation (30), one must know the frequency dependence of the vector $`𝐜`$ in the narrow bandwidth $`ff+\mathrm{\Delta }f`$, which can be achieved by analyzing the multi-frequency data. One important remark is that the signal-to-noise ratio in each reconstructed multipoles might be influenced by the sampling frequencies and/or the data analysis strategy. This point will be discussed later in Sec. IV.3.
### III.2 Reconstruction based on the time-series representation
So far, we have discussed the reconstruction method based on the harmonic-Fourier representation (14). The main advantage of the harmonic-Fourier representation is that it gives a simple algebraic equation suitable for applying the least-squares solution by SVD. However, the harmonic-Fourier representation implicitly assumes that the space interferometer orbits around the sun under keeping their configuration rigidly. In reality, rigid adiabatic treatment of the spacecraft motion is no longer valid due to the intrinsic variation of arm-length caused by the Keplerian motion of three space crafts, as well as the tidal perturbation by the gravitational force of solar system planets Cornish and Hellings (2003); Tinto et al. (2004); Shaddock (2004). In a rigorous sense, the time dependence of antenna pattern function cannot be described by the Euler rotation of antenna pattern function at rest frame (see footnote 2). Although the influence of arm-length variation is expected to be small in the low-frequency regime, alternative approach based on the other representation would be helpful to consider the GWB skymap beyond the low-frequency regime.
Here, we briefly discuss the reconstruction method based on the time-series representation (5). The time-series representation is mathematically equivalent to the harmonic-Fourier representation under the rigid adiabatic treatment, but in general, no additional assumption for space craft configuration is invoked to derive the expression (5), except for the premise that the time-dependence of the antenna pattern functions, i.e., the motion of each spacecraft, is well-known theoretically and/or observationally. Moreover, as we see below, the method based on the time-series representation allows us to reconstruct directly the skymap $`S_h`$, without going through intermediate variables, e.g. $`p_\mathrm{}m`$. In this sense, the time-series representation would be superior to the harmonic-Fourier representation for practical purpose, although there still remains the same problem as discussed in Appendix B concerning the degeneracy of the multipole coefficients.
In principle, the perturbative reconstruction scheme given in Sec.III.1 can be applicable to the map-making problem based on the time-series representation. To apply this, we discretize the integral expression (5) so as to reduce it to a matrix form like equation (19). For example, we discretize the celestial sphere $`(\theta ,\varphi )`$ into a regular $`N`$ meshes and also the continuous time-series into the regular $`M`$ grids. Then we have
$$\stackrel{~}{C}(t_i,f)=\underset{j=1}{\overset{N}{}}S_h(|f|,𝛀_j)^E(f,𝛀_j;t_i)\frac{\mathrm{\Delta }𝛀_j}{4\pi },(i=1,2,\mathrm{},M),$$
(33)
where $`\mathrm{\Delta }𝛀_i=\mathrm{sin}\theta _i\mathrm{\Delta }\theta _i\mathrm{\Delta }\varphi _i`$. Here, we have ignored the noise contribution and dropped the subscript <sub>IJ</sub>. For a reconstruction of low-frequency skymap, a large number of mesh and/or grid are not necessary. Typically, for the angular resolution up to $`\mathrm{}=5`$, it is sufficient to set $`M=16`$ and $`N=16`$ (see Sec.IV.3). The above discretization procedure is repeated for the five output data of the correlation signals. Then, simply following the same procedure as in the case of the harmonic-Fourier representation, the least-squares method by SVD is applied to get the luminosity distribution of GWBs. Note that the meanings of the matrix $`𝐀^{(i)}`$, the vectors $`𝐜^{(i)}`$ and $`𝐩^{(i)}`$ in the perturbative expansion are appropriately changed as follows. For the matrix $`𝐀^{(i)}`$, we have
$$𝐀^{(i)}=\frac{1}{4\pi }\left(\begin{array}{cccc}^{(i)}(𝛀_1,t_1)\mathrm{\Delta }𝛀_1& ^{(i)}(𝛀_2,t_1)\mathrm{\Delta }𝛀_2& \mathrm{}& ^{(i)}(𝛀_N,t_1)\mathrm{\Delta }𝛀_N\\ ^{(i)}(𝛀_1,t_2)\mathrm{\Delta }𝛀_1& ^{(i)}(𝛀_2,t_2)\mathrm{\Delta }𝛀_2& \mathrm{}& ^{(i)}(𝛀_N,t_2)\mathrm{\Delta }𝛀_N\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ ^{(i)}(𝛀_1,t_M)\mathrm{\Delta }𝛀_1& ^{(i)}(𝛀_2,t_M)\mathrm{\Delta }𝛀_2& \mathrm{}& ^{(i)}(𝛀_N,t_M)\mathrm{\Delta }𝛀_N\end{array}\right),(i=2,\mathrm{\hspace{0.17em}3},\mathrm{})$$
(34)
The corresponding vectors $`𝐜_{(i)}`$ and $`𝐩_{(i)}`$ becomes
$$𝐜^{(i)}=\left(\begin{array}{c}\stackrel{~}{C}^{(i)}(t_1)\\ \stackrel{~}{C}^{(i)}(t_2)\\ \mathrm{}\\ \stackrel{~}{C}^{(i)}(t_M)\end{array}\right),𝐩^{(i)}=\left(\begin{array}{c}S_h^{(i)}(\mathrm{\Omega }_1)\\ S_h^{(i)}(\mathrm{\Omega }_2)\\ \mathrm{}\\ S_h^{(i)}(\mathrm{\Omega }_N)\end{array}\right),$$
(35)
where $`\stackrel{~}{C}^{(i)}(t)`$ and $`S_h^{(i)}(\mathrm{\Omega })`$ the perturbative coefficients of $`\stackrel{~}{C}(t,f)`$ and $`S_h(|f|,\mathrm{\Omega })`$ in power of $`\widehat{f}^i`$, respectively. Using these expressions, the least-squares solution is constructed as $`𝐩_{\mathrm{approx}}=𝐩_{\mathrm{approx}}^{(2)}+𝐩_{\mathrm{approx}}^{(3)}+\mathrm{}`$, with a help of equation (32). Then, the resultant expression $`𝐩_{\mathrm{approx}}`$ directly gives a GWB skymap in the ecliptic coordinate, i.e., $`S_h(|f|,𝛀)`$, not the multipole coefficients.
## IV Demonstration: skymap of Galactic background
Perturbative reconstruction scheme presented in the previous section is applicable to a general map-making problem for any kind of GWB sources. In this section, to see how our general scheme works in practice, we demonstrate the reconstruction of a GWB skymap focusing on a specific source of anisotropic GWB. An interesting example for LISA detector is a confusion-noise background produced by the Galactic population of unresolved binaries. After describing a model of Galactic GWB in Sec.IV.1, the expected signals for time-modulation data of self- and cross-correlated TDIs are calculated in Sec.IV.2. Based on these, the detectable Fourier components for time-modulation signals are discussed evaluating the signal-to-noise ratios. In Sec.IV.3, a reconstruction of the GWB skymap is performed based on the harmonic-Fourier representation and the time-series representation. The resultant values of the multipole coefficients for Galactic GWBs are compared with those from the full-resolution skymap, taking account of the influence of the noises.
### IV.1 A model of Galactic GWB
For our interest of GWBs in the low-frequency regime with $`ff_{}9.54`$ mHz, it is reasonable to assume that the anisotropic GWB spectrum $`S_h(f,𝛀)`$ is separately treated as $`S_h(f,𝛀)=H(f)P(𝛀)`$, as we mentioned. The power spectral density $`H(f)`$ is approximated by a power-law function like $`H(f)=\mathrm{N}f^\alpha `$, whose amplitude will be explicitly given later. For illustrative purpose, we consider the simplest model of luminosity distribution $`P(𝛀)`$, in which the Galactic GWB is described by an incoherent superposition of gravitational waves produced by compact binaries whose spatial structure just traces the Galactic stellar distribution observed via infrared photometry. We use the fitting model of Galactic stellar distribution given in Binney et al. (1997), which consists of triaxial bulge and disk components (see also Seto (2004)). The explicit functional form of the density distribution $`\rho (𝐱)`$ written in the Galactic coordinate system becomes
$`\rho (𝐱)=\rho _{\mathrm{disk}}(𝐱)+\rho _{\mathrm{bulge}}(𝐱);`$ (36)
$`\rho _{\mathrm{bulge}}={\displaystyle \frac{\rho _0}{(1+a/a_0)^{1.8}}}e^{(a/a_\mathrm{m})^2},`$
$`\rho _{\mathrm{disk}}=\left({\displaystyle \frac{e^{|z|/z_0}}{z_0}}+\alpha {\displaystyle \frac{e^{|z|/z_1}}{z_1}}\right)R_se^{R/R_\mathrm{s}}`$
with the quantities $`a`$ and $`R`$ defined by $`a[x^2+(y/\eta )^2+(z/\zeta )^2]^{1/2}`$ and $`R(x^2+y^2)^{1/2}`$. Note that the $`z`$-axis is oriented to the north Galactic pole and the direction of the $`x`$-axis is $`20^{}`$ different from the Sun-center line. Here the parameters are given as follows: $`\rho _0=624`$, $`a_\mathrm{m}=1.9`$ kpc, $`a_0=100`$ pc, $`R_\mathrm{s}=2.5`$ kpc, $`z_0=210`$ pc, $`z_1=42`$ pc, $`\alpha =0.27`$, $`\eta =0.5`$ and $`\zeta =0.6`$. Provided the three-dimensional structure of stellar distribution $`\rho (𝐱)`$, the angular distribution $`P(𝛀)`$ is obtained by projecting it onto the sphere in observed frame, i.e., ecliptic coordinate:
$$P(𝛀)=C𝑑r\mathrm{\hspace{0.17em}\hspace{0.17em}4}\pi r^2\frac{\rho (r,𝛀)}{r^2},$$
(37)
where $`C`$ is a numerical constant normalized by $`𝑑𝛀P(𝛀)=1`$. The integration over $`r`$ is performed along a line-of-sight direction from a location of space interferometer to infinity.
In Fig.3, the projected intensity distribution of GWB is numerically obtained specifically in ecliptic coordinate and it is then transformed into Galactic coordinate, shown as the Hammer-Aitoff map. A strong intensity peak is found around the Galactic center and the disk-like structure can be clearly seen. Notice that while the result depicted in Fig.3 represents a full-resolution skymap, it cannot be attained from the perturbative reconstruction scheme in low-frequency regime. For the antenna pattern functions of cross-correlated TDI signals up to the order $`𝒪(\widehat{f}^3)`$, the detectable multipole moments of GWB are limited to $`\mathrm{}5`$. Moreover, in the low-frequency limit $`𝒪(\widehat{f}^2)`$, only the even modes $`\mathrm{}=0`$, $`2`$ and $`4`$ are measurable. Thus, the reconstructed skymap would be rather miserable compared to the full skymap which contains the higher multipoles $`\mathrm{}30`$ (see Fig.12 in Appendix). Taking account of these facts, in right panel of Fig.4, we plot the low-resolution skymap, which was obtained from the full-resolution skymap just dropping the higher multipole moments with $`\mathrm{}>5`$. In Appendix D, with a help of the Fortran package of spherical harmonic analysis (Adams and Swarztrauber (2003), see Appendix D), the numerical values of the multipole coefficients $`p_\mathrm{}m`$ up to $`\mathrm{}=5`$ are computed and are summarized in Table 3. Also in the left panel, the odd modes are further subtracted and the remaining multipole moments are only $`\mathrm{}=0`$, $`2`$ and $`4`$. As a result, the fine structure around the bulge and the disk components is coarse-grained and the intensity of the GWB diminishes. It also shows some fake patterns with negative intensity. Nevertheless, one can clearly see the anisotropic structure of GWB , which is mainly contributed from the gravitational-wave sources around the Galactic disk. With a perturbative reconstruction of low-frequency up to $`𝒪(\widehat{f}^3)`$, one can roughly infer that the main sources of Galactic GWB comes from the Galactic center.
### IV.2 Time-modulation signals and signal-to-noise ratio
Given the intensity distribution of GWB, one can calculate the cross-correlation signals observed via LISA, which are inherently time-dependent due to the orbital motion of the LISA detector. The effect of the annual modulation of the Galactic binary confusion noise on the LISA data analysis had been previously studied in Seto (2004) in the low-frequency limit $`\widehat{f}1`$. Recently, Monte Carlo simulations of Galactic GWB were carried out by several groups and the annual modulation of GWB intensity has been confirmed in a realistic setup with specific detector output Benacquista et al. (2004); Edlund et al. (2005); Timpano et al. (2005).
Owing to the expression (5), the time-modulation signals $`\stackrel{~}{C}(t,f)`$ neglecting the instrumental noises are computed for optimal TDIs at the frequency $`\widehat{f}=0.1`$, i.e., $`f1\mathrm{m}\mathrm{H}\mathrm{z}`$ using the full expression of antenna pattern function (7). The results are then plotted as function of orbital phase. In Fig. 5, six outputs of the self- and cross-correlation signals normalized to its time-averaged value, $`\stackrel{~}{C}(t,f)/|\stackrel{~}{C}_0(f)|`$ are shown. The solid and dashed lines denote the real and imaginary parts of the correlation signals, respectively.
The time modulation of these outputs basically reflects the spatial structure of GWB. As LISA orbits around the Sun, the direction normal to LISA’s detector plane, which is the direction sensitive to the gravitational waves, sweeps across the Galactic plane. Since the response of the LISA detector to the gravitational waves along the $`\widehat{𝛀}`$ direction give the same response to the waves along the $`\widehat{𝛀}`$ direction (see Eq.(9) and the brief comment there), the time-modulation signal is expected to have a bimodal structure, like $`AA`$\- and $`EE`$-correlations. However, the actual modulation signals are more complicated, depending on the angular sensitivity of their antenna pattern functions, as well as the observed frequency. Further, the cross-correlation data can be generally complex variables, whose behaviors are different between real- and imaginary-parts.
Notice that the time-modulation signals presented in Fig. 5 are the results in an idealistic situation free from the noise contribution. In presence of the random noise, some of the correlation data which contain gravitational wave signals are overwhelmed by noise, which cannot be used for the reconstruction of GWB skymap. Hence, one must consider the signal-to-noise ratio (SNR) to discriminate available correlation data. To evaluate this, the output signals depicted in Fig.5 are first transformed to its Fourier counterpart, $`\stackrel{~}{C}_k(f)`$. The resultant Fourier components for each signal are then compared to the noise contributions. The SNR for each component is expressed in the following form (Seto and Cooray (2004), see also Giampieri and Polnarev (1997a); Ungarelli and Vecchio (2001b)):
$$\left(\frac{S}{N}\right)_k=\sqrt{2\mathrm{\Delta }fT}\frac{|\stackrel{~}{C}_k|}{N_k}.$$
(38)
Here we set the observational time to $`T=10^8`$sec and the bandwidth to $`\mathrm{\Delta }f=10^3`$Hz. The quantity $`N_k`$ represents the noise contribution. The important remark is that the noise contribution comes from not only the detector noise but also the randomness of the signal itself. The details of the analytic expression for $`N_k`$ will be given elsewhere Kudoh and Taruya (2005b). Here, as a crude estimate, we evaluate the noise contribution $`N_k`$ as
$$N_k=\sqrt{2\mathrm{\Delta }fT}S_n^{II},(I=A,E,T)$$
(39)
for $`k=0`$ component of self-correlation signals and
$$N_k=\sqrt{\text{max}(\stackrel{~}{C}_{II,0}\stackrel{~}{C}_{JJ,0},\stackrel{~}{C}_{II,0}S_n^{JJ},\stackrel{~}{C}_{JJ,0}S_n^{II},S_n^{II}S_n^{JJ})},(I,J=A,E,T)$$
(40)
for cross-correlation signals and $`k0`$ component of self-correlation signals. To estimate SNR, one further needs the spectral density for instrumental noise, i.e., $`S_n^{\mathrm{AA}}(f)`$, $`S_n^{\mathrm{EE}}(f)`$ and $`S_n^{\mathrm{TT}}(f)`$. We use the expressions given in equation (58) of Paper I (see also Cornish (2002b); Prince et al. (2002)) <sup>6</sup><sup>6</sup>6In equation (58) of Paper I, there are some typos in the numerical values of $`S_{\mathrm{shot}}(f)`$ and $`S_{\mathrm{accel}}(f)`$. In the present paper, we adopt $`S_{\mathrm{shot}}(f)=1.6\times 10^{41}`$Hz<sup>-1</sup> and $`S_{\mathrm{accel}}(f)=2.31\times 10^{41}(\text{mHz}/f)^4`$Hz<sup>-1</sup> according to Refs.Bender and et. al., (1998); Cornish (2002b).. Note that all the cross-correlated noise spectra such as $`S_n^{\mathrm{AE}}(f)`$ and $`S_n^{\mathrm{AT}}(f)`$ are exactly canceled.
In Fig.6, the SNRs for six output data are evaluated and are shown in the histogram as function of Fourier component, $`k`$. In these panels, thick-dotted lines show the detection limit of $`(S/N)_k=5`$, while the thin-dotted lines mean $`(S/N)_k=1`$. When evaluating the SNR, we specifically consider the two cases:
realistic case in which the rms amplitude of GWB spectrum is given by $`S_h^{1/2}=5\times 10^{19}`$Hz<sup>-1/2</sup> at the frequency $`f=1`$mHz <sup>7</sup><sup>7</sup>7The amplitude of GWB spectrum might be reduced by a factor of $`25`$ according to the recent numerical simulations Edlund et al. (2005); Timpano et al. (2005). In case A, however, the noise contributions in SNR (38) used for reconstruction analysis are basically determined by the $`k=0`$ components of self-correlation signals, i.e., $`N_k=\sqrt{\stackrel{~}{C}_{II,0}\stackrel{~}{C}_{JJ,0}}`$. Hence, a slight change of the GWB amplitude would not alter the final results..
optimistic case in which the rms amplitude of GWB spectrum is ten times larger than that in the realistic case, i.e., $`S_h^{1/2}=5\times 10^{18}`$Hz<sup>-1/2</sup> at $`f=1`$mHz.
The results of SNR are then shown in solid (case A) and dashed lines (case B), respectively.
As anticipated from the sensitivity of antenna pattern function in Table 1, the SNR for $`TT`$-correlation is much less than unity. Even in the optimistic case, the SNR is about ten times smaller than unity and thus the $`TT`$-correlation data cannot be used for reconstruction of the GWB skymap. Apart from this, the SNRs for the self-correlation signals $`AA`$ and $`EE`$ as well as for the cross-correlation signal $`AE`$ are generally good compared to the cross-correlation data $`AT`$ and $`ET`$. With sufficient higher SNR of $`(S/N)_k5`$, the available Fourier components of $`AA`$-, $`EE`$\- and $`AE`$-correlations become $`k=2,1,\mathrm{\hspace{0.17em}0},+1,+2`$ in both realistic and optimistic cases. This is consistent with the previous estimates Seto and Cooray (2004). <sup>8</sup><sup>8</sup>8The precise numerical values of the SNR for $`AA`$\- and $`EE`$-correlations slightly differ from Ref. Seto and Cooray (2004). This is mainly because the Euler rotation angles for the LISA orbital motion given in Sec.II.1 are different from those in Seto and Cooray (2004). On the other hand, with a large amplitude of GWB spectrum (case B), only the $`k=0`$ component is accessible in the signal combinations of $`AT`$ and $`ET`$. This is mainly due to the fact that the sensitivity of the $`T`$-variable is poor at low-frequency and thereby the noise contribution $`N_k`$ becomes $`\{\stackrel{~}{C}_{AA,0}S_n^{\mathrm{TT}}\}^{1/2}`$ or $`\{\stackrel{~}{C}_{EE,0}S_n^{\mathrm{TT}}\}^{1/2}`$, which is much larger than $`\{\stackrel{~}{C}_{AA,0}\stackrel{~}{C}_{TT,0}\}^{1/2}`$ or $`\{\stackrel{~}{C}_{EE,0}\stackrel{~}{C}_{TT,0}\}^{1/2}`$.
Thus, the available Fourier components of correlation data used for the reconstruction of the skymap would be severely restricted in practice. Under such restricted situation, the deconvolution problem of the linear system (14) tends to be under-determined. Nevertheless, as it will be shown below, one can determine the $`\mathrm{}=0`$, $`2`$ and $`4`$ modes of multipole coefficients of the Galactic GWB with sufficiently small errors. In addition, with the $`k=0`$ components of $`AT`$\- and $`ET`$-correlations, the odd moments $`\mathrm{}=1`$ and $`3`$ can be recovered.
### IV.3 Reconstruction of a skymap
Keeping the remarks on the SNR estimation in previous subsection in mind, we now proceed to a reconstruction analysis and test the validity of perturbative reconstruction method presented in Sec. III. For this purpose, in addition to the analysis in the under-constrained cases (case A and B) mentioned above, we also consider the over-constrained case as an illustrative example.
#### IV.3.1 Over-determined case
(I) Harmonic-Fourier representation Let us first focus on a very idealistic situation that the noise contributions are entirely neglected. In such a case, all the components of self- and cross-correlation data $`\stackrel{~}{C}_k`$ are available to the reconstruction analysis. In practice, however, it is sufficient to consider some restricted components among all available data. Here, to make a skymap with multipoles of $`\mathrm{}5`$, we use the $`k=2+2`$ components of self-correlation signals $`AA`$ and $`EE`$, the $`k=4+4`$ components of cross-correlation signals $`AT`$ and $`ET`$, and $`k=8+8`$ components of $`AE`$-signal (see Table 2). Even in this case, the linear system (14) is still over-determined and the least-squares approximation by SVD is potentially powerful to obtain the multipole coefficients of anisotropic GWB. The procedure of reconstruction analysis is the same one as presented in Fig.2. For the output signals of harmonic-Fourier representation, we use the data $`\stackrel{~}{C}_k`$ observed at the frequencies $`\widehat{f}=0.05`$ and $`0.15`$, in addition to the data for our interest at $`\widehat{f}=0.1`$. Collecting these multi-frequency data, the perturbative expansion form of the vector $`𝐜(f)`$ is specified up to the third order in $`\widehat{f}`$ and the coefficients $`𝐜^{(i)}`$ are determined (Appendix C).
In Fig.7, the reconstructed results of multipole coefficients $`p_\mathrm{}m`$ are converted to the projected intensity distribution (9) and are shown as Hammer-Aitoff map in Galactic coordinate. Left panel shows the skymap reconstructed from the lowest-order signals of self- and cross-correlation data $`𝐜^{(2)}`$, which only includes the $`\mathrm{}=0`$, $`2`$ and $`4`$ modes, while right panel is the result taking account of the leading-order correction $`𝐜^{(3)}`$. Comparing those with the expected skymap shown in Fig.4, the reconstruction seems almost perfect. In Fig. 8, numerical values of the reconstructed multipole coefficients are compared with the true values listed in Table 3. The agreement between the reconstruction results (open circles) and the true values (crosses) is quite good and the fractional errors are well within a few percent except for $`p_{20}`$ and $`p_{50}`$. A remarkable fact is that the monopole and the quadrupole values can be reproduced reasonably well despite the presence of the degeneracy mentioned in Appendix.B. This is just an accidental result. The (small) discrepancy in the “recovered” $`p_{20}`$ can be ascribed to the expression (56). After all, the least-squares method by SVD provides a robust reconstruction method when the linear system (14) becomes over-determined.
(II) Time-series representation The successful reconstruction of the GWB skymap can also be achieved by the alternative approach based on the time-series representation (5). Following the procedure presented in Sec.III.2, the intensity skymap of the GWB is directly obtained and the results taking account of the lowest-order and the leading-order contributions to the antenna pattern functions are shown in left and right panels in Fig.9, respectively. Here, to create the discretized data set (33), the number of grid and/or mesh was specified as $`M=16`$ in time and $`N=17\times 32`$ in spherical coordinate. The time-series data of antenna pattern functions were numerically generated based on the full analytic expressions given in Sec.II.2 under assuming that the arm-length of the three space crafts are rigidly kept fixed. The reconstructed skymap reasonably agrees with Fig.7 as well as Fig.4. Although the situation considered here is very idealistic and thus the results in Fig.9 should be regarded as just a preliminary one, one expects that the methodology based on the time-series representation is potentially powerful even when the rigid adiabatic treatment of space craft motion becomes inadequate. To discuss its effectiveness, a further investigation is needed. The details of the analysis including the effects of arm-length variation will be presented elsewhere.
#### IV.3.2 Under-determined case
Turning to focus on the analysis based on the harmonic-Fourier representation, we next consider the under-constrained case in which the number of available Fourier components is restricted due to the noises (case A and B discussed in Sec.IV.2). Fig.10 shows the reconstructed images of GWB intensity map in Galactic coordinate, free from the noises but \] restricting the number of Fourier components according to Table 2. The top panel is the intensity map obtained from the lowest-order signals $`𝐜^{(2)}`$. Since the accessible Fourier components are basically the same in the lowest-order analysis, the same results are obtained between case A and B. On the other hand, the bottom panels of Fig.10 represent the skymap reconstructed from both $`𝐜^{(2)}`$ and $`𝐜^{(3)}`$, which indicate that the different images of GWB skymap are obtained depending on the available number of Fourier components (or the amplitude of GWB spectrum); case A (left) and case B (right). In Fig.11, numerical values of the reconstructed multipoles $`p_\mathrm{}m`$ in ecliptic coordinate are summarized together with the statistical errors. The statistical errors were roughly estimated according to the discussion in Appendix E based on the SNR \[Eq. (38)\].
It turns out that the case A fails to reconstruct all the dipole moments, while they can be somehow reproduced in the case B. This is because the cross-correlation data $`AT`$ and $`ET`$ which include the information about $`\mathrm{}=1`$ modes were not used in the reconstruction analysis in case A. Although the $`\mathrm{}=5`$ modes of multipole coefficients were not reproduced well in both cases, their contributions to the intensity distribution are not large. Hence, the reconstructed GWB skymap in case B roughly matches the expected intensity map in Fig.4 and the visual impression becomes better than case A. This readily implies that large amplitude of the GWB spectrum is required for a correct reconstruction of a skymap, in practice. However, it should be emphasized that the present reconstruction technique can work well in the under-determined cases. Even in the realistic situation with smaller amplitude of GWB spectrum (case A), the reconstructed skymap including the multipoles $`\mathrm{}<5`$ still shows a disk-like structure, which may provide an important clue to discriminate between the Galactic and the extragalactic GWBs.
## V Conclusion and Discussion
In this paper, we have presented a perturbative reconstruction method to make a skymap of GWB observed via space interferometer. The orbital motion of the detector makes the output signals of GWB time-dependent due to the presence of anisotropies of GWB. Since the output signals of GWB are obtained through an all-sky integral of primary signals convolving with an antenna pattern function of gravitational-wave detectors, the time dependence of output data can be used to reconstruct the luminosity distribution of GWBs under full knowledge of detector’s antenna pattern functions. Focusing on the low-frequency regime, we have explicitly given a non-parametric reconstruction method based on both the harmonic-Fourier and the time-series representation. With a help of low-frequency expansion of the antenna pattern functions, the least-squares approximation by SVD enables us to obtain the multipole coefficients of GWB or direct intensity map even when the system becomes over-determined or under-determined. For illustrative purpose, the reconstruction analysis of the GWB skymap has been demonstrated for the confusion-noise background of Galactic binaries around the low-frequency $`f=1`$mHz. It then turned out that the space interferometer LISA free from the noises is capable of making a skymap of Galactic GWB with angular resolution up to the multipoles, $`\mathrm{}=5`$. For more realistic case based on the estimation of signal-to-noise ratios, the system tends to become under-determined and the number of Fourier components used in reconstruction analysis would be severely restricted. Nevertheless, the resultant skymap still contains the information of the multipoles up to $`\mathrm{}<5`$, from which one can infer that the main sources of GWB come from the Galactic center.
Although the present paper focuses on the reconstruction of the GWB skymap from the LISA, the methodology discussed in Sec.III is quite general and is also applicable to the reconstruction of low-frequency skymap obtained from the ground detectors as well as the other space interferometers, provided their antenna pattern functions. Despite a wide applicability of the present method, however, accessible multipole moments of low-frequency GWB are generally restricted to lower multipoles due to the properties of antenna pattern functions. This would be generally true in any gravitational-wave detectors at the frequencies $`f<f_{}=c/(2\pi L)`$, where $`L`$ is arm-length of single detector or separation between the two detectors. Hence, with the limited angular resolution, only the present methodology may not give a powerful constraint on the luminosity distribution of GWBs. In this respect, we need to combine the other techniques such as the parametric reconstruction method, in which we specifically assume an explicit functional form of the luminosity distribution characterized by the finite number of model parameters and determine them through the likelihood analysis. The data analysis strategy to give a tight constraint on the luminosity distribution is definitely a very important issue and it must be considered urgently. Nevertheless, we note that detectable multipole moments depend practically on the signal-to-noise ratios. Since the signal-to-noise ratios are determined both from the amplitude of GWB and detector’s intrinsic noises, a further feasibility study is necessary in order to clarify the detectable multipole moments correctly. With improved data analysis strategy, it might be even possible that the angular resolution of GWB skymap becomes better than that of LISA considered in the present paper.
Another important issue on the map-making problem is to consider the reconstruction of skymap beyond the low-frequency regime, where the antenna pattern functions give a complicated response to the anisotropic GWBs and thereby the angular resolution of GWB map can be improved Kudoh and Taruya (2005a). Since the low-frequency expansion cannot be used there, a new reconstruction technique should be devised to extract the information of anisotropic GWBs. Further, in the case of LISA, the effect of arm-length variation becomes important and the rigid adiabatic treatment of the detector response cannot be validated. Hence, as emphasized in Sec.III.2, it would become essential to consider the reconstruction method based on the time-series representation, in which no additional assumptions for spacecraft configuration and/or motion are needed to compute the antenna pattern functions. The analysis concerning this issue will be presented elsewhere.
###### Acknowledgements.
We would like to thank N. Seto for valuable comments on the estimation of signal-to-noise ratios. We also thank Y. Himemoto and T. Hiramatsu for discussions and comments, T. Takiwaki and K. Yahata for the technical support to plot the skymap. The work of H.K. is supported by the Grant-in-Aid for Scientific Research of Japan Society for Promotion of Science.
## Appendix A Multipole coefficients of antenna pattern functions in low-frequency regime
In this appendix, based on the analytic expression (21) of paper I, the multipole moments for antenna pattern functions of optimal TDIs are calculated at detector’s rest frame. Using the low-frequency approximation $`\widehat{f}1`$, we present the perturbative expressions up to the forth order in $`\widehat{f}`$.
To evaluate the antenna pattern functions, we must first specify the directional unit vectors $`𝐚`$, $`𝐛`$ and $`𝐜`$, which connect respective spacecrafts. Here, we specifically choose (Fig. 1 and Eq.(28) of Paper I):
$`𝐚={\displaystyle \frac{\sqrt{3}}{2}}𝐱+{\displaystyle \frac{1}{2}}𝐲,𝐛=𝐲,𝐜={\displaystyle \frac{\sqrt{3}}{2}}𝐱+{\displaystyle \frac{1}{2}}𝐲,`$
where the vectors $`𝐱`$ and $`𝐲`$ respectively denote the unit vectors parallel to the $`x`$\- and $`y`$-axes in detector’s rest frame (see Fig.1 of Paper I). Then, based on the expression (15), the antenna pattern functions of the optimal TDIs, $`_{IJ}`$ $`(I,J=A,E,T)`$ are analytically computed at detector’s rest frame. Their multipole moments become
$$a_\mathrm{}m(\widehat{f})=_0^\pi 𝑑\theta _0^{2\pi }𝑑\varphi \mathrm{sin}\theta Y_\mathrm{}m^{}(\theta ,\varphi )(\widehat{f},\theta ,\varphi ),$$
(41)
which are expressed as a function of normalized frequency $`\widehat{f}=f/f_{}`$. Here, for definiteness, we write down the explicit form of the harmonic functions $`Y_\mathrm{}m`$:
$`Y_{\mathrm{}}^m(\theta ,\varphi )=\sqrt{{\displaystyle \frac{2\mathrm{}+1}{4\pi }}{\displaystyle \frac{(\mathrm{}m)!}{(\mathrm{}+m)!}}}P_{\mathrm{}}^m(\mathrm{cos}\theta )e^{im\varphi }.`$ (42)
The analytic expressions for the multipole moments $`a_\mathrm{}m`$ are generally intractable due to the complicated form of the antenna pattern functions. In the low-frequency regime, however, the perturbative treatment regarding $`\widehat{f}`$ as a small expansion parameter is applied to derive an analytic expression of multipole moments. Below, we summarize the perturbation results up to the fourth order in $`\widehat{f}`$:
$`_{AA}:`$ $`a_{00}={\displaystyle \frac{2\sqrt{\pi }}{5}}\widehat{f}^2{\displaystyle \frac{211\sqrt{\pi }}{1260}}\widehat{f}^4,a_{20}={\displaystyle \frac{4}{7}}\sqrt{{\displaystyle \frac{\pi }{5}}}\widehat{f}^2{\displaystyle \frac{16}{63}}\sqrt{{\displaystyle \frac{\pi }{5}}}\widehat{f}^4a_{22}={\displaystyle \frac{1i\sqrt{3}}{756}}\sqrt{{\displaystyle \frac{15\pi }{2}}}\widehat{f}^4,`$ (45)
$`a_{40}={\displaystyle \frac{\sqrt{\pi }}{105}}\widehat{f}^2+{\displaystyle \frac{13\sqrt{\pi }}{9240}}\widehat{f}^4,a_{42}={\displaystyle \frac{13(1i\sqrt{3})}{5544}}\sqrt{{\displaystyle \frac{\pi }{10}}}\widehat{f}^4,a_{44}={\displaystyle \frac{1+i\sqrt{3}}{6}}\sqrt{{\displaystyle \frac{\pi }{70}}}\widehat{f}^2+{\displaystyle \frac{13(1+i\sqrt{3})}{792}}\sqrt{{\displaystyle \frac{5\pi }{14}}}\widehat{f}^4`$
$`a_{60}={\displaystyle \frac{1}{1848}}\sqrt{{\displaystyle \frac{\pi }{13}}}\widehat{f}^4,a_{62}={\displaystyle \frac{(1i\sqrt{3})}{4752}}\sqrt{{\displaystyle \frac{3\pi }{455}}}\widehat{f}^4,a_{64}={\displaystyle \frac{1+i\sqrt{3}}{396}}\sqrt{{\displaystyle \frac{\pi }{182}}}\widehat{f}^4,`$
for self-correlated antenna pattern, $`_{AA}`$. Note that the multipole moments of $`_{EE}`$ are related to those of $`_{AA}`$ through the relations, $`a_\mathrm{}m^{\mathrm{EE}}=a_\mathrm{}m^{\mathrm{AA}}`$ for $`m=0,\pm 6,\pm 12,\pm 18,\mathrm{}`$ and $`a_\mathrm{}m^{\mathrm{EE}}=a_\mathrm{}m^{\mathrm{AA}}`$ for $`m=\pm 2,\pm 4,\pm 8,\mathrm{}`$ (see Paper I). The multipole moments of cross-correlation signals are
$`_{AE}:`$ $`a_{22}={\displaystyle \frac{3+i\sqrt{3}}{756}}\sqrt{{\displaystyle \frac{5\pi }{2}}}\widehat{f}^4,a_{33}={\displaystyle \frac{1}{18}}\sqrt{{\displaystyle \frac{7\pi }{15}}}\widehat{f}^3,a_{42}={\displaystyle \frac{13(\sqrt{3}+i)}{5544}}\sqrt{{\displaystyle \frac{\pi }{10}}}\widehat{f}^4,`$ (48)
$`a_{44}={\displaystyle \frac{\sqrt{3}i}{6}}\sqrt{{\displaystyle \frac{\pi }{70}}}\widehat{f}^2{\displaystyle \frac{13(\sqrt{3}i)}{792}}\sqrt{{\displaystyle \frac{5\pi }{14}}}\widehat{f}^4,a_{53}={\displaystyle \frac{1}{18}}\sqrt{{\displaystyle \frac{\pi }{1155}}}\widehat{f}^3,a_{62}={\displaystyle \frac{\sqrt{3}+i}{4752}}\sqrt{{\displaystyle \frac{3\pi }{455}}}\widehat{f}^4,`$
$`a_{64}={\displaystyle \frac{\sqrt{3}i}{396}}\sqrt{{\displaystyle \frac{\pi }{182}}}\widehat{f}^4,`$
$`_{AT}:`$ $`a_{11}={\displaystyle \frac{(1+i\sqrt{3})\sqrt{\pi }}{168}}\widehat{f}^3,a_{22}={\displaystyle \frac{\sqrt{3}+i}{864}}\sqrt{{\displaystyle \frac{3\pi }{5}}}\widehat{f}^4,a_{31}={\displaystyle \frac{1+i\sqrt{3}}{72}}\sqrt{{\displaystyle \frac{\pi }{14}}}\widehat{f}^3,`$ (51)
$`a_{42}={\displaystyle \frac{\sqrt{3}+i}{3168}}\sqrt{{\displaystyle \frac{\pi }{5}}}\widehat{f}^4,a_{44}={\displaystyle \frac{17(\sqrt{3}i)}{3168}}\sqrt{{\displaystyle \frac{\pi }{35}}}\widehat{f}^4,a_{51}={\displaystyle \frac{1+i\sqrt{3}}{1008}}\sqrt{{\displaystyle \frac{\pi }{55}}}\widehat{f}^3,`$
$`a_{55}={\displaystyle \frac{1i\sqrt{3}}{72}}\sqrt{{\displaystyle \frac{3\pi }{154}}}\widehat{f}^3,a_{62}={\displaystyle \frac{\sqrt{3}+i}{4752}}\sqrt{{\displaystyle \frac{3\pi }{910}}}\widehat{f}^4,a_{64}={\displaystyle \frac{\sqrt{3}i}{3168}}\sqrt{{\displaystyle \frac{\pi }{91}}}\widehat{f}^4.`$
The multipole moments of antenna pattern function $`_{ET}`$ are also obtained from those of $`_{AT}`$ using the relation, $`a_\mathrm{}m^{ET}=(i/\sqrt{3})\mathrm{tan}\left(m\pi /3\right)a_\mathrm{}m^{AT}`$, as shown in Paper I.
Finally, we also list the multipole moments of self-correlated antenna pattern $`_{TT}`$, which are all higher-order contribution with $`𝒪(\widehat{f}^4)`$:
$`_{TT}:`$ $`a_{00}={\displaystyle \frac{\sqrt{\pi }}{504}}\widehat{f}^4,a_{40}={\displaystyle \frac{\sqrt{\pi }}{1584}}\widehat{f}^4,a_{60}={\displaystyle \frac{1}{11088}}\sqrt{{\displaystyle \frac{\pi }{13}}}\widehat{f}^4,a_{66}={\displaystyle \frac{1}{48}}\sqrt{{\displaystyle \frac{\pi }{3003}}}\widehat{f}^4.`$ (52)
## Appendix B Remarks on the degeneracy between monopole and quadrupole moments for LISA measurement of GWB anisotropy
The reconstruction method based on the harmonic-Fourier representation (14) presented in Sec.III.1 has been first discussed by Cornish Cornish (2001, 2002a). Later, Seto & Cooray Seto and Cooray (2004) considered the reconstruction of a skymap in the low-frequency limit using the optimal set of TDI signals $`A`$ and $`E`$. In this case, the accessible multipole moments $`p_\mathrm{}m`$ are $`\mathrm{}=0,\mathrm{\hspace{0.17em}2}`$ and $`4`$ (see Table 1). Seto & Cooray explicitly wrote down the expressions for linear equation (14) in the case of the self-correlation signals (i.e., AA-, EE-correlations) and showed that the output data with sufficient statistical significance are only $`\stackrel{~}{C}_0`$, $`\stackrel{~}{C}_{\pm 1}`$ and $`\stackrel{~}{C}_{\pm 2}`$. Further, they found that the multipole coefficients $`p_{00}`$ and $`p_{20}`$ cannot be separately determined due to the degeneracy associated with a specific combination between the Wigner $`D`$ matrices and the antenna pattern functions.
Here, we explicitly point out the origin of this degeneracy and discuss its influences on the reconstruction of GWB skymap. Using the multipole coefficients of antenna patterns in Appendix A, the lowest-order contribution to the linear system (14) for AA- and EE-correlations becomes
$`\stackrel{~}{C}_{AA,0}^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{4\sqrt{\pi }}}\left\{{\displaystyle \frac{2}{5}}p_{00}{\displaystyle \frac{1}{14\sqrt{5}}}p_{20}{\displaystyle \frac{37}{13440}}p_{40}{\displaystyle \frac{27}{512\sqrt{70}}}{\displaystyle \underset{m=4,4}{}}b_{4m}^{\left(\mathrm{AA}\right)}p_{4m}\right\},`$ (53)
$`\stackrel{~}{C}_{EE,0}^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{4\sqrt{\pi }}}\left\{{\displaystyle \frac{2}{5}}p_{00}{\displaystyle \frac{1}{14\sqrt{5}}}p_{20}{\displaystyle \frac{37}{13440}}p_{40}+{\displaystyle \frac{27}{512\sqrt{70}}}{\displaystyle \underset{m=4,4}{}}b_{4m}^{\left(\mathrm{AA}\right)}p_{4m}\right\}.`$ (54)
with $`b_{44}^{\left(\mathrm{AA}\right)}=1i\sqrt{3}=[b_{4,4}^{\left(\mathrm{AA}\right)}]^{}`$. Here, we only show the relevant components of $`\stackrel{~}{C}_k`$ which contains the multipole coefficients $`p_{00}`$ and $`p_{20}`$. The above expressions include the multipole coefficients of $`\mathrm{}=4`$, which are all determined separately from the lowest-order contribution of $`AE`$-correlation, $`\stackrel{~}{C}_{AE,k}^{(2)}`$. Thus, apart from the octupole moments, equations (53) and (54) clearly show the presence of degeneracy between the remaining multipole coefficients, $`p_{00}`$ and $`p_{20}`$.
In a language of the least-squares method by SVD, this degeneracy implies that the sub-system in the matrix equation (31) containing the coefficients $`p_{00}`$ and $`p_{20}`$ becomes under-determined. In this case, the least-squares method by SVD cannot correctly produce the approximate solutions for $`p_{00}`$ and $`p_{20}`$, although it still provides some “approximate” solution. From equations (53) and (54), the only meaningful equation for $`p_{00}`$ and $`p_{20}`$ is now reduced to
$$\stackrel{~}{C}=\frac{1}{4\sqrt{\pi }}\left(\frac{2}{5}p_{00}\frac{1}{14\sqrt{5}}p_{20}\right),$$
(55)
where the numerical constant $`\stackrel{~}{C}`$ represents a collection of the irrelevant terms of $`\mathrm{}=4`$ modes, which are separately determined from the linear system in $`AE`$-correlation. If one naively applies the least-squares method to the above system, a very tight relation is obtained:
$`p_{00}={\displaystyle \frac{7840\sqrt{\pi }}{789}}\stackrel{~}{C},p_{20}={\displaystyle \frac{280\sqrt{5\pi }}{789}}\stackrel{~}{C}.`$ (56)
The presence of degenerate coefficients may be a big obstacle in constructing the skymap as well as in determining the normalization factor of GWB spectrum. In principle, this degeneracy can be broken when we consider the higher-order terms of $`𝒪(\widehat{f}^4)`$. However, these terms are generally small and irrelevant for the reconstruction analysis in the low-frequency regime. In this sense, the reconstruction of low-frequency skymap is, in nature, incomplete and the other additional information for monopole or quadrupole moment is required to make a full skymap. Nevertheless, it should be emphasized that with a high signal-to-noise ratio, the other remaining multipole coefficients can be all determined by the least-squares solution by SVD, independently of the above degeneracy. Moreover, in cases with $`p_{00}p_{20}`$, which is usually satisfied, the least-squares solution (56) provides a modest estimate of the degenerate coefficients $`p_{00}`$ and $`p_{20}`$ because of the hierarchy of the coefficients in Eq. (56). This point has been explicitly demonstrated in Sec.IV.
## Appendix C Data sets and least-squares method
In this appendix, we describe in more details how to obtain the multipole moments of GWBs from many data streams by applying the least-squares method. Here we specifically focus on the situation considered in Sec. IV.3.1. That is, the data sets that we use are (i) the $`k=2+2`$ components of self-correlation signals $`AA`$ and $`EE`$, (ii) $`k=8+8`$ components of $`AE`$-signal, and (iii) the $`k=4+4`$ components of cross-correlation signals $`AT`$ and $`ET`$. According to the general strategy given in Sec. III.1 (see Eq. (31)), we first combine the leading data streams of (i) and (ii) which correspond to $`i=2`$ in Eq. (30). The combined data sets consist of the following matrices.
$$\left(\begin{array}{c}\stackrel{~}{C}_{+2,AA}^{(2)}\\ \stackrel{~}{C}_{+1,AA}^{(2)}\\ \stackrel{~}{C}_{1,AA}^{(2)}\\ \stackrel{~}{C}_{2,AA}^{(2)}\\ \stackrel{~}{C}_{+2,EE}^{(2)}\\ \mathrm{}\\ \stackrel{~}{C}_{2,EE}^{(2)}\end{array}\right)=\left(\begin{array}{ccccccccccccc}0& 0& 0& & & & & & & & & & \\ 0& 0& & 0& & & & & & & & & \\ 0& & 0& 0& & & & & & & & & \\ & 0& 0& 0& & & & & & & & & \\ 0& 0& 0& & & & & & & & & & \\ 0& 0& & 0& & & & & & & & & \\ 0& & 0& 0& & & & & & & & & \\ & 0& 0& 0& & & & & & & & & \end{array}\right)\left(\begin{array}{c}p_{22}\\ p_{21}\\ p_{2,1}\\ p_{2,2}\end{array}\right),\left(\begin{array}{c}\stackrel{~}{C}_{+8,AE}^{(2)}\\ \stackrel{~}{C}_{+7,AE}^{(2)}\\ \mathrm{}\\ \stackrel{~}{C}_{0,AE}^{(2)}\\ \mathrm{}\\ \stackrel{~}{C}_{7,AE}^{(2)}\\ \stackrel{~}{C}_{8,AE}^{(2)}\end{array}\right)=\left(\begin{array}{ccccccccccccc}0& & & & & & & & & & & & \\ & & & \text{}& & & & & & & & & \\ & \text{}& & & & & & & & & & & \\ & & & & & & & & & & & & \\ & & & \text{}& & & & & & & & & \\ & \text{}& & & & & & & & & & & \\ & & & & 0& & & & & & & & \end{array}\right)\left(\begin{array}{c}p_{44}\\ p_{43}\\ \mathrm{}\\ p_{40}\\ \mathrm{}\\ p_{4,3}\\ p_{4,4}\end{array}\right),$$
(57)
where $``$ represents a non-vanishing complex number. Note that for simplicity we have not included the data $`\stackrel{~}{C}_{0,AA}^{(2)}`$ and $`\stackrel{~}{C}_{0,EE}^{(2)}`$, which contain degeneracy between the multipole moments (see Sec. B). The matrices in the right hand side correspond to $`𝐀^{(i)}`$ in Eq. (30). These sparse forms of matrix are typical for the present problem. We perform the singular value decomposition with respect to these sparse matrices, and construct pseudo-inverse matrices which give the least-squares solutions of $`p_\mathrm{}m`$. In a similar way the least-squares solutions of $`\mathrm{}=\mathrm{odd}`$ modes are obtained from the following matrix equation,
$$\left(\begin{array}{c}\stackrel{~}{C}_{+8,AE}^{(3)}\\ \mathrm{}\\ \stackrel{~}{C}_{8,AE}^{(3)}\\ \stackrel{~}{C}_{4,AT}^{(3)}\\ \mathrm{}\\ \stackrel{~}{C}_{4,AT}^{(3)}\\ \stackrel{~}{C}_{4,ET}^{(3)}\\ \mathrm{}\\ \stackrel{~}{C}_{4,ET}^{(3)}\end{array}\right)=\left(\begin{array}{ccccccccccccc}& & & & & & & & & & & & \\ & & 𝐀_{\mathrm{AE}}^{(3)}& & & & & & & & & & \\ & & & & & & & & & & & & \\ & & 𝐀_{\mathrm{AT}}^{(3)}& & & & & & & & & & \\ & & & & & & & & & & & & \\ & & 𝐀_{\mathrm{ET}}^{(3)}& & & & & & & & & & \end{array}\right)\left(\begin{array}{c}p_{11}\\ p_{10}\\ p_{1,1}\\ p_{33}\\ \mathrm{}\\ p_{3,3}\\ p_{55}\\ \mathrm{}\\ p_{5,5}\end{array}\right),$$
(58)
where we have symbolically written the matrix $`𝐀^{(3)}`$ due to space limitation.
## Appendix D Multipole coefficients for gravitational-wave backgrounds
Here, we give the multipole coefficients for normalized intensity distribution $`P(𝛀)`$ of galactic GWB presented in Sec.IV.1. To evaluate this, we first create the intensity map $`P(𝛀)`$ with $`129\times 256`$ regular grids of celestial sphere $`(\theta ,\varphi )`$. Then, the spherical harmonic expansion of the intensity map is numerically carried out using the SPHEREPACK 3.1 package Adams and Swarztrauber (2003). Table 3 summarizes the multipole coefficients $`p_\mathrm{}m`$ up to $`\mathrm{}=5`$, which are relevant to the analysis in Sec.IV. Note also that $`p_{\mathrm{},m}=(1)^mp_\mathrm{}m^{}`$.
To characterize the contribution of the $`\mathrm{}`$-th moment to the galactic GWB, we introduce the angular power defined by
$`\sigma _{\mathrm{}}=\left\{{\displaystyle \frac{1}{2\mathrm{}+1}}{\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}|p_\mathrm{}m|^2\right\}^{1/2},`$ (59)
which is rotationally invariant Kudoh and Taruya (2005a). Figure 12 shows the normalized angular power, $`\sigma _{\mathrm{}}/\sigma _0`$ up to $`\mathrm{}=30`$. The dominant contribution to the intensity of galactic GWB comes from the multipoles with $`\mathrm{}4`$, however, the asymptotic behavior at higher multipoles is very slow and can be fitted by $`\sigma _{\mathrm{}}/\sigma _01.85e^{0.005\mathrm{}}/\mathrm{}`$ (dotted line in Fig.12), which turns out to be a good approximation even to much higher multipoles, $`\mathrm{}100`$.
## Appendix E Estimation of statistical error in reconstruction analysis
In the reconstruction analysis based on the harmonic-Fourier representation, the statistical errors plotted in Fig.11 are roughly estimated as follows. In the presence of the noises, the least-squares solution given in equation (32) becomes
$$𝐩_{\mathrm{approx}}^{(i)}=\left[𝐀^{(i)}\right]^+\left\{𝐜^{(i)}+𝐬_\mathrm{n}^{(i)}\right\},$$
(60)
where the additional term $`𝐬_\mathrm{n}^{(i)}`$ represents the noise contributions to the $`i`$-th order coefficient of perturbative expansion for $`𝐜(f)`$. The root-mean-square amplitude of the error $`\mathrm{\Delta }𝐩^{(i)}`$ is then defined by taking the ensemble average of the random noises as $`\mathrm{\Delta }𝐩^{(i)}|𝐩^{(i)}𝐩^{(i)}|^2`$, which gives the errors in multipole coefficient $`\mathrm{\Delta }p_\mathrm{}m`$. The $`j`$-th components of the vector $`\mathrm{\Delta }𝐩^{(i)}`$ becomes
$$|\mathrm{\Delta }𝐩_{\mathrm{approx},j}^{(i)}|^2=\left[𝐀^{(i)}\right]_{jk}^+\left[𝐀^{(i)}\right]_{jk}^+|𝐬_{\mathrm{n},k}^{(i)}|^2.$$
(61)
Here, the variance $`|𝐬_{\mathrm{n},\mathrm{k}}^{(i)}|^2`$ roughly corresponds to the quadrature of the vector $`𝐜^{(i)}`$ divided by the signal-to-noise ratio:
$$𝐒_{\mathrm{n},k}^{(i)}=\left\{\alpha \frac{𝐜_k^{(i)}}{(𝐬/𝐫)_k}\right\}^2.$$
(62)
In the above expression, the vector $`(𝐬/𝐫)`$ represents the SNR, each component of which is the quantity $`(S/N)_k`$ defined in (38) just for the same component of the vector $`𝐜^{(i)}`$. Notice that the factor $`\alpha `$ is multiplied in equation (62) by hand. The reason why we have introduced the factor $`\alpha `$ is as follows. First note that the quantity $`(S/N)_k`$ basically reflects the signal-to-noise ratio for the most dominant term in the perturbative expansion for the signal $`\stackrel{~}{C}_k(f)`$ (see Eq.(30)). For the $`AE`$-correlation, $`(S/N)_k`$ gives the signal-to-noise ratio for second-order term, i.e., $`𝐜^{(2)}`$. For the $`AT`$-correlation, $`(S/N)_k`$ reflects the signal-to-noise ratio for $`𝐜^{(3)}`$. In our reconstruction analysis, the higher-order contributions to the $`AE`$-correlation, $`𝐜^{(3)}`$ is used to make the skymap with multipoles $`\mathrm{}5`$, which can be estimated by analyzing the multi-frequency data. In general, the signal from higher-order contribution is weaker than the lowest-order term, leading to the calibration error. The significance of this error would be enhanced by the factor roughly proportional to $`\widehat{f}^{(i2)}`$ for the $`i`$-th order terms of $`AE`$-correlation. Hence, in order to mimic this, the factor $`\alpha `$ is multiplied and set to $`\alpha =\widehat{f}^{(i2)}`$. In the case examined in Fig.11, we adopt $`\alpha =\widehat{f}^1=10`$ for third-order cross-correlation data $`\stackrel{~}{C}_{\mathrm{AE},k}^{(3)}`$. Otherwise we set $`\alpha =1`$. As a result, statistical errors of $`\mathrm{}=3`$ modes become larger than those in $`\mathrm{}=`$even modes (see Fig.11). Note that the multipole coefficients with $`\mathrm{}=1`$ in case A and with $`\mathrm{}=5`$ in both case are completely degenerate and cannot be recovered by reconstruction analysis. Hence, the statistical errors were not evaluated.
## Appendix F Computational method of multipole coefficients
In this Appendix, we give a brief description of the numerics of calculating the multipole coefficients using the SPHEREPACK 3.1 package Adams and Swarztrauber (2003). The traditional spherical harmonic transform of a scalar function is
$`\mathrm{\Psi }(\theta ,\varphi )={\displaystyle \underset{\mathrm{},m}{}}p_\mathrm{}mY_\mathrm{}m(\theta ,\varphi ).`$ (63)
The spherical harmonics are given by equation (42) and the associated Legendre polynomials are
$`P_{\mathrm{}}^m(x)={\displaystyle \frac{(1)^m}{2^{\mathrm{}}\mathrm{}!}}(1x^2)^{m/2}{\displaystyle \frac{d^{m+\mathrm{}}}{dx^{m+\mathrm{}}}}(x^21)^{\mathrm{}}.`$ (64)
On the other hand, numerical computation of spherical harmonic expansion is performed with the subroutine `shaec` in SPHEREPACK 3.1 package, which directly gives the following expansion form:
$`\mathrm{\Psi }(\theta ,\varphi )`$ $`=`$ $`\sqrt{{\displaystyle \frac{\pi }{2}}}{\displaystyle \underset{\mathrm{}=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{}_{}{}^{}P_{\mathrm{}}^{m}(\mathrm{cos}\theta )\gamma _\mathrm{}m\left[\alpha _\mathrm{}m\mathrm{cos}m\varphi \beta _\mathrm{}m\mathrm{sin}m\varphi \right],`$ (65)
where the prime notation on the sum indicates that the first term corresponding to $`m=0`$ is multiplied by the factor $`1/2`$. To relate the coefficients $`\alpha _\mathrm{}m`$ and $`\beta _\mathrm{}m`$ with $`p_\mathrm{}m`$, the expansion (63) is compared with the expression (65), leading to
$`p_\mathrm{}m`$ $`=`$ $`\sqrt{{\displaystyle \frac{\pi }{2}}}(\alpha _\mathrm{}m+i\beta _\mathrm{}m),`$ (66)
$`(1)^mp_{\mathrm{},m}`$ $`=`$ $`\sqrt{{\displaystyle \frac{\pi }{2}}}(\alpha _\mathrm{}mi\beta _\mathrm{}m),`$ (67)
for $`m>0`$, and
$$p_\mathrm{}0=\sqrt{\frac{\pi }{2}}\alpha _\mathrm{}0,$$
(68)
for $`m=0`$. Hence, with a help of these expressions, one can read off the multipole coefficients $`p_\mathrm{}m`$ from the expansion formula (65). Notice that the associated Legendre polynomials used in the SPHEREPACK 3.1 package differ from (64) by a factor $`(1)^m`$ so that one must multiply this factor to the transformation law (67) and (68). |
warning/0507/hep-ph0507219.html | ar5iv | text | # Process 𝑒⁺𝑒⁻→3𝜋(𝛾) with final state radiative corrections
## I Introduction
The problem of taking into account of radiative corrections (RC) of lowest order of perturbation theory to process of three pion production in annihilation channel of colliding $`e^+e^{}`$ at moderately high energies becomes urgent for precision measurement of hadronic contribution to muon anomalous magnetic moment $`(g2)_\mu `$ Bennett:2004pv . This is the motivation of this paper. The similar calculation for production of 2 $`\pi `$ and $`\mu ^+\mu ^{}`$ at the annihilation channel was performed recently Bystritskiy:2005ib .
We shall use the pion sector of chiral perturbation theory (ChPT) to perform the calculations of interaction of pions with electromagnetic field. This theory is given by Wess-Zumino-Witten effective lagrangian Witten:1983tx ; Wess:1971yu . The relevant piece of this lagrangian is reproduced below:
$``$ $`=`$ $`{\displaystyle \frac{f_\pi ^2}{4}}Sp\left[D_\mu U\left(D^\mu U\right)^++\chi U^++U\chi ^+\right]`$ (1)
$``$ $`{\displaystyle \frac{e}{16\pi ^2}}\epsilon ^{\mu \nu \alpha \beta }A_\mu Sp\left[Q\left\{\left(_\nu U\right)\left(_\alpha U^+\right)\left(_\beta U\right)U^+\left(_\nu U^+\right)\left(_\alpha U\right)\left(_\beta U^+\right)U\right\}\right]`$
$``$ $`{\displaystyle \frac{ie^2}{8\pi ^2}}\epsilon ^{\mu \nu \alpha \beta }\left(_\mu A_\nu \right)A_\alpha Sp[Q^2\left(_\beta U\right)U^++Q^2U^+\left(_\beta U\right)+`$
$`+{\displaystyle \frac{1}{2}}QUQU^+\left(_\beta U\right)U^+{\displaystyle \frac{1}{2}}QU^+QU\left(_\beta U^+\right)U],`$
where $`f_\pi =94`$ MeV is the pion decay constant, $`U=exp\left(i\frac{\sqrt{2}\mathrm{\Phi }}{f_\pi }\right)`$, $`D_\mu U=_\mu U+ieA_\mu [Q,U]`$, $`Q=diag(\frac{2}{3},\frac{1}{3},\frac{1}{3})`$ is the quark charge matrix and terms with $`\chi =Bdiag(m_u,m_d,m_s)`$ introduce explicit chiral symmetry breaking due to nonzero quark masses. The constant B has dimension of mass and is determined through the equation $`Bm_q=M^2`$, where $`m_q=m_um_d`$ and $`M`$ is the pion mass. The pseudoscalar meson matrix $`\mathrm{\Phi }`$ has its standard form:
$`\mathrm{\Phi }=\left(\begin{array}{ccc}\frac{1}{\sqrt{2}}\pi ^0+\frac{1}{\sqrt{6}}\eta & \pi ^+& K^+\\ \pi ^{}& \frac{1}{\sqrt{2}}\pi ^0+\frac{1}{\sqrt{6}}\eta & K^0\\ K^{}& \overline{K}^0& \frac{2}{\sqrt{6}}\eta \end{array}\right).`$ (5)
## II Born approximation
We consider reaction of $`e^+e^{}`$ annihilation into three pions:
$`e^{}(p_{})+e^+(p_+)\pi ^{}(q_{})+\pi ^+(q_+)+\pi ^0(q_0).`$ (6)
Matrix element for this process in Born approximation have a form (see Fig. 1):
$`M^{(0)}={\displaystyle \frac{i\alpha }{\pi f_\pi ^{(0)3}}}{\displaystyle \frac{1}{q^2}}\overline{v}(p_+)\gamma _\mu u(p_{})(\mu q_+q_{}q_0),`$ (7)
where $`f_\pi ^{(0)}`$ is the unrenormalized pion decay constant, $`s=q^2=(p_++p_{})^2`$ – invariant mass of initial state, $`(\mu q_+q_{}q_0)\epsilon _{\mu \nu \alpha \beta }q_+^\nu q_{}^\alpha q_0^\beta `$.
Squaring this matrix element and performing summation over initial leptons spin states we get the total cross section for this process:
$`\sigma _B^{(0)}={\displaystyle \frac{\alpha ^2s^2}{2^83\pi ^5f_\pi ^{(0)6}}}{\displaystyle _{x_+^{min}}^{x_+^{max}}}𝑑x_+{\displaystyle _{x_{}^{min}}^{x_{}^{max}}}𝑑x_{}G(x_+,x_{}),`$ (8)
here
$`G(x_+,x_{})=4(x_+^2\mu ^2)(x_{}^2\mu ^2)\left(12x_+2x_{}+2x_+x_{}+\mu ^2\right)^2,`$ (9)
with $`\mu ^2=M^2/s`$, $`M`$ – is pion mass, $`x_\pm =\epsilon _\pm /\sqrt{s}`$, $`x_0=\epsilon _0/\sqrt{s}`$ – are fractions of final pion’s energies, while $`x_++x_{}+x_0=1`$. Limits of integration in (8) follows from kinematical constrains:
$`x_+^{min}`$ $`=`$ $`\mu ,x_+^{max}={\displaystyle \frac{1}{2}}\left(1{\displaystyle \frac{3M^2}{s}}\right)={\displaystyle \frac{1}{2}}\left(13\mu ^2\right),`$
$`x_{}^{max,min}`$ $`=`$ $`{\displaystyle \frac{1}{2(12x_++\mu ^2)}}\left((1x_+)(12x_++\mu ^2)\pm R(x_+)\right),`$
where $`R^2(x_+)=(x_+^2\mu ^2)(12x_++\mu ^2)(12x_+3\mu ^2)`$.
## III Virtual photon emission
Radiative corrections from emission of virtual photon can be represented by 11 Feynman diagrams (FD) (see Fig. 2).
First we shall notice that FDs 4, 5, 6, 7, 10, 11 gives zero contribution ($`\delta _4=\delta _5=\delta _6=\delta _7=\delta _{10}=\delta _{11}=0`$) due to
$`{\displaystyle \frac{d^4k}{i\pi ^2}\frac{(\mu kq(2q_{}k))}{k^2(k^22kq_{})(k^22kq+q^2M^2)}}0.`$ (10)
FDs 2 and 3 are the contributions from pion wave function renormalization and they are equal to Chang :
$`\delta _c=\delta _2+\delta _3={\displaystyle \frac{\alpha }{\pi }}\left(L_\mathrm{\Lambda }+\mathrm{ln}{\displaystyle \frac{M^2}{\lambda ^2}}{\displaystyle \frac{3}{4}}\right),`$ (11)
where $`L_\mathrm{\Lambda }=\mathrm{ln}(\mathrm{\Lambda }^2/M^2)`$ and $`\mathrm{\Lambda }`$ – is the ultraviolet cut-off parameter, $`\lambda `$ – is the fictitious photon mass. Considering contributions of FDs 1, 8, 9 we get:
$`\delta _v`$ $`=`$ $`\delta _1+\delta _8+\delta _9={\displaystyle \frac{\alpha }{\pi }}[1+{\displaystyle \frac{1}{2}}L_\mathrm{\Lambda }{\displaystyle \frac{1+\beta ^2}{2\beta }}\mathrm{ln}\left({\displaystyle \frac{1+\beta }{1\beta }}\right)\mathrm{ln}{\displaystyle \frac{M^2}{\lambda ^2}}+`$ (12)
$`+{\displaystyle \frac{1+\beta ^2}{4}}s_1\text{Re}{\displaystyle _0^1}{\displaystyle \frac{dx}{q_x^2}}(\mathrm{ln}{\displaystyle \frac{q_x^2}{M^2}}2)+Q],`$
$`Q`$ $`=`$ $`{\displaystyle _0^1}𝑑x{\displaystyle _0^1}𝑑yy\mathrm{ln}{\displaystyle \frac{x(1xy)+xx_+(1y)+\mu ^2y}{x(1xy)+xx_{}(1y)+\mu ^2y}}.`$ (13)
Here $`s_1=(q_++q_{})^2`$ is the invariant mass of charged pions pair, $`\beta ^2=14M^2/s_1`$, $`q_x^2=M^2s_1x(1x)i0`$. The integrals in (12) can be calculated explicitly:
$`\text{Re}{\displaystyle _0^1}{\displaystyle \frac{dxs_1}{q_x^2}}`$ $`=`$ $`{\displaystyle \frac{2}{\beta }}L,`$ (14)
$`\text{Re}{\displaystyle _0^1}{\displaystyle \frac{dxs_2}{q_x^2}}\mathrm{ln}{\displaystyle \frac{q_x^2}{M^2}}`$ $`=`$ $`{\displaystyle \frac{4}{\beta }}\left[L\mathrm{ln}\left({\displaystyle \frac{1+\beta }{2\beta }}\right){\displaystyle \frac{1}{4}}L^2+\mathrm{Li}_2\left({\displaystyle \frac{1\beta }{1+\beta }}\right)+2\xi _2\right],`$ (15)
where $`L=\mathrm{ln}\frac{1+\beta }{1\beta }`$, $`\mathrm{Li}_2\left(x\right)=_0^x\frac{dt}{t}\mathrm{ln}(1t)`$ and $`\xi _2=\frac{\pi ^2}{6}`$.
## IV Soft real photon emission
The standard calculation of contribution of real soft photon emission by final pions
$`\delta _s={\displaystyle \frac{\sigma _{soft}}{\sigma _B}}={\displaystyle \frac{\alpha }{4\pi ^2}}{\displaystyle \frac{d^3k}{\omega }\left(\frac{q_{}}{kq_{}}\frac{q_+}{kq_+}\right)^2}|_{\omega <\mathrm{\Delta }\epsilon },`$ (16)
where $`\mathrm{\Delta }\epsilon `$ is the maximum energy of soft photon (i.e. $`\omega <\mathrm{\Delta }\epsilon `$), leads to:
$`\delta _s`$ $`=`$ $`{\displaystyle \frac{2\alpha }{\pi }}\{(\mathrm{ln}\mathrm{\Delta }{\displaystyle \frac{1}{2}}\mathrm{ln}(x_+x_{})+\mathrm{ln}{\displaystyle \frac{M}{\lambda }})(1+{\displaystyle \frac{1+\beta ^2}{2\beta }}L)+`$ (17)
$`+{\displaystyle \frac{1+\beta ^2}{4\beta }}[g{\displaystyle \frac{1}{2}}L^2+L\mathrm{ln}\left({\displaystyle \frac{4}{1\beta ^2}}\right)\xi _22\mathrm{L}\mathrm{i}_2({\displaystyle \frac{1\beta }{1+\beta }})]\},`$
where $`\mathrm{\Delta }=\mathrm{\Delta }\epsilon /\sqrt{s}`$ and the quantity $`g`$ is defined by
$`g=2\beta {\displaystyle _0^1}{\displaystyle \frac{dt}{1\beta ^2t^2}}\mathrm{ln}\left(1+{\displaystyle \frac{1t^2}{4}}{\displaystyle \frac{\left(x_+x_{}\right)^2}{x_+x_{}}}\right).`$ (18)
## V Hard real photon emission
Consider the contribution of radiative corrections which arises from the emission of additional hard photon by final particles, i.e. the process:
$`e^{}(p_{})+e^+(p_+)\pi ^{}(q_{})+\pi ^+(q_+)+\pi ^0(q_0)+\gamma (k).`$ (19)
Amplitude for this process can be written in the form:
$`M=(4\pi \alpha )^2{\displaystyle \frac{1}{q^2}}\overline{v}(p_+)\gamma _\mu u(p_{}){\displaystyle \frac{1}{4\pi ^2f_\pi ^3}}T_{\mu \nu }e^\nu (k),`$ (20)
where $`k`$, $`e_\mu (k)`$ are the momenta and the polarization vector of final real photon. $`T_{\mu \nu }`$ is the tensor corresponding to the $`\gamma ^{}(\mu ,Q)\pi ^+(q_+)\pi ^{}(q_{})\pi ^0(q_0)\gamma (\nu ,k)`$ vertex, which follows from (1):
$`T^{\mu \nu }`$ $`=`$ $`(\mu \nu Qk)A+(\mu \nu (Q+k)q_0)+(\mu \lambda Qq_0)\left({\displaystyle \frac{q_{}^\nu q_+^\lambda }{(q_{}k)}}+{\displaystyle \frac{q_+^\nu q_{}^\lambda }{(q_+k)}}\right)`$ (21)
$`(\nu \lambda kq_0)\left({\displaystyle \frac{(2q_{}Q)^\mu q_+^\lambda }{Q^22(q_{}Q)}}+{\displaystyle \frac{(2q_+Q)^\mu q_{}^\lambda }{Q^22(q_+Q)}}\right),`$
where $`Q=q_++q_{}+q_0+k`$, (thus $`Q_0=\sqrt{s}`$), $`A=1(s_1M^2)/((Qk)^2M^2)`$. This tensor satisfies the gauge-invariance condition for both photon legs:
$`Q_\mu T^{\mu \nu }=k_\nu T^{\mu \nu }=0.`$ (22)
Thus we may write the cross section in a form:
$`\sigma ^{hard}={\displaystyle \frac{(4\pi \alpha )^3}{8s}}{\displaystyle \frac{1}{s^2}}Sp\left[\widehat{p}_+\gamma ^\mu \widehat{p}_{}\gamma ^\nu \right]\left({\displaystyle \frac{1}{4\pi ^2f_\pi ^3}}\right)^2{\displaystyle \frac{1}{3}}\left(g_{\mu \nu }{\displaystyle \frac{Q_\mu Q_\nu }{Q^2}}\right){\displaystyle 𝑑\mathrm{\Gamma }_4\underset{\lambda }{}\left|T^{\alpha \beta }e_\beta ^\lambda \right|^2}.`$ (23)
Phase volume for final state has a form:
$`d\mathrm{\Gamma }_4`$ $`=`$ $`(2\pi )^4\delta ^4\left(p_++p_{}q_+q_{}q_0k\right){\displaystyle \frac{d^3q_+}{(2\pi )^32\epsilon _+}}{\displaystyle \frac{d^3q_{}}{(2\pi )^32\epsilon _{}}}{\displaystyle \frac{d^3q_0}{(2\pi )^32\epsilon _0}}{\displaystyle \frac{d^3k}{(2\pi )^32\omega }}=`$ (24)
$`=`$ $`(2\pi )^8{\displaystyle \frac{s^2\pi ^2}{16}}xdxdx_+dx_{}dO_\gamma ,`$
where $`x=\omega /\sqrt{s}`$$`x+x_++x_{}+x_0=1`$. And now the cross section takes the form:
$`\sigma ^{hard}={\displaystyle \frac{\alpha ^3}{2^83\pi ^7f_\pi ^6}}{\displaystyle 𝑑xx𝑑x_+𝑑x_{}𝑑O_\gamma \left(\underset{\lambda }{}\left|T^{\mu \nu }\right|^2\right)}|_{x>\mathrm{\Delta }}.`$ (25)
We shall notice, that the sum of hard photon emission $`\delta _h=\sigma ^{hard}/\sigma _B`$ and soft real photon emission $`\delta _s`$ contributions ($`\delta _h+\delta _s`$) does not contain the auxiliary parameter $`\mathrm{\Delta }`$. To see this explicitly let us consider the small $`x=\omega /\sqrt{s}`$ limit of $`\sigma ^{hard}`$. Really if we consider the case $`\mathrm{\Delta }\sqrt{s}<\omega \sqrt{s}`$ then (see (21)):
$`T^{\mu \nu }e_\nu (k)|_{\omega \sqrt{s}}(\mu q_+q_{}q_0)\left({\displaystyle \frac{(q_{}e)}{(q_{}k)}}{\displaystyle \frac{(q_+e)}{(q_+k)}}\right).`$ (26)
We can calculate the hard photon emission contribution in this limit:
$`{\displaystyle \frac{d^3k}{2\pi \omega }\underset{\lambda }{}\left(\left|T^{\mu \nu }e_\nu (k)\right|^2\right)}|_{\omega \mathrm{\Delta }\sqrt{s}}=4\mathrm{ln}{\displaystyle \frac{1}{\mathrm{\Delta }}}\left({\displaystyle \frac{1+\beta ^2}{2\beta }}\mathrm{ln}\left({\displaystyle \frac{1+\beta }{1\beta }}\right)1\right),`$ (27)
and we get for $`\delta _h`$:
$`\delta _h|_{\omega \mathrm{\Delta }\sqrt{s}}{\displaystyle \frac{2\alpha }{\pi }}\left[\mathrm{ln}{\displaystyle \frac{1}{\mathrm{\Delta }}}\left({\displaystyle \frac{1+\beta ^2}{2\beta }}\mathrm{ln}\left({\displaystyle \frac{1+\beta }{1\beta }}\right)1\right)+O(\mathrm{\Delta })\right].`$ (28)
We redefine the contributions in the following manner:
$`\delta _s+\delta _h\overline{\delta }_s+\overline{\delta }_h,`$ (29)
$`\overline{\delta }_s`$ $`=`$ $`\delta _s+{\displaystyle \frac{2\alpha }{\pi }}\mathrm{ln}{\displaystyle \frac{1}{\mathrm{\Delta }}}\left({\displaystyle \frac{1+\beta ^2}{2\beta }}\mathrm{ln}\left({\displaystyle \frac{1+\beta }{1\beta }}\right)1\right),`$ (30)
$`\overline{\delta }_h`$ $`=`$ $`\delta _h{\displaystyle \frac{2\alpha }{\pi }}\mathrm{ln}{\displaystyle \frac{1}{\mathrm{\Delta }}}\left({\displaystyle \frac{1+\beta ^2}{2\beta }}\mathrm{ln}\left({\displaystyle \frac{1+\beta }{1\beta }}\right)1\right),`$ (31)
where both $`\overline{\delta }_s`$ and $`\overline{\delta }_h`$ are not dependent on $`\mathrm{\Delta }`$ any more.
## VI Conclusion
The final result is
$`\sigma ^{ee3\pi (\gamma )}`$ $`=`$ $`\sigma _B^{(0)}\left(1+\delta _{sw}\right)\left(1+\delta \right)=\sigma _B\left(1+\delta \right),\sigma _B=\sigma _B^{(0)}\left(f_\pi ^{(0)}f_\pi \right)`$ (32)
$`\delta `$ $`=`$ $`\left(\delta _c+\delta _v+\overline{\delta }_s+\overline{\delta }_h\right)|_{L_\mathrm{\Lambda }=0},\delta _{sw}={\displaystyle \frac{3\alpha }{2\pi }}L_\mathrm{\Lambda },`$ (33)
where we extracted the short-distance contributions in form $`\left(1+\delta _{sw}\right)`$ and used this factor to renormalize $`f_\pi ^{(0)}`$ pion decay constant in form $`f_\pi ^{(0)6}\left(1+\delta _{sw}\right)=f_\pi ^6`$ Holstein:ua .
The explicit form of $`\delta \overline{\delta }_h`$ is:
$`\delta \overline{\delta }_h`$ $`=`$ $`{\displaystyle \frac{\alpha }{\pi }}\{{\displaystyle \frac{1}{2}}\mathrm{ln}\left(x_+x_{}\right)(1+{\displaystyle \frac{1+\beta ^2}{2\beta }}L)+{\displaystyle \frac{1}{4}}+Q+`$ (34)
$`+{\displaystyle \frac{1+\beta ^2}{4\beta }}\left[g{\displaystyle \frac{1}{2}}L^2+L\mathrm{ln}{\displaystyle \frac{4}{1\beta ^2}}\xi _22\mathrm{L}\mathrm{i}_2\left({\displaystyle \frac{1\beta }{1+\beta }}\right)\right]+`$
$`+{\displaystyle \frac{1+\beta ^2}{\beta }}[L+L\mathrm{ln}{\displaystyle \frac{1+\beta }{2\beta }}{\displaystyle \frac{1}{4}}L^2+\mathrm{Li}_2({\displaystyle \frac{1\beta }{1+\beta }})+2\xi _2]\},`$
with $`L=\mathrm{ln}\left(\frac{1+\beta }{1\beta }\right)`$, $`g`$ defined in (18) and $`Q`$ defined in (13). We note that $`g(x_+=x_{})=Q(x_+=x_{})=0`$. The form of $`\overline{\delta }_h`$ depends on the experimental conditions of final state particles registration and not considered here.
In Fig. 3 we present the value $`\delta \overline{\delta }_h`$ for typical experimental situation.
###### Acknowledgements.
We are grateful to G. V. Fedotovich for attracting our attention to this problem. We are also grateful to Z. K. Silagadze for valuable discussions. |
warning/0507/gr-qc0507025.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The theory of unification of interactions has as its aim to elucidate the basic properties of the micro-world at high energies. An interesting question is the following: is there not a tendency of gravitational and electromagnetic interactions to approach (at least in quantitative terms) already at a macro-level in situations where the interacting objects have very high relative velocities? In other words, are we not getting an indication, in the framework of general theory of relativity and classical electromagnetism, of a tendency of diminution of the difference between these interactions when relative velocities of the interacting classical particles are ultrarelativistic?
Certainly, investigations of analogies between gravitational and electromagnetic interactions have their own long history on the level of Newton’s Law of universal gravitation and Coulomb’s Law. In the last decades fundamentally new common traits of these interactions were brought to light. In particular, terms such as ”gravitoelectric field”, ”gravitomagnetic field”, and ”gravitoelectromagnetism” have gained currency within the general theory of relativity, along with the elaboration of their context, . Particularly worthy of attention is an early publication .
Interesting aspects of analogies between gravitation and electromagnetism are revealed when one looks at the behavior of a classical test particle with spin in a gravitational field . Examining the Mathisson-Papapetrou (MP) equations in a Kerr field, considers gravitational spin-orbit and spin-spin interactions in respective approximations in power of $`1/c`$ and compares them with analogous electromagnetic interactions.
The concept of the electromagnetic field was preceded by concepts of two independent entities, the electric and the magnetic interaction, which became unified in Maxwell’s theory. On the other hand, in the case of Einstein’s theory of gravitation just the opposite took place: the general theory of relativity was from the start a theory of a single gravitational field and only with time did many investigators begin to feel the need to treat as separate (although closely linked) two of its components: gravitoelectric and gravitomagnetic. The separation of gravitoelectric and gravitomagnetic components of the gravitational field in the general theory of relativity is carried out with the Riemann tensor as the basic characteristic of the field . In the local orthonormal basis, the gravitoelectric components of the gravitational field $`E_{(k)}^{(i)}`$ are determined by the relationship
$$E_{(k)}^{(i)}=R^{(i)(4)}{}_{(k)(4)}{}^{},$$
(1)
where $`R^{(i)(4)}_{(k)(4)}`$ denotes local components of the Riemann tensor. (Here and in the following the indices of the orthogonal tetrads are placed in round parentheses; Latin indices run through values 1, 2, 3, while Greek indices through 1, 2, 3, 4). Correspondingly, for gravitomagnetic components $`B_{(k)}^{(i)}`$ we have
$$B_{(k)}^{(i)}=\frac{1}{2}R^{(i)(4)}{}_{(m)(n)}{}^{}\epsilon _{}^{(m)(n)}{}_{(k)}{}^{},$$
(2)
where $`\epsilon ^{(m)(n)}_{(k)}`$ is the Levi-Civita tensor.
We shall begin our investigation by analyzing relationships (1) and (2) in the concrete case of a gravitational field created by a Schwarzschild mass moving relative to an observer with arbitrary velocity.
We point out that in interesting paper the similar to a certain extent problem was considered. Namely, ”…it is shown that the gravitational field of a fast-moving mass bears an increasing resemblance to a plane gravitational wave, the greater the speed of the mass” , p. 96. However, in this paper the gravitational field of a fast-moving Schwarzschild mass was considered only in the context of investigations of the gravitational waves in the general theory of relativity. The influence of components (1), (2) on the other masses was not under investigation in .
The gravitational field of a massless particle which moves with the velocity of light was considered in . It was shown that the gravitational field of this particle is nonvanishing only on a plane containing the particle and orthogonal to the direction of motion. The results of were used in for investigating the ultrarelativistic collision of two black holes. The detail elucidation of this problem one can find in . Passing from the infinite Lorentz $`\gamma `$\- factor (the case of a massless particle with the velocity of light) to the large but finite $`\gamma `$ (a massive ultrarelativistic particle) is described taking into account the small parameter $`\gamma ^1`$ and the corresponding small corrections to the metric of a massless particle . This approach to the description of the dependence of the gravitational field of a moving massive particle on the $`\gamma `$ differs from the method and results of . In the Riemann tensor components were calculated for any value $`1<\gamma <\mathrm{}`$ without the consideration of the case of the infinite $`\gamma `$ as the initial approximation.
In what follows, we shall use a system of units where $`c=G=1`$.
## 2 Dependence of the gravitoelectromagnetic field of a moving Schwarzschield source on the Lorentz $`\gamma `$-factor
The results of this Section may be considered as a direct development of . The main idea of is the comparison of the canonical forms of the Riemann tensor for a plane gravitational wave and a fast-moving mass (the $`3\times 3`$ matrices $`P`$ and $`Q`$ in the notation of ). Our purpose is the analysis of the action of different components (1), (2), as measured by a fast-moving observer in the Schwarzschild field, on the test masses.
We shall label the reference frame, which moves with respect to the source of the Schwarzschild field in an arbitrary direction and with arbitrary velocity, by a set of corresponding terads $`\lambda _{(\beta )}^\alpha `$. The Schwarzschild metric we consider in standard coordinates $`x^1=r,x^2=\theta ,x^3=\phi ,x^4=t`$. For expediency and without loss of generality we assume the directions of the space axes of the ortho-reference to be as follows: The first axis is perpendicular to the plane determined by the direction of observer motion and the radial direction to the field source. (Obviously, in the particular case of radial motion there is freedom of choice). The second axis coincides with the direction of motion. As a consequence, we note that the following tetrad components have zero components: $`\lambda _{(1)}^1,\lambda _{(1)}^3,\lambda _{(2)}^2,\lambda _{(3)}^2`$. For evaluation of other components we shall use the general relationship between terad components and the metric tensor $`g^{\alpha \pi }`$:
$$\lambda _{(\beta )}^\alpha \lambda _{(\rho )}^\pi \eta ^{(\beta )(\rho )}=g^{\alpha \pi },$$
(3)
where $`\eta ^{(\beta )(\rho )}=diag(1,1,1,1)`$ is the Minkowski tensor.
For the Schwarzschild metric, where only the diagonal elements of the tensor $`g^{\alpha \pi }`$ are different from zero, the system of the ten algebraic equations involving tetrad components of (3) may be separated into subsystems of lower dimensionality, permitting to determined all components that are different from zero:
$$\lambda _{(1)}^2=\sqrt{g^{22}},\lambda _{(2)}^1=u^1u^4\sqrt{\frac{g_{44}}{u_4u^41}},\lambda _{(2)}^3=u^3u^4\sqrt{\frac{g_{44}}{u_4u^41}},$$
$$\lambda _{(2)}^4=\sqrt{\frac{u_4u^41}{g_{44}}},\lambda _{(3)}^1=u^3\sqrt{\frac{g^{11}g_{33}}{u_4u^41}},\lambda _{(3)}^3=u^1\sqrt{\frac{g^{33}g_{11}}{u_4u^41}},$$
$$\lambda _{(4)}^1=u^1,\lambda _{(4)}^3=u^3,\lambda _{(4)}^4=u^4,$$
(4)
where $`u^\mu `$ is the 4-vector of the observer velocity. (The $`\theta `$ angle is measured such that the observer moves in the plane $`\theta =\pi /2`$, which entails $`u^2=0`$. Expressions (4) were used in while dealing with another problem).
According to (1) and (2), in order to evaluate $`E_{(k)}^{(i)}`$ and $`B_{(k)}^{(i)}`$ it is necessary to have the values of the local components of the Riemann tensor, which are connected to its global components by the well known relation
$$R_{(\alpha )(\beta )(\gamma )(\delta )}=\lambda _{(\alpha )}^\mu \lambda _{(\beta )}^\nu \lambda _{(\gamma )}^\rho \lambda _{(\delta )}^\sigma R_{\mu \nu \rho \sigma }.$$
(5)
Non-zero components of the Riemann tensor, expressed in standard Schwarzschild coordinates for $`\theta =\pi /2`$, are
$$R_{1212}=R_{1313}=\frac{m}{r2m},R_{2323}=2mr,$$
$$R_{1414}=\frac{2m}{r^3},R_{2424}=R_{3434}=\frac{m}{r}\left(1\frac{2m}{r}\right).$$
(6)
Using (4)–(6), we find the following non-zero components of the Riemann tensor that are present in (1) and (2):
$$R^{(1)(4)}{}_{(1)(4)}{}^{}=\frac{m}{r^3}(3u_3u^31),R^{(2)(4)}{}_{(2)(4)}{}^{}=\frac{2m}{r^3}\frac{u^1u^1}{g_{44}(u_4u^41)}$$
$$+\frac{m}{r}\frac{u^3u^3}{u_4u^41},R^{(2)(4)}{}_{(3)(4)}{}^{}=R^{(3)(4)}{}_{(2)(4)}{}^{}=\frac{3m}{r^2}\frac{u^1u^3u^4}{u_4u^41},$$
$$R^{(3)(4)}{}_{(3)(4)}{}^{}=\frac{m}{r^3}(u_4u^41)+\frac{2m}{r^3}u_4u^4\frac{u_3u^3u_1u^1}{u_4u^41},$$
$$R^{(1)(4)}{}_{(1)(2)}{}^{}=\frac{3mu^3u^3u^4}{r\sqrt{u_4u^41}}(1\frac{2m}{r})^{1/2},$$
$$R^{(1)(4)}{}_{(1)(3)}{}^{}=\frac{3mu^1u^3}{r^2\sqrt{u_4u^41}}(1\frac{2m}{r})^{1/2},$$
$$R^{(2)(4)}{}_{(2)(3)}{}^{}=\frac{3mu^1u^3}{r^2\sqrt{u_4u^41}}(1\frac{2m}{r})^{1/2},$$
$$R^{(3)(4)}{}_{(2)(3)}{}^{}=\frac{3mu^3u^3u^4}{r\sqrt{u_4u^41}}(1\frac{2m}{r})^{1/2}.$$
(7)
Using (7) in (1), we obtain the following non-zero components of the gravitoelectric field:
$$E_{(1)}^{(1)}=\frac{m}{r^3}(1+3u_{}^2),E_{(2)}^{(2)}=\frac{2m}{r^3}+\frac{3m}{r^3}\frac{u_{}^2}{u_4u^41},$$
$$E_{(3)}^{(2)}=E_{(2)}^{(3)}=\frac{3m}{r^3}\frac{u_{}u_{}u^4}{u_4u^41},E_{(3)}^{(3)}=\frac{m}{r^3}\frac{3m}{r^3}\frac{u_{}^2u_4u^4}{u_4u^41}.$$
(8)
where $`u_{}=u^1`$ is the radial component of the 4-velocity, $`u_{}=ru^3`$ is its tangential component. Because of the condition $`u_\mu u^\mu =1`$, here we have the following relationship:
$$u_4u^41=u_{}^2+\left(1\frac{2m}{r}\right)^1u_{}^2.$$
(9)
Similarly, using (7) in (2), we obtain the non-zero components of the gravitomagnetic field,
$$B_{(2)}^{(1)}=B_{(1)}^{(2)}=\frac{3mu_{}u_{}}{r^3\sqrt{u_4u^41}}\left(1\frac{2m}{r}\right)^{1/2},$$
$$B_{(3)}^{(1)}=B_{(1)}^{(3)}=\frac{3mu_{}^2u^4}{r^3\sqrt{u_4u^41}}\left(1\frac{2m}{r}\right)^{1/2}.$$
(10)
Let us stress that relationships (8) and (10) hold true for any arbitrary velocity of the observer.
We shall begin the examination of the components of (8) by simply noting that they have non-zero values even in the Newtonian limit, when $`|u_{}|1`$, $`|u_{}|1`$, $`u^41`$. This had to be expected, inasmuch as in the Newtonian theory there is correspondence between the $`E_{(k)}^{(i)}`$ components and the so-called tidal matrix $`E_{ij}`$, where
$$E_{ij}=\frac{}{x_i}\frac{}{x_j}\phi (\stackrel{}{x},t),$$
(11)
i.e. the second derivatives of the Newtonian potential . The denomination $`E_{ij}`$ is not fortuitous in view of the fact that the components of the tidal acceleration $`a_{i_{tidal}}`$ in the Newtonian theory are, see ,
$$a_{i_{tidal}}=E_{ij}r_j.$$
(12)
In the general theory of relativity the components $`E_{(k)}^{(i)}`$ are also linked with the tidal acceleration, more exactly with the equation of deviation of geodesic lines. Taking (1) into account, this may be written as
$$\frac{D^2l^{(i)}}{ds^2}=E_{(k)}^{(i)}l^{(k)},$$
(13)
where $`s`$ is the proper time, $`l^{(i)}`$ is the vector of relative deviation of two neighboring geodesic lines. It is the equation of deviation of geodesics that was used in , Sec. 31.2, for the analysis of tidal forces felt by an observer while falling onto a Schwarzschild black hole.
According to (13) we have
$$a_{tidal}^{(i)}=E_{(k)}^{(i)}l^{(k)},$$
(14)
Let us note that in , Sec. 31.2, the components $`E_{(k)}^{(i)}`$ are not explicitly mentioned in the analysis of tidal force, but only Riemann tensor components which, according to (1), correspond to $`E_{(k)}^{(i)}`$. (At another place in , in Sec. 1.6, there is mention of the analogy between one part of the Riemann tensor components and the electric field components, and between the other part and the magnetic field components, but the relationships (1) and (2) are not given explicitly). Let us also note that in , Sec. 31.2, the analysis of tidal forces is limited to the case of radial motion, when $`u_{}=0`$. For such motion, according to (8), the components $`E_{(k)}^{(i)}`$ assume the following values:
$$E_{(1)}^{(1)}=\frac{m}{r^3},E_{(2)}^{(2)}=\frac{2m}{r^3},E_{(3)}^{(2)}=E_{(2)}^{(3)}=0,E_{(3)}^{(3)}=\frac{m}{r^3},$$
(15)
that is, they appear completely independent of $`u_{}`$. In view of (14), in such a case the tidal acceleration also does not depend on the velocity of radial motion. The fact that in a radial fall of the observer the tidal forces felt by him are independent of this velocity is, in essence, noted in , Sec. 31.2, (while noting at the same time the analogy with electromagnetism). As a consequence, it is stated in that for a radial falling observer the tidal forces increase sharply only at $`r0`$. The question which remaind unanswed in was the following: what will change if the fall is non-radial? In view of expressions (8), this question can be readily answered. According to (8), the expressions giving the components of the gravitoelectric field in the case of non-radial motion differ significantly from expressions (15) only when the velocity becomes ultrarelativistic. Indeed, inasmuch as $`u_{}`$, $`u_{}`$, $`u^4`$ are proportional to the Lorentz relativistic $`\gamma `$-factor, (8) gives us, for $`|u_{}|1`$, $`|u_{}|1`$, $`u^41`$
$$E_{(1)}^{(1)}\frac{3m}{r^3}\gamma ^2,E_{(2)}^{(2)}\frac{3m}{r^3},E_{(3)}^{(2)}=E_{(2)}^{(3)}\frac{3m}{r^3}\gamma ,E_{(3)}^{(3)}\frac{3m}{r^3}\gamma ^2.$$
(16)
Comparing (16) with (15), we see that while the components in (15) assume arbitrarily large values only when $`r0`$, the components $`E_{(1)}^{(1)}`$, $`E_{(3)}^{(2)}=E_{(2)}^{(3)}`$, $`E_{(3)}^{(3)}`$ of (16) become arbitrarily large already at finite values of $`r`$, provided $`\gamma \mathrm{}`$. (Here we leave aside the question of how to impart to an observer, in practice, a velocity corresponding to large values of $`\gamma `$). Thus, according to (14), an observer in a Schwarzschild field runs the risk of being torn apart by tidal forces not only at $`r0`$, i.e. under the surface of the horizon (which is described in , Sec. 31.2), but even at large values of $`r`$ if his velocity becomes ultrarelativistic.
Obviously, the expressions for the components of the gravitoelectric field (8) and (16) are independently valid, without having to be linked with equations of deviation of geodesic lines and tidal forces. The significance of (8) and (16) resides primarily in characterizing the gravitational field created by a moving Schwarzschild source.
Let us now look at the components of the gravitomagnetic field, (10). It is easy to see that components (10) are different from zero only when $`u_{}0`$, that is only when the observer is moving non-radially. (As is well known, a similar situation arises in electrodynamics for components of the vector of magnetic field intensity of a moving electric charge). Quite generally, the values of components (10) depend signoficantly on observer motion. In the low relativistic region, with $`|u_{}|1`$, $`|u_{}|1`$, $`u^41`$, the common multiplier $`m/r^3`$ of components is further multiplied by corresponding small factors. Whereas in the ultrarelativistic region, where $`|u_{}|1`$, $`|u_{}|1`$, $`u^41`$, this multiplier is further multiplied by large factors, because in this case, according to (10), we have:
$$B_{(2)}^{(1)}=B_{(1)}^{(2)}\frac{3m}{r^3}\gamma ,B_{(3)}^{(1)}=B_{(1)}^{(3)}\frac{3m}{r^3}\gamma ^2.$$
(17)
As we have seen, in the Newtonian limit the components of the gravitomagnetic field (10) have zero values, in contrast to the gravitoelectric field. Moreover, at low relativistic velocities the absolute values of components (10) are considerably smaller than the components (8). Yet at ultrarelativistic velocities , the largest components $`B_{(k)}^{(i)}`$ from (17) and $`E_{(k)}^{(i)}`$ from (16) are of the same order of magnitude, determined by the factor $`3m\gamma ^2/r^3`$. (We point out that this factor is present in the expression for the amplitude of the gravitational wave from ). Consequently, we can conclude that the two aspects of the single gravitational interaction, the gravitoelectric and the gravitomagnetic, show, in the ultrarelativistic range, a tendency of qualitatively drawing together, even though at low velocities they differ substantially.
Befor comparing the gravitational interaction with the electromagnetic in the ultrarelativistic range, we shall examine an important physical situation which evinces the role of the gravitomagnetic interaction.
## 3 A classical particle with spin in an ultrarelativistic gravitational field
Relationships (13), (14) describe a simple experiment in which an observer moving in a Schwarzschild field can determine the values of components of the gravitoelectric field in his own frame of reference. The question arises which experiment would permit this observer to determine the components of the gravitomagnetic field. As is shown in , such an experiment can be carried out by observing the behavior of a test particle with spin. Indeed, according to (9) in , the local components of the 3-acceleration $`a_{(i)}`$ with which the test particle with spin deviates from free geodesic fall in an arbitrary gravitational field is given by the relationship
$$a_{(i)}=\frac{s_{(1)}}{M}R_{(i)(4)(2)(3)},$$
(18)
where $`M`$ is the mass of a test particle. (Here the space axes of the reference frame are chosen such that the spin points along the first axis, so that the spin components are $`s_{(2)}=s_{(3)}=0`$. Expression (18) is a direct cosequence of the MP equations). Taking (2) into consideration, it is not difficult to see that the right side of (18) contains the components of the gravitomagnetic field.
For our concrete case of observer motion in a Schwarzschild field, characterized by the set of tetrads (4), we have
$$a_{(i)}=\frac{s_{(1)}}{M}B_{(i)}^{(1)}.$$
(19)
This result is obtained from (18) taking into account the appropriate components of the curvature tensor from (7) and expressions (10). The non-zero values of $`B_{(k)}^{(i)}`$ in (18) come from (10). Even though the right sides of relationships (19) and (14) have a similar appearance, what is essential is that they contain different components of the gravitational field: in (19) gravitomagnetic and in (14) gravitoelectric. (Let us stress that we are dealing with one and the same reference frame, one connected with an observer moving in a Schwarzschild field). Correspondingly, the nature of forces that cause accelerations (14) and (19) is different: in case of (14) these are tidal forces, and in (19) it is the spin-orbit force. (Detailed discussion of this question may be found in ). Referring to (10) we obtain the magnitude of the 3-acceleration $`|\stackrel{}{a}|`$ with components (19):
$$|\stackrel{}{a}|=\frac{3m}{r^3}\frac{|s_{(1)}u_{}|}{M}\sqrt{1+u_{}^2}.$$
(20)
According to (17), the acceleration components (19) depend, in the case of ultrarelativistic non-radial motion, on the Lorentz $`\gamma `$-factor such that $`a_{(2)}\gamma `$, $`a_{(3)}\gamma ^2`$. The component $`a_{(1)}`$ remains equal to zero at any velocity in the case at hand, where the spin is directed along the first spacial vector of the reference frame, which direction is perpendicular to the plane determined by the direction of observer motion and the radial direction. This is so because in this case the corresponding component of the gravitational field is zero also. Expression (20) also shows that $`|\stackrel{}{a}|\gamma ^2`$.
Thus, the fact that at ultrarelativistic velocity the largest component of the gravitomagnetic field (17) is proportional to $`\gamma ^2`$ entails, on account of (19), that the spin-orbit acceleration is also proportional to $`\gamma ^2`$.
Both, the MP equations and relationships (18)–(20) which follow from them, are rigorously valid for the model of a point test particle with spin, with tidal forces not coming into play. Certainly, for any real macroscopic test particle with rotational motion, tidal and spin-orbit forces become important.
Together, relationships (14) and (19) permit to evaluate these forces. The most significant conclusion of these evaluations lies in the result that for ultrarelativistic non-radial motions the values of these forces in the proper frame of the particle are both proportional to $`\gamma ^2`$.
## 4 Comparison of gravitational and electromagnetic interactions of two particles moving with an ultrarelativistic relative velocity
We shall examine two situations where two particles are mutually interacting. In the first case, we consider two electrically neutral particles with the mass of one being considerably grater than the mass of the other. The particle with smaller mass is endowed with classical spin (internal angular momentum). Thus, we can consider this particle to be the test particle with spin, moving in the gravitational field of the more massive particle, which in its own frame is described by a Schwarzschild metric.
In the second case, the particles carry electric charge, with one charge being considerably larger than the other. Moreover, the particle with the smaller charge has a magnetic moment arising from its internal rotation. The masses of these particles are such that at low relative velocities, when the Coulomb Law and Newton’s Law of gravitational attraction hold true, the force of the electric interaction is very much larger than the gravitational attraction. Again, the particle with the smaller charge and the magnetic moment may be regarded as the test particle. Thus, we may consider that the first pair of particles interacts only gravitationally, and the second pair only electrically.
Let us inquire, in the two cases, how the forces resulting from the gravitational and the electromagnetic interactions, respectively, depend on the magnitude of the relative velocities of the particles. We shall assume that the particles are sufficiently distant one from the other to be able to neglect the respective gravitational and electromagnetic radiation. In accordance with the analysis carried out in the preceding sections, in the first case the force is due to spin and it increases with increasing relative velocity proportionally to $`\gamma ^2`$, as long as the test mass is not moving radially, i.e. it is not moving along the line joining the two masses.
At the same time, in the second case, classical electrodynamics tells us that the force acting on the test particle with the magnetic field is proportional to $`\gamma `$. This means that, no matter how small the gravitational interaction may be in comparison with the electromagnetic interaction in the subrelativistic range of velocities, in passing into the ultrarelativistic range the ratio of the respective forces could, in principle, change to an extent of both forces becoming of the same order of magnitude, provided $`\gamma `$ becomes large enough.
Similar conclusions may be drawn in a third situation, where a model proton interacts with a model electron. (Here we consider two classical particles with masses and charges of a proton and an electron, respectively. In this case, for the description of the gravitational field of the proton as the more massive particle, ane has to refer to the Reissner-Nordstrom metric, rather the Schwarzschild metric, but this does not change the conclusion in principle).
In the Introduction we asked the question, in the framework of general relativity and classical electrodynamics, if there is a tendency of gravitational and electromagnetic interactions to approach quantitatively with increasing relative velocity of interacting particles. We have shown that this question may be answered in the affirmative.
## 5 A case of the ultrarelativistic motion of a classical spinning particle in a Schwarzschild field and the corresponding solution of the Dirac equation
The physical measurements with the ultrarelativistic macroscopic masses are, at least at present, unattainable. Nonetheless, the results from Sections 3 and 4 are not merely academic. If only because the known fact that the general covariant Dirac equation passes, in a quasiclassical approximation, into the MP equations. A concrete problem that should be tackled, is obtaining solutions of the general covariant Dirac equation in a Schwarzschild field corresponding to ultrarelativistic electrons. It is not difficult to check that the Mathisson-Papapetrou equations in a Schwarzschild field have a strict partial solution describing the circular motion of a spinning test particle around the field source on the orbit with $`r=3m`$. The relationship between the components of the particle’s 4-velocity $`u_{}r\dot{\phi }`$ and the 3-vector spin component $`S_2S_\theta `$ is
$$u_{}=\frac{3mM}{S_\theta }$$
(21)
(as above, here we use the standard Schwarzschild coordinates; spin is perpendicular to the plane of motion $`\theta =\pi /2`$, therefore $`S_1=0`$, $`S_3=0`$). It is necessary to take into account the condition for a spinning test particle $`|S_0|/Mr1`$ where $`|S_0|`$ is the value of the spin of a test particle as measured by the comoving observer (in our case there is the relation $`|S_\theta |=ru_4|S_0|`$). Therefore, for the value $`|u_{}|`$ from (21) we have $`|u_{}|1`$, i.e. for the motion on the circular orbit with $`r=3m`$ a particle must possess the ultrarelativistic velocity, the higher the spin is smaller. Formally, at $`S_\theta =0`$ the value $`|u_{}|`$ in (21) must be infinitely large, that is the particle must move with the speed of light. This fact corresponds to the known result following from the geodesic equations in a Schwarzschild field: the circular nonisotropic geodesic orbits exist only at $`r>3m`$ and, formally, for the motion on the orbit $`r=3m`$ a test particle without spin must possess the speed of light. In practice it means that only the beam of light can move on the orbit with $`r=3m`$.
So, according to the MP Eqs. the spin of a test particle allows its ultrarelativistic motion on the circular orbit $`r=3m`$. The calculations of the gravitational spin-orbit acceleration on the orbit $`r=3m`$ according to (20), (21) give
$$|\stackrel{}{a}|=\frac{\sqrt{3}}{9m}.$$
(22)
For the quantitatively comparison, we point out that value (22) is close to the Newtonian value of the free fall acceleration for the mass $`m`$ at the distance $`r=3m`$ (in the used system of unites this Newtonian acceleration is equal to $`1/9m`$).
It is interesting that the orbit $`r=3m`$ is a common solution of the MP Eqs. at the two known variants of the auxiliary conditions for these Eqs., namely, the condition of Pirani and Tulczyjew-Dixon . Generally the solutions of the MP Eqs. at the different condition do not coincide.
It is clear that the solution of the MP Eqs. describing the orbit $`r=3m`$ in a Schwarzschild field is interesting mainly in the theoretical sense because in practice one cannot deal with a macroscopic particle moving with the ultrarelativistic velocity relatively the field source. There is much more perspective situation with the high-energy elementary particles, e.g. electrons or protons. In this connection the question arises: does the Dirac equation in a Schwarzschild field have a solution which corresponds, in the certain meaning, to the considered solution of the MP Eqs. with $`r=3m`$? For answer this question let us analyse the components of the 4-spinor $`\mathrm{\Psi }_\mu `$ which by the known procedure of separation of the variables in the Dirac equation in a Schwarzschild field (see, e.g., , Ch. 10) take the form
$$\mathrm{\Psi }_1=\frac{1}{r\sqrt{2}}R_{1/2}(r)S_{1/2}(\theta )exp[i(\sigma t+m^{}\phi )],$$
$$\mathrm{\Psi }_2=R_{+1/2}(r)S_{+1/2}(\theta )exp[i(\sigma t+m^{}\phi )],$$
$$\mathrm{\Psi }_3=R_{+1/2}(r)S_{1/2}(\theta )exp[i(\sigma t+m^{}\phi )],$$
$$\mathrm{\Psi }_4=\frac{1}{r\sqrt{2}}R_{1/2}(r)S_{+1/2}(\theta )exp[i(\sigma t+m^{}\phi )].$$
(23)
For the radial functions $`R_{+1/2}(r)`$, $`R_{1/2}(r)`$ we have the expressions
$$R_{+1/2}(r)=\frac{1}{\sqrt{r^22mr}}\psi _{+1/2}(r)exp\left(\frac{i}{2}\mathrm{arctan}\frac{Mr}{\lambda }\right),$$
$$R_{1/2}(r)=\psi _{1/2}(r)exp\left(+\frac{i}{2}\mathrm{arctan}\frac{Mr}{\lambda }\right),$$
(24)
where $`\lambda `$ is the parameter of separation of the variables depended on the orbital moment, $`\sigma `$ is the value of energy, and the functions $`\psi _{+1/2}(r)`$, $`\psi _{1/2}(r)`$ can be find from the two differential equations written in . We shall consider these Eqs. for the case $`mM1`$, that is when the Schwarzschild source is, e.g., an ordinary (not microscopic) black hole. When perfoming concrete calculations one can take into account different values of $`\sigma `$, $`\lambda `$ and investigate the corresponding quantum states. Here we consider the case when $`\sigma `$ and $`\lambda `$ are equal to the values of the energy and moment of the classical electron following from the MP Eqs. for the above considered circular orbit with $`r=3m`$. Then for the functions $`\psi _{+1/2}`$, $`\psi _{1/2}`$ we obtain the equations
$$\frac{d\psi _{+1/2}}{dx}iA\left(1\frac{2}{x}\right)^1\psi _{+1/2}+3^{3/2}\frac{A}{x}\left(1\frac{2}{x}\right)^{1/2}\psi _{1/2}=0,$$
$$\frac{d\psi _{1/2}}{dx}+iA\left(1\frac{2}{x}\right)^1\psi _{1/2}+3^{3/2}\frac{A}{x}\left(1\frac{2}{x}\right)^{1/2}\psi _{+1/2}=0,$$
(25)
where $`xr/m`$, $`A2m^3M^3/\sqrt{3}`$ (Eqs. (25) are written for the values of $`x`$ which are not in the small neighborhood of $`x=2`$). The analysis of the solutions of (25) shows that the property $`|\psi _{+1/2}|=|\psi _{1/2}|`$ takes place and the maximum value of $`|\psi _{\pm 1/2}|`$ is achieved at $`x=3`$. We stress that just the values $`|\psi _{\pm 1/2}|^2`$ together with (23), (24) determine the probability to find an electron in the certain space region because for the components of the Dirac carrent $`J^\mu `$ the relationship takes place :
$$J^\mu =\sqrt{2}[l^\mu (|\mathrm{\Psi }_1|^2+|\mathrm{\Psi }_4|^2)+n^\mu (|\mathrm{\Psi }_2|^2+|\mathrm{\Psi }_3|^2)$$
$$m^\mu (\mathrm{\Psi }_1\mathrm{\Psi }_2^{}\mathrm{\Psi }_3\mathrm{\Psi }_4^{})m^\mu (\mathrm{\Psi }_1^{}\mathrm{\Psi }_2\mathrm{\Psi }_3^{}\mathrm{\Psi }_4)],$$
(26)
where $`l^\mu ,n^\mu ,m^\mu ,m^\mu `$ are the known isotropic vectors in the Newman-Penrose formalism. Taking into account (23), (26) and the relation $`|\psi _{+1/2}|=|\psi _{1/2}|`$ it is easy to find that $`J^\phi 0`$, $`J^t0`$, whereas $`J^r=0`$, $`J^\theta =0`$. It follows that the current circulates exactly on the circle and the maximum values of $`|J^\phi |`$, $`|J^t|`$ are achieved at $`r=3m`$. The width of the peak of the curve $`|\psi _{\pm 1/2}|^2`$ decreases when $`A`$ grows, and at $`A\mathrm{}`$ we have the classical circular orbit with $`r=3m`$.
So, the considered solution of the Dirac equation in a Schwarzschild field describes the quantum state corresponding to the classical orbit with $`r=3m`$. We point out that the parameters $`\sigma `$ and $`\lambda `$ of this state are equal to the energy and moment of the classical electron on this orbit. As we stress above the circular orbit with $`r=3m`$ is an example when the gravitational ultrarelativistic spin-orbit acceleration becomes significant. Further on it is interesting to investigate other solutions of the Dirac equation which can show the role of the ultrarelativistic gravitation in the astrophysical processes.
## 6 Conclusions
Relationships (16), (17), and the conclusions drawn from them in concrete physical situations described by expressions (13), (19), (20), point unequivocally to the need for investigators to direct more attention to the gravitational interaction at ultrarelativistic relative velocities. In view of the correspondence principle, there are reasons to infer that some important relationships of gravitational interaction of classical (non-quantum) objects will, to a certain extent, hold for particles of the micro-world, where ultrarelativistic relative velocities and high energies are ubiquitous.
An important question asks whether the analyses, carried out above, might not be helpful in delving into the specifics of inclusion of the gravitational interaction into the scheme of unification of interactions. We think they might be. If only because a purely classical examination affords a deeper insight into the gravitational interaction in the micro-world at high energies.
Not less important is the need to elucidate how the entity which in classical terms is denoted as ”gravitational ultrarelativistic spin-orbit interaction” should be expressed in the scheme of second quantization. |
warning/0507/physics0507105.html | ar5iv | text | # Lancret helices
## 1 Introduction
Helical filaments are tridimensional structures commonly found in Nature. They can be seen in microscopic systems, as biomolecules , bacterial fibers and nanosprings , and in macroscopic ones, as ropes, strings and climbing plants . Usually, the axis of all these objects is modeled as a circular helix, i. e. a 3D-space curve whose mathematical geometric properties, namely the curvature, $`k_F`$, and the torsion, $`\tau _F`$, are constant . This kind of helical structure has been shown to be a static solution of the Kirchhoff rod model .
The Kirchhoff rod model has been proved to be a good framework to study the statics and dynamics of long, thin and inextensible elastic rods. Applications of the Kirchhoff model range from Biology to Engineering and, recently, to Nanoscience . In most cases, the rod or filament is considered as being homogeneous, but the case of nonhomogeneous rods have also been considered in the literature. It has been shown that nonhomogeneous Kirchhoff rods may present spatial chaos . In the case of planar rods, Domokos and collaborators have provided some rigorous results for non-uniform elasticae and for constrained Euler buckling . Deviations of the helical structure of rods due to periodic variation of the Young’s modulus were verified numerically by da Fonseca, Malta and de Aguiar . Nonhomogeneous rods subject to given boundary conditions were studied by da Fonseca and de Aguiar in . The effects of a nonhomogeneous mass distribution in the dynamics of unstable closed rods have been analyzed by Fonseca and de Aguiar . Goriely and McMillen studied the dynamics of cracking whips and Kashimoto and Shiraishi studied twisting waves in inhomogeneous rods.
The stability analysis of helical structures is of great importance in the study of the elastic behavior of filamentary systems and has been performed both experimentally and theoretically . It has been also shown that the type of instability in twisted rods strongly depends on the anisotropy of the cross section .
Here, we consider a rod with nonhomogeneous bending and twisting coefficients varying along its arclength $`s`$, $`B(s)`$ and $`C(s)`$, respectively. We are concerned with the following question: is there any helical solution for the stationary Kirchhoff equations in the case of an inhomogeneous rod ? The answer is ‘yes’ and it will be shown that the helical solution for an inhomogeneous rod with varying bending coefficient cannot be the well known circular helix, for which the curvature, $`k_F`$, and torsion, $`\tau _F`$, are constant. To this purpose, we shall derive a set of differential equations for the curvature and the torsion of the centerline of an inhomogeneous rod and then apply the condition that a space curve must satisfy to be helical: the Lancret’s theorem. We shall obtain the simplest helical solutions satisfying the Lancret’s theorem and show that they are free standing helices, i.e., helices that are not subjected to axial forces . A resulting helical structure different from the circular helix, from now on, will be called a Lancret helix.
According to the fundamental theorem for space curves , the curvature $`k_F(s)`$, and the torsion, $`\tau _F(s)`$, completely determine a space curve, but for its position in space. We shall show that the $`k_F(s)`$ and $`\tau _F(s)`$ of a Lancret helix depend directly on the bending coefficient, $`B(s)`$, an expected result since the centerline of the rod does not depend on the twisting coefficient (see for example, Neukirch and Henderson ).
Some motivations for this work are related to defects and distortions in biological molecules. These defects and distortions could be modeled as inhomogeneities along a continuous elastic rod.
In Sec. II we review the general definition of a space curve, the Frenet basis and the so-called Lancret’s theorem. In Sec. III we present the static Kirchhoff equations for an intrinsically straight rod with varying stiffness, and derive the differential equations for the curvature and torsion of the rod. In Sec. IV we use the Lancret’s theorem for obtaining helical solutions of the static Kirchhoff equations and we show that they cannot be circular helices if the bending coefficient is not constant. As illustration, we compare a homogeneous rod with two simple cases of inhomogeneous rods: (i) linear and (ii) periodic bending coefficient varying along the rod. The circular helix has a well known relation of the curvature and torsion with the radius and pitch of the helix. In Sec. V we define a function involving all these variables in such a way that for the circular helix its value is identically null. We have verified, numerically, that this function approaches zero for the inhomogeneous cases considered here. In Sec. VI we analyse the cases of null torsion (straight and planar rods). Since helical solutions of intrinsically straight rods are not dynamically stable , in Sec. VII we consider a rod with a given helical intrinsic curvature and we obtain, for this case, a helical solution of the static Kirchhoff equations similar to that of an intrinsically straight inhomogeneous rod. In Sec. VIII we summarize the main results.
## 2 Curves in space
A curve in space can be considered as a path of a particle in motion. The rectangular coordinates $`(x,y,z)`$ of the point on a curve can be expressed as function of a parameter $`u`$ inside a given interval:
$$x=x(u),y=y(u),z=z(u),u_1uu_2.$$
(1)
We define the vector $`𝐱(u)(x(u),y(u),z(u))`$. If $`u`$ is the time, $`𝐱(u)`$ represents the trajectory of a particle.
### 2.1 The Frenet frame and the Frenet-Serret equations
The vector tangent to the space curve at a given point $`P`$ is simply $`d𝐱/du`$. It is possible to show that if the arclength $`s`$ of the space curve is considered as its parameter, the tangent vector at a given point $`P`$ of the curve $`𝐱(s)`$ is a unitary vector. So, using the arclength $`s`$ to parametrize the curve, we shall denote by $`𝐭`$ its tangent vector
$$𝐭=\frac{d𝐱}{ds},$$
(2)
$`𝐭=1`$. The tangent vector $`𝐭`$ points in the direction of increasing $`s`$.
The plane defined by the points $`P_1`$, $`P_2`$ and $`P_3`$ on the curve, with $`P_2`$ and $`P_3`$ approaching $`P_1`$, is called the osculating plane of the curve at $`P_1`$ . Given a point $`P`$ on the curve, the principal normal at $`P`$ is the line, in the osculating plane at $`P`$, that is perpendicular to the tangent vector at $`P`$. The normal vector $`𝐧`$ is the unit vector associated to the principal normal (its sense may be chosen arbitrarily, provided it is continuous along the space curve).
From $`𝐭.𝐭=1`$, differentiating with respect to $`s`$ (indicated by a prime) it follows that:
$$𝐭.𝐭^{}=0,$$
(3)
so that $`𝐭`$ and $`𝐭^{}`$ are orthogonal. It is possible to show that $`𝐭^{}`$ lies in the osculating plane, consequently $`𝐭^{}`$ is in the direction of $`𝐧`$. This allows us to write
$$𝐭^{}=k_F𝐧,$$
(4)
$`k_F`$ being called the curvature of the space curve at $`P`$.
The curvature measures the rate of change of the tangent vector when moving along the curve. In order to measure the rate of change of the osculating plane, we introduce the vector normal to this plane at $`P`$: the binormal unit vector $`𝐛`$. At a point $`P`$ on the curve, $`𝐛`$ is defined in such a way that
$$𝐛=𝐭\times 𝐧.$$
(5)
The frame $`\{𝐧,𝐛,𝐭\}`$ can be taken as a new frame of reference and forms the moving trihedron of the curve. It is commonly called the Frenet frame. The rate of change of the osculating plane is expressed by the vector $`𝐛^{}`$. It is possible to show that $`𝐛^{}`$ is anti-parallel to the unit vector $`𝐧`$ . So we can write
$$𝐛^{}=\tau _F𝐧,$$
(6)
$`\tau _F`$ being called the torsion of the space curve at $`P`$.
The rate of variation of $`𝐧`$ can be obtained straightforwardly. It is given by
$$𝐧^{}=k_F𝐭+\tau _F𝐛.$$
(7)
The set of differential equations for $`\{𝐭,𝐧,𝐛\}`$ is
$$\begin{array}{c}𝐭^{}=k_F𝐧,\hfill \\ 𝐧^{}=k_F𝐭+\tau _F𝐛,\hfill \\ 𝐛^{}=\tau _F𝐧,\hfill \end{array}$$
(8)
and are known as the formulas of Frenet or the Serret-Frenet equations .
### 2.2 The Fundamental theorem of space curves
A space curve parametrized by its arclength $`s`$ is defined by a vectorial function $`𝐱(s)`$. The form of $`𝐱(s)`$ depends on the choice of the coordinate system. Nevertheless, there exists a form of characterization of a space curve given by a relation that is independent of the coordinates. This relation gives the natural equation for the curve.
$`k_F(s)`$ gives the natural equation in the case of planar curves. Indeed, if $`\phi `$ is the angle between the tangent vector of the planar curve and the $`x`$-axis of the coordinate system, it is possible to show that $`k_F=d\phi /ds`$. Since $`\mathrm{cos}(\phi )=dx/ds`$ and $`\mathrm{sin}(\phi )=dy/ds`$, knowing $`k_F(s)`$, then $`\phi (s)`$, $`x(s)`$, and $`y(s)`$ of the planar curve can be obtained immediatly:
$$\phi (s)=_{s_0}^sk_F(s)𝑑s,x(s)=_{s_0}^s\mathrm{cos}\phi (s)𝑑s,y(s)=_{s_0}^s\mathrm{sin}\phi (s)𝑑s.$$
(9)
In the case of non-planar curves, if we have two single valued continuous functions $`k_F(s)`$ and $`\tau _F(s),s>0,`$ then there exists one and only one space curve, determined but for its position in space, for which $`s`$ is the arclength, $`k_F(s)`$ the curvature, and $`\tau _F(s)`$ the torsion. It is the Fundamental theorem for space curves . The functions $`k_F(s)`$ and $`\tau _F(s)`$ provide the natural equations of the space curve.
### 2.3 Curves of constant slope: the Lancret’s theorem
A space curve $`𝐱(s)`$ is a helix if the lines tangent to $`𝐱`$ make a constant angle with a fixed direction in space (the helical axis) . Denoting by $`𝐚`$ the unit vector of this direction, a helix satisfies
$$𝐭.𝐚=\mathrm{cos}\alpha =\text{constant}.$$
(10)
Differentiating Eq. (10) with respect to $`s`$ gives $`𝐚.𝐧=0`$. Therefore $`𝐚`$ lies in the plane determined by the vectors $`𝐭`$ and $`𝐛`$:
$$𝐚=𝐭\mathrm{cos}\alpha +𝐛\mathrm{sin}\alpha .$$
(11)
Differentiating Eq. (11) with respect to $`s`$, gives
$$0=(k_F\mathrm{cos}\alpha \tau _F\mathrm{sin}\alpha )𝐧,$$
or
$$\frac{k_F}{\tau _F}=\mathrm{tan}\alpha =\text{constant.}$$
(12)
This result says that for curves of constant slope the ratio of curvature over torsion is constant. Conversely, given a regular curve for which the equation (12) is satisfied, it is possible to find a constant angle $`\alpha `$ such that
$$𝐧\left(k_F\mathrm{cos}\alpha \tau _F\mathrm{sin}\alpha \right)=0,$$
$$\frac{d}{ds}\left(𝐭\mathrm{cos}\alpha +𝐛\mathrm{sin}\alpha \right)=0,$$
implying that the vector $`𝐚=𝐭\mathrm{cos}\alpha +𝐛\mathrm{sin}\alpha `$ is the unit vector along the axis. Moreover, $`𝐚.𝐭=\mathrm{cos}\alpha `$=constant, so that the curve has constant slope. This result can be expressed as:
A necessary and suficient condition for a space curve to be a curve of constant slope (a helix) is that the ratio of curvature over torsion be constant. It is the well known Lancret’s theorem, dated of 1802 and first proved by B. de Saint Venant .
If a helical curve $`𝐱(s)`$ is projected onto the plane perpendicular to $`𝐚`$, the vector $`𝐱_1(s)`$ representing this projection is given by
$$𝐱_1(s)=𝐱(𝐱.𝐚)𝐚.$$
(13)
It is possible to show that the curvature $`k_1`$ of the projected curve is given by:
$$k_1(s)=\frac{k_F(s)}{\mathrm{sin}^2\alpha }.$$
(14)
The shape of the planar curve obtained by projecting a helical curve onto the plane perpendicular to its axis is used to characterize it. For example, the well known circular helix projects a circle onto the plane perpendicular to its axis. The spherical helix projects an arc of an epicycloid onto a plane perpendicular to its axis . The logarithmic spiral is the projection of a helical curve called conical helix .
## 3 The static Kirchhoff equations
The statics and dynamics of long and thin elastic rods are governed by the Kirchhoff rod model. In this model, the rod is divided in segments of infinitesimal thickness to which the Newton’s second law for the linear and angular momentum are applied. We derive a set of partial differential equations for the averaged forces and torques on each cross section and for a triad of vectors describing the shape of the rod. The set of PDE are completed with a linear constitutive relation between torque and twist.
The central axis of the rod, hereafter called centerline, is represented by a space curve $`𝐱`$ parametrized by the arclength $`s`$. A Frenet frame is defined for this space curve as described in the previous section. For a physical filament the use of a local basis, $`\{𝐝_1,𝐝_2,𝐝_3\}`$, to describe the rod has the advantage of taking into account the twist deformation of the filament. This local basis is defined such that $`𝐝_3`$ is the vector tangent to the centerline of the rod ($`𝐝_3=𝐭`$), and $`𝐝_1`$ and $`𝐝_2`$ lie on the cross section plane. The local basis is related to the Frenet frame $`\{𝐧,𝐛,𝐭\}`$ through
$$(𝐝_1𝐝_2𝐝_3)=(𝐧𝐛𝐭)\left(\begin{array}{ccc}\mathrm{cos}\xi & \mathrm{sin}\xi & 0\\ \mathrm{sin}\xi & \mathrm{cos}\xi & 0\\ 0& 0& 1\end{array}\right),$$
(15)
where the angle $`\xi `$ is the amount of twisting of the local basis with respect to $`𝐭`$.
In this paper, we are concerned with equilibrium solutions of the Kirchhoff model, so our study departs from the static Kirchhoff equations . In scaled variables, for intrinsically straight isotropic rods, these equations are:
$`𝐅^{}=0,`$ (16)
$`𝐌^{}=𝐅\times 𝐝_3,`$ (17)
$`𝐌=B(s)k_1𝐝_1+B(s)k_2𝐝_2+C(s)k_3𝐝_3,`$ (18)
the vectors $`𝐅`$ and $`𝐌`$ being the resultant force, and corresponding moment with respect to the centerline of the rod, respectively, at a given cross section. As in the previous section, $`s`$ is the arclength of the rod and the prime denotes differentiation with respect to $`s`$. $`k_i`$ are the components of the twist vector, $`𝐤`$, that controls the variations of the director basis along the rod through the relation
$$𝐝_i^{}=𝐤\times 𝐝_i,i=1,2,3.$$
(19)
$`k_1`$ and $`k_2`$ are related to the curvature of the centerline of the rod $`(k_F=\sqrt{k_1^2+k_2^2})`$ and $`k_3`$ is the twist density. $`B(s)`$ and $`C(s)`$ are the bending and twisting coefficients of the rod, respectively. In the case of macroscopic filaments the bending and twisting coefficients can be related to the cross section radius and the Young’s and shear moduli of the rod. Writing the force $`𝐅`$ in the director basis,
$$𝐅=f_1𝐝_1+f_2𝐝_2+f_3𝐝_3,$$
(20)
the equations (1618) give the following differential equations for the components of the force and twist vector:
$`f_1^{}f_2k_3+f_3k_2=0,`$ (21)
$`f_2^{}+f_1k_3f_3k_1=0,`$ (22)
$`f_3^{}f_1k_2+f_2k_1=0,`$ (23)
$`(B(s)k_1)^{}+(C(s)B(s))k_2k_3f_2=0,`$ (24)
$`(B(s)k_2)^{}(C(s)B(s))k_1k_3+f_1=0,`$ (25)
$`(C(s)k_3)^{}=0.`$ (26)
The equation (26) shows that the component $`M_3=C(s)k_3`$ of the moment in the director basis (also called torsional moment), is constant along the rod, consequently the twist density $`k_3`$ is inversely proportional to the twisting coefficient $`C(s)`$
$$k_3(s)=\frac{M_3}{C(s)}.$$
(27)
In order to look for helical solutions of the Eqs. (2126) the components of the twist vector $`𝐤`$ are expressed as follows:
$`k_1`$ $`=`$ $`k_F(s)\mathrm{sin}\xi ,`$ (28)
$`k_2`$ $`=`$ $`k_F(s)\mathrm{cos}\xi ,`$ (29)
$`k_3`$ $`=`$ $`\xi ^{}+\tau _F(s),`$ (30)
where $`k_F(s)`$ and $`\tau _F(s)`$ are the curvature and torsion, respectively, of the space curve that defines the centerline of the rod and $`\xi `$ is given by Eq. (15). If the rod is homogeneous, the helical solution has constant $`k_F`$ and $`\tau _F`$, and $`\xi ^{}`$ is proved to be null .
Substituting Eqs. (2830) in Eqs. (2126), extracting $`f_1`$ and $`f_2`$ from Eqs. (25) and (24), respectively, differentiating them with respect to $`s`$, and substituting in Eqs. (21), (22) and (23), gives the following set of nonlinear differential equations:
$`[M_3k_F(s)B(s)k_F(s)\tau _F(s)]^{}(B(s)k_F(s))^{}\tau _F(s)=0,`$ (31)
$`(B(s)k_F(s))^{\prime \prime }+k_F(s)\tau _F(s)[M_3B(s)\tau _F(s)]f_3(s)k_F(s)=0,`$ (32)
$`(B(s)k_F(s))^{}k_F(s)+f_3^{}(s)=0.`$ (33)
Appendix A presents the details of the derivation of Eqs. (3133).
The Eqs. (3133) for the curvature, $`k_F`$, and torsion, $`\tau _F`$ do not depend on the twisting coefficient, $`C(s)`$. Therefore, the centerline of an inhomogeneous rod does not depend on the twisting coefficient like in the case of homogeneous rods (see, for example, Eqs. (13) and (14) of Ref. ).
Langer and Singer have obtained a set of first-order ordinary differential equations for the curvature and torsion of the centerline of a homogeneous rod that contains terms proportional to $`k_F^2`$ and $`\tau _F^2`$. The Eqs. (3133) have the advantage of involving only terms linear in $`k_F`$ and $`\tau _F`$.
## 4 Helical solutions of inhomogeneous rods
In order to find helical solutions for the static Kirchhoff equations, we apply the Lancret’s theorem to the general equations (3133). We first rewrite the Lancret’s theorem in the form:
$$k_F(s)=\beta \tau _F(s),$$
(34)
with $`\beta 0`$. From Eq. (12),
$$\beta \mathrm{tan}\alpha =\text{Constant}.$$
(35)
Substituting Eq. (34) in Eq. (31) we obtain
$$\tau _F^{}(M_3B\tau _F)2\tau _F(B\tau _F)^{}=0.$$
(36)
Substituting Eq. (34) in Eq. (32) and extracting $`f_3`$, we obtain
$$f_3=\frac{(B\tau _F)^{\prime \prime }}{\tau _F}+\tau _F(M_3B\tau _F).$$
(37)
Differentiating $`f_3`$ with respect to $`s`$ and substituting in Eq. (33) we obtain the following differential equation for $`\tau _F`$:
$$\frac{(B\tau _F)^{\prime \prime \prime }}{\tau _F}\frac{(B\tau _F)^{\prime \prime }\tau _F^{}}{\tau _F^2}+(\beta ^2+1)\tau _F(B\tau _F)^{}=0,$$
(38)
where the Eq. (36) was used to simplify the above equation. One immediate solution for this differential equation is
$$(B\tau _F)^{}=0,$$
(39)
that substituted in Eq. (36) gives
$$\tau _F^{}(M_3B\tau _F)=0.$$
(40)
For non-constant $`\tau _F`$, the Eq. (40) gives the following solution for $`\tau _F`$:
$$\tau _F(s)=\frac{M_3}{B(s)}.$$
(41)
Substituting the Eqs. (39) and (41) in Eq. (37) we obtain that
$$f_3(s)=0.$$
(42)
Substituting Eq. (41) in (34) we obtain:
$$k_F(s)=\beta \frac{M_3}{B(s)}.$$
(43)
Substituting Eq. (15) in Eq. (20), the force $`𝐅`$ becomes
$$𝐅=(f_1\mathrm{cos}\xi f_2\mathrm{sin}\xi )𝐧+(f_1\mathrm{sin}\xi +f_2\mathrm{cos}\xi )𝐛+f_3𝐭,$$
(44)
where $`\{𝐧,𝐛,𝐭\}`$ is the Frenet basis. Using the Eqs. (88) and (89) for $`f_1`$ and $`f_2`$ (Appendix A), we obtain
$$𝐅=(Bk_F)^{}𝐧+k_F[M_3B\tau _F]𝐛+f_3𝐭,$$
(45)
where $`f_3`$, in the inhomogeneous case, must satisfy the Eq. (33).
Substituting the Eqs. (4143) in the Eq. (45), and using Eq. (35), it follows $`𝐅=0`$. Therefore, the helical solutions satisfying (39) are free standing.
Now, we prove that a circular helix cannot be a solution of the static Kirchhoff equations for a rod with varying bending stiffness. If a helix is circular, $`k_F^{}=0`$ and $`\tau _F^{}=0`$, and from Eq. (31) we obtain:
$$2k_F\tau _FB^{}=0.$$
(46)
Since $`B^{}(s)0`$, Eq. (46) will be satisfied only if $`k_F=0`$ and/or $`\tau _F=0`$. Therefore, it is not possible to have a circular helix as a solution for a rod with varying bending coefficient.
The solutions for the curvature $`k_F`$, Eq. (43), and the torsion $`\tau _F`$, Eq. (41), can be used to obtain the unit vectors of the Frenet frame (by integration of the Eqs. (8)). From Eqs. (43), (35) and (11), we can obtain $`\alpha `$ and $`𝐚`$ once $`k_F(0)`$, $`M_3`$ and $`B(s)`$ are given. By choosing the $`z`$-direction of the fixed cartesian basis as the direction of the unit vector $`𝐚`$, we can integrate $`𝐭`$ in order to obtain the three-dimensional configuration of the centerline of the rod.
Figure 1 displays the helical solution of the static Kirchhoff equations for rods with bending coefficients given by
$`\text{Fig }\text{1}\text{a:}B_a(s)`$ $`=`$ $`1,`$ (47)
$`\text{Fig }\text{1}\text{b:}B_b(s)`$ $`=`$ $`1+0.007s,`$ (48)
$`\text{Fig }\text{1}\text{c:}B_c(s)`$ $`=`$ $`1+0.1\mathrm{sin}(0.04s+2).`$ (49)
The case of constant bending (47) produces the well known circular helix displayed in Fig. 1a. Figs. 1b–1c show that non-constant bending coefficients (Eqs. (4849)) do not produce a circular helix.
The helical solutions displayed in Fig. 1 satisfy the Lancret’s theorem, Eq. (12). The tridimensional helical configurations displayed in Fig. 1 were obtained by integrating the Frenet-Serret equations (8) using the following initial conditions for the Frenet frame: $`𝐭(s=0)=(0,\mathrm{sin}\alpha ,\mathrm{cos}\alpha )`$, $`𝐧(s=0)=`$ $`(1,0,0)`$ and $`𝐛(s=0)=`$ $`(0,\mathrm{cos}\alpha ,\mathrm{sin}\alpha )`$. This choice ensures that the $`z`$-axis is parallel to the direction of the helical axis, vector $`𝐚`$. The centerline of the helical rod is a space curve $`𝐱(s)=(x(s),y(s),z(s))`$ that is obtained by integration of the tangent vector $`𝐭(s)`$. We have taken the helical axis as the $`z`$-axis and placed the initial position of the rod at $`x(0)=1/k_1(0)`$, $`y(0)=0`$ and $`z(0)=0`$ (in scaled units), where $`k_1(0)`$ is the curvature of the planar curve at $`s=0`$ obtained by projecting the space curve onto the plane perpendicular to the helical axis (Eq. (14)). From Eq. (14) we have
$$k_1(0)=\frac{k_F(0)}{\mathrm{sin}^2\alpha }.$$
(50)
Using the Eq. (35), it follows that
$$\mathrm{sin}^2\alpha =\frac{\beta ^2}{1+\beta ^2}.$$
(51)
From Eq. (43), setting $`s=0`$, we get
$$\beta =\frac{k_F(0)B(0)}{M_3}.$$
(52)
Substituting Eqs. (51) and (52) in Eq. (50), we obtain:
$$x(0)=\frac{1}{k_1(0)}=\frac{k_F(0)B^2(0)}{M_3^2+k_F^2(0)B^2(0)}.$$
(53)
$`k_F(0)`$ and $`M_3`$ are free parameters that have been chosen so that the helical solutions displayed in Fig. 1 have the same angle $`\alpha `$. The parameters $`k_F(0)=0.24`$ and $`M_3=0.05`$ give $`x(0)4`$ for the helical solutions displayed in Figs. 1a and 1b, and the parameters $`k_F(0)=0.22`$ and $`M_3=0.05`$ give $`x(0)4.36`$ for the helical solution displayed in the Fig. 1c.
For short the projection of the space curve onto the plane perpendicular to the helical axis will be called projected curve. As mentioned in Sec. II, the circle is the projected curve of the most common type of helix, the circular helix.
Fig. 2 displays the projected curves related to the helical solutions displayed in Fig. 1. Fig. 2a shows that the helical solution of the inhomogeneous rod with constant bending coefficient projects a circle onto the plane perpendicular to the helical axis.
If required, the natural equations for the projected curves displayed in Fig. 2 are easily obtained, for instance, by substitution of the solution (Eq. (43)) for the curvature $`k_F(s)`$ of the helical rod into Eq. (14). The natural equation of the projected curve is given by its curvature,
$$k_1(s)=\frac{\beta M_3}{B(s)}\mathrm{sin}^2\alpha ,$$
(54)
where $`\beta `$ and $`\mathrm{sin}^2\alpha `$ can be obtained by Eqs. (51) and (52). Then, in the Eq.(54), setting $`B(s)=B_i(s)`$, $`i=a,b,c`$, as given in Eqs. (4749), produces the natural equation for the corresponding projected curve displayed in Fig. 2. The helical rod displayed in Fig. 1b is a conical helix since the radius of curvature of its projected curve (inverse of $`k_1(s)`$) is a logarithmic spiral ($`1/k_1(s)`$ is a linear function of $`s`$ ).
From Eqs. (2830), (27) and (41) we obtain the variation of the angle $`\xi `$ between the local basis, $`𝐝_i`$, $`i=1,2,3`$, and the Frenet frame, $`\{𝐧,𝐛,𝐭\}`$:
$$\xi ^{}=k_3(s)\tau _F(s)=\left(\frac{M_3}{C(s)}\frac{M_3}{B(s)}\right)=M_3\frac{B(s)C(s)}{B(s)C(s)}.$$
(55)
Eq. (55) shows that $`\xi ^{}0`$ for the general case of $`B(s)C(s)`$, i. e. helical filaments corresponding to inhomogeneous rods are not twistless. The circular helix is a helicoidal solution for the centerline of an inhomogeneous rod having constant bending coefficient. We emphasize that the inhomogeneous rod is not twistless in contrast with the homogeneous case where it has been proved that $`\xi ^{}=0`$ .
A homogeneous rod has $`B(s)`$ and $`C(s)`$ constant so that $`k_3=`$ Constant (from Eq. (26)). Since $`\xi ^{}`$ has been proved to be null for a helical solution of a homogeneous rod (see reference ), Eq. (30) shows that the torsion $`\tau _F`$ must be a constant. In order to satisfy the Lancret’s theorem (Eq. (12)) the curvature $`k_F`$ of the helical solution must also be a constant. Therefore, the only type of helical solution for a homogeneous rod is the circular helix, while an inhomogeneous rod may present other types of helical structures.
## 5 Radius and Pitch of the helical solution
The radius $``$ of a helix is defined as being the distance of the space curve to its axis. The pitch $`𝒫`$ of a helix is defined as the height of one helical turn, i.e., the distance along the helical axis of the initial and final points of one helical turn.
For a circular helix, $``$, $`𝒫`$, $`k_F`$ and $`\tau _F`$ are constant, and it is easy to prove that
$$\lambda =\left(\sqrt{^2+𝒫^2/(4\pi ^2)}\right)^1=\sqrt{k_F^2+\tau _F^2}.$$
(56)
For other types of helix, it constitutes a very hard problem in differential geometry to obtain the relation between the curvature $`k_F`$ and the torsion $`\tau _F`$ with the radius $``$ and the pitch $`𝒫`$. We have seen in Sec. II that the definitions of curvature and torsion involve the calculation of the modulus of the tangent and normal vectors derivative with respect to the arclength of the rod. We also saw that the Frenet-Serret differential equations for the Frenet frame depend on the curvature and torsion. The difficulty of integration of the Frenet-Serret equations for the general case where $`k_F(s)`$ and $`\tau _F(s)`$ are general functions of $`s`$ poses the problem of finding an analytical solution for the centerline of the rod, thus the difficulty of relating non constant curvature and torsion with non constant radius and pitch. Due to this difficulty we shall test the possibility of generalizing the relation (56) to the present inhomogeneous case. In order to do so, from the equation (56) we define:
$$g_\lambda (s)\left(\sqrt{^2(s)+𝒫^2(s)/(4\pi ^2)}\right)^1\sqrt{k_F^2(s)+\tau _F^2(s)},$$
(57)
where $`(s)`$ and $`𝒫(s)`$ are the radius and the pitch of the helical structure as function of $`s`$. In the case of a circular helix, from Eq. (56), $`g_\lambda (s)=0`$ for all $`s`$.
Since the $`z`$-axis is defined as being the axis of the helical solution we can calculate the radius $`(s)`$ through: $`(s)=\sqrt{x^2(s)+y^2(s)}`$, where $`x(s)`$ and $`y(s)`$ are the $`x`$ and $`y`$ components of the vector position of the centerline of the helical rod.
The pitch of the helix is the difference between the $`z`$-coordinate of the initial and final positions of one helical turn. A helical turn can be defined such that the projection of the vector position of the spatial curve along the $`xy`$-plane (vector $`𝐱_1`$ of Eq. (13)), rotates of $`2\pi `$ around the $`z`$-axis.
Fig. 3 shows $`g_\lambda (s)`$ for the free standing helix of Fig. 1b. We see that $`g_\lambda (s)`$ oscillates, its maximum amplitude being smaller than $`0.006`$. For the helical shape displayed in Fig. 1c we found that the maximum value of $`g_\lambda (s)`$ is smaller than $`0.008`$ (data not shown). While for a circular helix $`g_\lambda =0`$, for the free standing helices displayed in Fig. 1b and Fig. 1c the function $`g_\lambda `$ oscillates around zero with small amplitude.
The small amplitude of these oscillations suggests that the relations $`(s)k_F(s)[\sqrt{k_F^2(s)+\tau _F^2(s)}]^1`$ and $`𝒫(s)2\pi \tau _F(s)[\sqrt{k_F^2(s)+\tau _F^2(s)}]^1`$, valid for circular helices, could be used to derive approximate functions for the radius and the pitch of different types of helical structures, but the oscillatory behavior indicates that these relations are not simple functions of the geometric features of the helix.
## 6 Straight and planar inhomogeneous rods
Straight rods ($`k_F=0`$), and planar rods ($`k_F0`$), have null torsion ($`\tau _F=0`$), and constitute particular cases of helices. In both cases there is at least one direction in space that makes a constant angle $`\alpha =\pi /2`$ with the vector tangent to the rod centerline.
The straight inhomogeneous rod is a solution of the static Kirchhoff equations that has non-constant twist density (Eq. (27)), in contrast with the homogeneous case for which the twist density is constant.
The twisted planar ring ($`k_F=`$ Constant) is a solution of the static Kirchhoff equations only if the bending coefficient can be written in the form:
$$B(s)=A_0\mathrm{cos}(k_Fs)+B_0\mathrm{sin}(k_Fs)+C_I/k_F^2,$$
(58)
with $`A_0`$, $`B_0`$ and $`C_I`$ constant. If $`k_F`$ is function of $`s`$ (instead of being a constant) there exist no solutions for Eqs. (2126). So, the existence of a planar solution related to the general form of the components of the twist vector given by equations (2830) requires $`k_F=`$ Constant.
## 7 Helical structure with intrinsic curvature
The helical shape displayed in Fig. 1b resembles that exhibited by the tendrils of some climbing plants. In these plants the younger parts have smaller cross section diameter, giving rise to non-constant bending coefficient. The main difference between the solution displayed in Fig. 1b and the tendrils of climbing plants is that the solution in Fig. 1b was obtained for an intrinsically straight rod while the tendrils have intrinsic curvature .
The tendrils of climbing plants are stable structures while the helical solution displayed in Fig. 1b is not stable because the rod is intrinsically straight . We shall show that a rod with intrinsic curvature may have a static solution of the Kirchhoff equations similar to that displayed in Fig. 1b.
The intrinsic curvature of a rod is introduced in the Kirchhoff model through the components of the twist vector, $`𝐤^{(0)}`$, in the unstressed configuration of the rod as
$`k_1^{(0)}`$ $`=`$ $`k_F^{(0)}(s)\mathrm{sin}\xi ,`$ (59)
$`k_2^{(0)}`$ $`=`$ $`k_F^{(0)}(s)\mathrm{cos}\xi ,`$ (60)
$`k_3^{(0)}`$ $`=`$ $`\xi ^{}+\tau _F^{(0)}(s),`$ (61)
where $`k_F^{(0)}(s)`$ and $`\tau _F^{(0)}(s)`$ are the curvature and torsion of the space curve that represents the axis of the rod in its unstressed configuration, simply called intrinsic curvature of the rod. We consider that the unstressed configuration of the axis of the rod forms a helical space curve with the intrinsic curvature satisfying
$`B(s)k_F^{(0)}(s)=K_0,`$ (62)
$`B(s)\tau _F^{(0)}(s)=T_0,`$ (63)
where $`K_0`$ and $`T_0`$ are constant and $`B(s)`$ is the bending coefficient of the rod.
The linear constitutive relation (Eq. (18)) becomes
$$𝐌=B(s)(k_1k_1^{(0)})𝐝_1+B(s)(k_2k_2^{(0)})𝐝_2+C(s)(k_3k_3^{(0)})𝐝_3,$$
(64)
where $`C(s)`$ is the twisting coefficient of the rod. The static Kirchhoff equations for this case, Eqs. (16), (17) and Eq. (64), are given by
$`f_1^{}f_2k_3+f_3k_2=0,`$ (65)
$`f_2^{}+f_1k_3f_3k_1=0,`$ (66)
$`f_3^{}f_1k_2+f_2k_1=0,`$ (67)
$`(B(s)(k_1k_1^{(0)}))^{}B(s)(k_2k_2^{(0)})k_3+C(s)(k_3k_3^{(0)})f_2=0,`$ (68)
$`(B(s)(k_2k_2^{(0)}))^{}+B(s)(k_1k_1^{(0)})k_3C(s)(k_3k_3^{(0)})+f_1=0,`$ (69)
$`(C(s)(k_3k_3^{(0)}))^{}+B(s)(k_1^{(0)}k_2k_2^{(0)}k_1)=0.`$ (70)
The components of the twist vector are expressed as:
$`k_1`$ $`=`$ $`k_F(s)\mathrm{sin}\chi ,`$ (71)
$`k_2`$ $`=`$ $`k_F(s)\mathrm{cos}\chi ,`$ (72)
$`k_3`$ $`=`$ $`\chi ^{}+\tau _F(s).`$ (73)
In order to obtain the simplest solution for the static Kirchhoff equations Eqs. (6570) with the intrinsic curvature given by Eqs. (5961) and (6263) we shall look for a solution such that $`\chi =\xi `$ in Eqs. (7173). This solution preserves the intrinsic twist density of the helical structure. In this case, the Eq. (70) becomes simply $`[C(s)(\tau _F\tau _F^{(0)})]^{}=0`$ or
$$M_3=C(s)(\tau _F\tau _F^{(0)})=\text{Constant},$$
(74)
and we obtain the following differential equations for the curvature $`k_F(s)`$, and the torsion $`\tau _F(s)`$, of the rod:
$$\begin{array}{cc}[M_3k_FB\tau _F(k_Fk_F^{(0)})]^{}[B(k_Fk_F^{(0)})]^{}\tau _F=0,\hfill & \\ [B(k_Fk_F^{(0)})]^{\prime \prime }+\tau _F[M_3k_FB\tau _F(k_Fk_F^{(0)})]f_3k_F=0,\hfill & \\ [B(k_Fk_F^{(0)})]^{}k_F+f_3^{}=0,\hfill & \end{array}$$
(75)
where we have omitted the dependence on $`s`$ to simplify the notation. In order to obtain a helical solution of these equations we apply the Lancret’s theorem, Eq. (12), to the Eqs. (75). We obtain the following results:
$`f_3(s)=0,`$ (76)
$`[B(s)(k_F(s)k_F^{(0)}(s))]^{}=0k_F(s)k_F^{(0)}(s)={\displaystyle \frac{K}{B(s)}},`$ (77)
$`[B(s)(\tau _F(s)\tau _F^{(0)}(s))]^{}=0\tau _F(s)\tau _F^{(0)}(s)={\displaystyle \frac{T}{B(s)}},`$ (78)
where $`K`$ and $`T`$ are integration constants. From Eqs. (74) and Eq. (78) we obtain
$$T=\frac{B(s)}{C(s)}M_3,$$
(79)
so that the ratio $`B(s)/C(s)`$ has to be constant.
From Eqs. (77), (78), (62), (63) and (12) we have
$$\frac{k_F(s)}{\tau _F(s)}=\frac{K+K_0}{T+T_0}=\mathrm{tan}\alpha .$$
(80)
From Eqs. (68), (69), (77) and (78) it follows that
$$\begin{array}{c}f_1=0,\hfill \\ f_2=0.\hfill \end{array}$$
(81)
Therefore, the helical solution given by Eqs. (7678) (obtained imposing $`\chi =\xi `$) is a free standing helix ($`𝐅=(f_1,f_2,f_3)=0`$).
It follows from Eqs. (7678), (62) and (63) that the solutions for the curvature $`k_F(s)`$, and the torsion $`\tau _F(s)`$, of the rod with helical intrinsic curvature are similar to those of intrinsically straight rods, Eqs. (4143). Therefore, rods with intrinsic curvature and a non-constant bending coefficient given by Eq. (48) (Eq. (49)) can have a three-dimensional configuration similar to that displayed in Fig. 1b (Fig. 1c).
## 8 Conclusions
The existence of helical configurations for a rod with non-constant stiffness has been investigated within the framework of the Kirchhoff rod model. Climbing and spiralling solutions of planar rods have been studied by Holmes et. al. . Here, we have shown that helical spiralling three-dimensional structures are possible solutions of the static Kirchhoff equations for an inhomogeneous rod.
From the static Kirchhoff equations, we derived the set of differential equations (3133) for the curvature and the torsion of the centerline of a rod whose bending coefficient is a function of the arclength $`s`$. We have shown that the circular helix is the type of helical solution obtained when $`B(s)`$ is constant, independently of the rod being homogeneous or inhomogeneous.
Though the differential equations for the curvature and torsion are general, we have obtained only the simplest helical solutions (Eqs. (39) and (4143)), obtained when the Lancret’s theorem is applied to the differential equations. We show that these solutions are free standing and that the curvature and torsion depend directly on the form of variation of the bending coefficient. Figures 1b and 1c are examples of helical solutions of inhomogeneneous rods whose bending coefficients are given by Eqs. (48) and (49). The helical structure displayed in Fig. 1b is a conical helix since the projected curve onto the plane perpendicular to the helical axis is a logarithmic spiral, i. e., $`1/k_1(s)`$ is a linear function of $`s`$ .
In the particular case of an inhomogeneous rod with the intrinsic curvature defined by Eqs. (5961) and (6263), with $`B(s)/C(s)`$ constant, we also obtain the helical solutions displayed in Figs. 1b and 1c. The tendrils of some climbing plants present a three-dimensional structure similar to that displayed in Fig. 1b. In these plants, the cross-section diameter of the tendrils varies along them, giving rise to non-constant bending coefficient, and the differential growth of the tendrils produces intrinsic curvature . The bending and twisting coefficients of a continuous filament with circular cross-section are proportional to its moment of inertia $`I`$. It implies that $`B(s)/C(s)`$ is constant for an inhomogeneous rod. Therefore, the tendrils of climbing plants can be well described by the Kirchhoff model for an inhomogeneous rod with a linear variation of the bending stiffness.
## Acknowledgements
This work was partially supported by the Brazilian agencies FAPESP, CNPq and CAPES. The authors would like to thank Prof. Manfredo do Carmo for valuable informations about the Lancret’s theorem.
## Appendix A Appendix: The differential equations for the curvature and torsion
Here, we shall derive the Eqs. (3133). Substitution of Eqs. (2830) into Eqs. (2126) gives:
$`f_1^{}f_2(\xi ^{}+\tau _F)+f_3k_F\mathrm{cos}\xi =0,`$ (82)
$`f_2^{}+f_1(\xi ^{}+\tau _F)f_3k_F\mathrm{sin}\xi =0,`$ (83)
$`f_3^{}f_1k_F\mathrm{cos}\xi +f_2k_F\mathrm{sin}\xi =0,`$ (84)
$`(B(s)k_F\mathrm{sin}\xi )^{}+(C(s)B(s))k_F\mathrm{cos}\xi (\xi ^{}+\tau _F)f_2=0,`$ (85)
$`(B(s)k_F\mathrm{cos}\xi )^{}(C(s)B(s))k_F\mathrm{sin}\xi (\xi ^{}+\tau _F)+f_1=0,`$ (86)
$`(C(s)(\xi ^{}+\tau _F))^{}=0.`$ (87)
First, we extract $`f_1`$ and $`f_2`$ from Eqs. (86) and (85), respectively:
$`f_1=(B(s)k_F)^{}\mathrm{cos}\xi +[M_3k_FB(s)k_F\tau _F]\mathrm{sin}\xi ,`$ (88)
$`f_2=(B(s)k_F)^{}\mathrm{sin}\xi +[M_3k_FB(s)k_F\tau _F]\mathrm{cos}\xi ,`$ (89)
where $`M_3=C(s)(\xi ^{}+\tau _F)`$ is the torsional moment of the rod that is constant by Eq. (87). Differentiating $`f_1`$ and $`f_2`$ with respect to $`s`$, substituting in Eqs. (82) and (83), respectively, and using Eqs. (88) and (89), gives the following equations:
$$\begin{array}{cc}\left\{(B(s)k_F)^{\prime \prime }\tau _F[M_3k_FB(s)k_F\tau _F]+f_3k_F\right\}\mathrm{cos}\xi +\hfill & \\ \left\{[M_3k_FB(s)k_F\tau _F]^{}\tau _F(B(s)k_F)^{}\right\}\mathrm{sin}\xi =0,\hfill & \end{array}$$
(90)
$$\begin{array}{cc}\left\{(B(s)k_F)^{\prime \prime }+\tau _F[M_3k_FB(s)k_F\tau _F]f_3k_F\right\}\mathrm{sin}\xi +\hfill & \\ \left\{[M_3k_FB(s)k_F\tau _F]^{}\tau _F(B(s)k_F)^{}\right\}\mathrm{cos}\xi =0.\hfill & \end{array}$$
(91)
Multiplying Eq. (90) (Eq. (91)) by $`\mathrm{sin}\xi `$ ($`\mathrm{cos}\xi `$) and then adding the resulting equations, we obtain the Eq. (31) for the curvature and torsion:
$$[M_3k_FBk_F\tau _F]^{}(Bk_F)^{}\tau _F=0.$$
(92)
Multiplying Eq. (90) (Eq. (91)) by $`\mathrm{cos}\xi `$ ($`+\mathrm{sin}\xi `$) and then adding the resulting equations, we obtain the Eq. (32):
$$(Bk_F)^{\prime \prime }+k_F\tau _F(M_3B\tau _F)f_3k_F=0.$$
(93)
Finally, the Eq. (33) is obtained by substituting Eqs. (88) and (89) in Eq. (84):
$$(Bk_F)^{}+f_3^{}=0.$$
(94) |
warning/0507/cond-mat0507112.html | ar5iv | text | # Critical aging of a ferromagnetic system from a completely ordered state
## Abstract
We adapt the non-linear $`\sigma `$ model to study the nonequilibrium critical dynamics of $`O(n)`$ symmetric ferromagnetic system. Using the renormalization group analysis in $`d=2+\epsilon `$ dimensions we investigate the pure relaxation of the system starting from a completely ordered state. We find that the average magnetization obeys the long-time scaling behavior almost immediately after the system starts to evolve while the correlation and response functions demonstrate scaling behavior which is typical for aging phenomena. The corresponding fluctuation-dissipation ratio is computed to first order in $`\epsilon `$ and the relation between transverse and longitudinal fluctuations is discussed.
Aging phenomena have been found in a broad variety of strongly disordered systems such as polymer and spin glasses hodge-95 , electronic system of an Anderson insulator vaknin-00 , array of flux lines pinned by disorder schehr-04 , etc. Recently much attention has been attracted by aging of pure (or weakly disordered) systems with slow relaxation dynamics governed, for example, by domain growth bray-94 as in a ferromagnetic system below the critical temperature $`T_c`$ or by critical slowing down as in a ferromagnetic system exactly at criticality godreche-00 . One expects that the critical aging phenomena can be cast into different universality classes of nonequilibrium critical dynamics. Most of studies consider the relaxation of a ferromagnetic system starting from a completely disordered state after quenching it to the fixed temperature $`TT_c`$. It was found that the response function $`R(t,s)`$ and the correlation function $`C(t,s)`$ depend nontrivially on the ratio $`x=t/s`$ similar to that found in glassy systems. Here $`s`$ and $`t`$ are waiting and observation times respectively. The distance from equilibrium can be measured by the fluctuation-dissipation ratio (FDR) $`X(t,s)=TR(t,s)/_sC(t,s)`$. It has been argued that for critical aging the limit $`X^{\mathrm{}}=lim_{t,s\mathrm{}}X(t,s)`$ is a novel universal quantity of critical phenomena godreche-00 . The FDR was computed for the $`d`$ dimensional spherical model godreche-00 , $`O(n)`$ symmetric ferromagnetic model calabrese-02 and diluted spin models calabrese-02b ; schehr-05 ; chen-05 . In all these systems $`X^{\mathrm{}}`$ has values ranging between $`0`$ and $`1/2`$. The mean field calculations show that the aging behavior is modified in the presence of long-range correlations in the initial disordered state newman-90 ; picone-02 . Much less known about the relaxation starting from an ordered state. The numerical simulations show that the correlation function demonstrates behavior which is typical for aging phenomena schehr-05 ; berthier-01 , while the magnetization obeys the long-time scaling behavior almost immediately after the system starts to evolve jaster-99 ; li-96 . This observation was used to develop new effective numerical methods to determine the critical exponents, which are based on the short-time critical dynamics and do not require the time-consuming equilibration of the system zheng-96 . However, up to now there is no any theoretical explanation why the long-time scaling behavior emerges already in the macroscopically early initial stage of relaxation and there is no theoretical framework which allows one to take properly into account the critical fluctuations in this regime of aging.
In this paper we study the nonequilibrium critical dynamics of a ferromagnetic system starting from a completely ordered state. The long-distance properties of the $`O(n)`$ symmetric system below the transition point can be related to the non-linear $`\sigma `$ model defined by the reduced Hamiltonian polyakov-75
$$=d^dx\left[\frac{1}{2}(𝐬)^2𝐡𝐬\right],|𝐬(x)|^2=1.$$
(1)
In principle, to describe properly the dynamics of a real isotropic magnet one has to consider the Larmor precession of the conserved field $`𝐬`$ in the local magnetic field $`𝐡`$, i.e. construct the low-temperature version of the model J trimper-78 . However in this work, which to our knowledge is the first analytical study of the critical aging of a finite range system from a completely ordered state, we are interested in the influence of the initial condition rather than in the complications caused by more realistic dynamic model. We restrict consideration to the pure relaxational dynamics (model A) with the non-conserved $`n`$-component order parameter $`𝐬(x)`$. It can be described by the Langevin equation $`_t𝐬=\lambda \delta /\delta 𝐬+𝜻,`$ where $`\lambda `$ is an Onsager coefficient and $`𝜻`$ is a random thermal noise dominicis-77 ; zinn-justin ; bausch-80 . Because of the constraint $`𝐬^2=1`$ the random noise can act only in the tangential direction and thus is $`𝐬`$ \- dependent. This dependence complicates the explicit expression for thermal noise distribution, however, one can simply formulate the dynamics in terms of a generating functional bausch-80 . The corresponding weight reads
$$P[𝐬,\widehat{𝐬}]=\delta (𝐬^21)\delta (𝐬\widehat{𝐬})\mathrm{exp}(𝒜),$$
(2)
where $`\widehat{𝐬}`$ is the response field and the action is given by
$$𝒜[𝐬,\widehat{𝐬}]=\frac{1}{T}d^dx𝑑t\left(\lambda \widehat{𝐬}^2+i\widehat{𝐬}\left[_t𝐬+\lambda \frac{\delta }{\delta 𝐬}\right]\right).$$
(3)
Note that $`\widehat{𝐬}`$ is a real variable which is connected with the corresponding quantity $`\stackrel{~}{𝐬}`$ introduced in Refs. calabrese-02 ; calabrese-02b ; schehr-05 ; chen-05 by $`\stackrel{~}{𝐬}=i\widehat{𝐬}`$. The first $`\delta `$-function in Eq. (2) is due to the constraint $`𝐬^2=1`$ and the additional constraint $`𝐬\widehat{𝐬}=0`$, imposed by the second $`\delta `$-function, ensures that the thermal noise $`𝜻`$ acts only in the tangential direction bausch-80 .
The effect of a macroscopic initial condition can be taken into account by averaging over the initial configurations $`𝐬_0(x)𝐬(t=0,x)`$ with a weight $`P[𝐬_0]=\mathrm{exp}(_0[𝐬_0]/T)`$ janssen-89 . Taking $`_0=d^dx𝐡_0𝐬_0(x)`$ we specify an initial condition that corresponds to the equilibrium state of the system subjected to an external magnetic field $`𝐡_0`$, which is switched off at $`t=0`$. There is a deep analogy between the considered model and the renormalization group (RG) description of the surface critical phenomena using expansion about the low critical dimension diehl-86 . We assume that the magnetic fields $`𝐡_0`$ and $`𝐡`$ act along the one direction which we will refer as to longitudinal one. We decompose the fields $`𝐬=(\sigma ,𝝅)`$ and $`\widehat{𝐬}=(\widehat{\sigma },\widehat{𝝅})`$ into $`(n1)`$-component transverse parts $`𝝅`$ and $`\widehat{𝝅}`$ and longitudinal parts $`\sigma =\sqrt{1\pi ^2}`$ and $`\widehat{\sigma }=𝝅\widehat{𝝅}/\sqrt{1\pi ^2}`$, where the constraints on the fields have been used. We now can write the weight functional $`\mathrm{exp}(𝒜_{\mathrm{tr}}[𝝅,\widehat{𝝅}])`$ for the transverse components only with the action given by
$`𝒜_{\mathrm{tr}}[𝝅,\widehat{𝝅}]={\displaystyle \frac{1}{T}}{\displaystyle }d^dx{\displaystyle _0^{\mathrm{}}}dt\{\lambda \widehat{𝝅}^2+i\widehat{𝝅}(_t\lambda ^2)𝝅`$
$`+\lambda {\displaystyle \frac{(𝝅\widehat{𝝅})^2}{1𝝅^2}}+{\displaystyle \frac{i}{2}}{\displaystyle \frac{𝝅\widehat{𝝅}}{1𝝅^2}}[(_t\lambda ^2)𝝅^2{\displaystyle \frac{\lambda }{2}}{\displaystyle \frac{(𝝅^2)^2}{1𝝅^2}}`$
$`+2\lambda h\sqrt{1𝝅^2}]\}{\displaystyle \frac{h_0}{T}}{\displaystyle }d^dx\sqrt{1𝝅_0^2}.`$ (4)
We will adopt the dimensional regularization scheme zinn-justin in which the terms generated by the measure in the path integral over $`𝐬`$ and $`\widehat{𝐬}`$ vanish so that we have omitted these terms in Eq. (4) from the beginning. Let us introduce the connected Green functions $`G_{k\widehat{k}}^{l\widehat{l}}:=[𝝅]^k[\widehat{𝝅}]^{\widehat{k}}[𝝅_0]^l[\widehat{𝝅}_0]^{\widehat{l}}`$, the low temperature expansions of which can be obtained with a loop expansion based on the action (4). The terms quadratic in $`𝝅`$ and $`\widehat{𝝅}`$ give us the free response function and the free correlator:
$`R_q^0(t,s)`$ $`=`$ $`i\mathrm{\Theta }(ts)G_q(ts),`$ (5)
$`C_q^0(t,s)`$ $`=`$ $`{\displaystyle \frac{G_q(|ts|)}{h+q^2}}+\left(h_0^1{\displaystyle \frac{1}{h+q^2}}\right)G_q(t+s),`$ (6)
where the notation $`G_q(t):=Te^{\lambda (h+q^2)t}`$ is introduced. The infinite number of higher order terms in the expansion of square roots in Eq. (4) will be treated as interactions. In each order in $`T`$ we have to take into account only a fixed number of such terms. The completely ordered initial state corresponds to the limit $`h_0\mathrm{}`$ in which Eq. (6) becomes a Dirichlet correlator. Although the action has infinite number of vertices proportional to $`h_0`$, which are located at time surface $`t=0`$, it is easy to show by direct inspection of Feynman diagrams that they do not contribute in the limit $`h_0\mathrm{}`$. The finite $`h_0^1`$ can be treated then as an additional perturbation. The great advantage of the considered nonequilibrium model in comparison with the equilibrium counterpart is that the Green functions are not spoiled by ir singularities, so that we can probe the critical domain just taking the limit $`h0`$ without using the less convenient matching procedure bausch-80 ; nelson-77 . However, the theory suffers of the uv divergences which can be converted into poles in $`\epsilon =d2`$ using dimensional regularization. Exploiting the ideas of Refs. polyakov-75 and zinn-justin one can prove the renormalizability of the model, which means that all uv divergences can be absorbed into finite number of Z-factors according to
$`\stackrel{̊}{𝝅}=Z^{1/2}𝝅,\stackrel{̊}{\widehat{𝝅}}=\widehat{Z}Z^{1/2}\widehat{𝝅},\stackrel{̊}{𝝅}_0=(ZZ_0)^{1/2}𝝅_0,`$
$`\stackrel{̊}{\widehat{𝝅}}_0=\widehat{Z}Z^{1/2}Z_0^{1/2}\widehat{𝝅}_0,\stackrel{̊}{\lambda }=Z\widehat{Z}^1\lambda ,\stackrel{̊}{T}=\mu ^\epsilon K_dZ_TT,`$
$`\stackrel{̊}{h}=\mu ^2Z_TZ^{1/2}h,\stackrel{̊}{h}_0=\mu ^2Z_T(ZZ_0)^{1/2}h_0.`$ (7)
Here circles denote the bare quantities, $`(2\pi )^dK_d`$ is the surface area of a $`d`$-dimensional unit sphere, and $`\mu `$ is an arbitrary momentum scale. The Z-factors except for $`Z_0`$ are the same as in equilibrium dominicis-77 ; zinn-justin and to order $`T`$ given by $`Z=\widehat{Z}=1+(n1)T/\epsilon `$ and $`Z_T=1+(n2)T/\epsilon `$. The new factor $`Z_0=1+(n3)T/\epsilon +O(T^2)`$ serves to cancel the divergences arising from the free correlator (6) for $`t+s0`$. The renormalized Green function reads
$$\stackrel{̊}{G}_{k\widehat{k}}^{l\widehat{l}}=Z^{k/2}(\widehat{Z}Z^{1/2})^{\widehat{k}}(ZZ_0)^{l/2}(Z_0^{1/2}\widehat{Z}Z^{1/2})^{\widehat{l}}G_{k\widehat{k}}^{l\widehat{l}}$$
(8)
and satisfies the RG equation
$`[\mu _\mu +\beta _T_T+(\widehat{\zeta }\zeta )\lambda _\lambda +\rho h_h+(\rho +\zeta _0/2)h_0_{h_0}`$
$`+(l+k\widehat{l}\widehat{k}){\displaystyle \frac{\zeta }{2}}+(\widehat{l}+\widehat{k})\widehat{\zeta }+(l+\widehat{l}){\displaystyle \frac{\zeta _0}{2}}]G_{k\widehat{k}}^{l\widehat{l}}=0,`$ (9)
with $`\beta _T=\mu _\mu T|_0`$, $`\zeta _i=\mu _\mu \mathrm{ln}Z_i|_0`$, $`\rho =\beta _T/T+\zeta /2d`$, where $`|_0`$ denotes the derivative at fixed bare parameters. Note that the equation similar to Eq. (9) holds for longitudinal and mixed Green functions.
The critical behavior of the non-linear $`\sigma `$ model is controlled by the unstable ir fixed point (FP). For $`n>2`$ and $`d>2`$ the beta function has a nontrivial zero $`T_c=\epsilon /(n2)+O(\epsilon ^2)`$, which is related to the bare critical temperature $`\stackrel{̊}{T}_c`$. The solution of Eq. (9), in conjunction with the simple dimensional analysis, yields the scaling behavior of the Green function at FP as
$`G_{k\widehat{k}}^{l\widehat{l}}(q,t;T,h,h_0)=\xi (T)^{d_G}M(T)^{l+k\widehat{l}\widehat{k}}\widehat{M}(T)^{\widehat{l}+\widehat{k}}`$
$`\times M_0(T)^{l+\widehat{l}}F_{k\widehat{k}}^{l\widehat{l}}(q\xi ,t\xi ^z;hM\xi ^d/T,h_0M_0/h)`$ (10)
where $`\xi (T)=\mu ^1T^{1/\epsilon }\mathrm{exp}(_0^T𝑑T^{}(1/\beta _T(T^{})1/\epsilon T^{}))`$ is the correlation length and $`d_G=d(k+\widehat{k}+l+\widehat{l}1)2(\widehat{l}+\widehat{k})`$ the canonical dimension of $`G_{k\widehat{k}}^{l\widehat{l}}`$. The scaling functions $`M`$, $`\widehat{M}`$ and $`M_0`$ are given by $`\mathrm{ln}M_i(T)=\frac{1}{2}_0^T𝑑T^{}\zeta _i(T^{})/\beta _T(T^{})`$. Note that in statics $`M`$ has a meaning of the spontaneous magnetization. For $`\tau =TT_c0^{}`$ we derive $`\xi |\tau |^\nu `$ and $`M|\tau |^\beta `$, with critical exponents $`\nu =1/\beta _T(T_c)`$ and $`\beta =\nu \zeta (T_c)/2`$. The dynamic exponent reads $`z=2+\widehat{\zeta }(T_c)\zeta (T_c)`$. To one loop order we have $`\nu =1/\epsilon `$, $`\beta =(n1)/2(n2)`$ and $`z=2+O(\epsilon ^2)`$.
We now focus on the scaling properties of two-times quantities which describe the critical system evolving from a completely ordered state. To that end we put in what follows $`h=h_0^1=0`$ and apply to Eq. (10) a short-time expansion of the fields $`𝝅`$ and $`\widehat{𝝅}`$ in terms of the initial fields $`𝝅_0`$ and $`\widehat{𝝅}_0`$ janssen-89 ; fedorenko-04 . As a result we obtain the scaling behavior of the response and correlation functions for $`s\xi ^z0`$ and $`t\xi ^z>0`$:
$`R_q(t,s)=A_Rs^{(z2+\eta )/z}(t/s)^{\overline{\theta }}f_R(\lambda q^z(ts),t/s),`$ (11)
$`C_q(t,s)=A_Cs^{(2\eta )/z}(t/s)^{\overline{\theta }}f_C(\lambda q^z(ts),t/s),`$ (12)
where $`\eta =2\beta /\nu \epsilon `$ is the Fisher exponent. The new dynamic critical exponent is defined as $`\overline{\theta }=(2\eta z\zeta _0(T_c)/2)/z`$, and should be distinguished from the initial slip exponent $`\theta `$ janssen-89 . To lowest order in $`\epsilon `$ it is given by $`\overline{\theta }=\epsilon (n1)/4(n2)+O(\epsilon ^2)`$ and coincides with $`\beta /\nu z`$. Although at this point we do not see any symmetry which ensures this identity, we argue that it holds to all orders in $`\epsilon `$ and below give some arguments supporting this conjecture. The functions $`f_R`$ and $`f_C`$ are regular functions of both arguments and finite for $`s0`$. We have explicitly computed the renormalized response and correlation functions to one-loop order:
$`R_q(t,s)`$ $`=`$ $`K_d^1\mathrm{\Omega }[t,s]R_q^0(t,s),`$ (13)
$`C_q(t,s)`$ $`=`$ $`K_d^1\mathrm{\Omega }[t,s]C_q^0(t,s)`$ (14)
$`+`$ $`(n2)T^2e^{\lambda q^2(t+s)}F(2\lambda q^2s)/2K_dq^2,`$
where we have introduced notations $`\mathrm{\Omega }[t,s]=1\frac{T}{2}(\gamma +\mathrm{ln}2\lambda \mu ^2s+\frac{n1}{2}\mathrm{ln}\frac{t}{s})`$ and $`F(y)=\gamma +\mathrm{ln}|y|\mathrm{Ei}(y)`$. Here $`\gamma `$ is the Euler constant and $`\mathrm{Ei}(x)`$ the exponential integral, so that $`F(0)=0`$ and $`F^{}(0)=1`$. Substituting the FP value $`T_c`$ in Eqs. (13) and (14) we find that they are consistent with scaling laws (11) and (12). The corresponding non-universal amplitudes to one loop order are given by $`A_R=K_d^1(iT_c)(1\gamma T_c/2)(2\lambda \mu ^2)^{T_c/2}`$, $`A_C=i\lambda A_R`$, and the universal (apart from the normalization) scaling functions to the same order read
$`f_R(u,v)`$ $`=`$ $`e^u,`$ (15)
$`f_C(u,v)`$ $`=`$ $`(v1)\{e^ue^{u(v+1)/(v1)}`$ (16)
$`\times `$ $`[1\epsilon F(2u/(v1))/2]\}/u.`$
Let us introduce the FDR for a particular mode $`q`$ as calabrese-02 $`X_q^1=_sC_q(t,s)/i\lambda R_q(t,s)`$. Using Eqs. (13) and (14) we obtain to one loop order that $`X_q^1=f_X(2\lambda q^2s)`$ with
$`f_X(u)`$ $`=`$ $`1+e^u+\epsilon [(n3)(1e^u)/(n2)u`$ (17)
$``$ $`e^u(F(u)2F^{}(u))]/2.`$
According to arguments of Ref. calabrese-02 the local FDR $`X^{\mathrm{}}`$ is identical to the global FDR $`X_{q=0}^{\mathrm{}}=lim_{t,s\mathrm{}}X_{q=0}(t,s)`$, which is given by
$$X_{q=0}^{\mathrm{}}=\frac{1}{2}(1+\frac{1}{4}\frac{n1}{n2}\epsilon )+O(\epsilon ^2).$$
(18)
Note that $`X_{q=0}^{\mathrm{}}>\frac{1}{2}`$ while for most systems, except for some urn models godreche-01 , the aging from a disordered state is characterized by the FDR $`X^{\mathrm{}}[0,\frac{1}{2}]`$. Thus the system aging from an ordered state is more close to equilibrium ($`X=1`$) than when it relaxes from a disordered state. We expect that this applies not only to the system under consideration but also to all systems with pure relaxation dynamics including the Ising model.
We are now in a position to discuss the relaxation of the magnetization. The asymptotic long-time behavior of the magnetization can be deduced directly from Eq. (10) for $`G_{10}^{00}`$ and reads $`M(t)t^{\beta /\nu z}`$. To describe the behavior of the magnetization during the macroscopically early initial stage of relaxation we derive the equation of motion for the magnetization. The latter can be written in the form of a Ward identity $`_0^{\mathrm{}}𝑑s\mathrm{\Gamma }_{q=0}(t,s)M(s)=0`$, which follows from the invariance of the generating functional under simultaneous rotation of fields $`𝐬`$ and $`\widehat{𝐬}`$ zinn-justin ; bausch-80 . Here $`\mathrm{\Gamma }_q(t,s)`$ is a vertex function which can be expressed as a sum of one-particle irreducible diagrams with amputated external lines $`𝝅`$ and $`\widehat{𝝅}`$. To one-loop order we have $`\mathrm{\Gamma }_q(t,s)=iT^1\delta (ts)𝒟_s(q)`$ with the renormalized operator $`𝒟_t=\widehat{Z}Z_T^1\stackrel{̊}{𝒟}_t`$ given for $`q=0`$ by
$$𝒟_t(0)=_t+\frac{T}{2}\left[\gamma +\mathrm{ln}2\lambda \mu ^2t\right]_t+\frac{n1}{4}\frac{T}{t}+O(T^2).$$
(19)
The equation of motion can be written as $`𝒟_t(0)M(t)=0`$, $`t>0`$ and $`T=T_c`$. Solving this equation to first order in $`\epsilon `$ we have to neglect the second time derivative in Eq. (19) since it carries a factor $`T`$ and hence in the equation of motion only gives rise to a correction of order $`\epsilon ^2`$. Having done that, we get $`M(t)=\overline{M}t^{\epsilon (n1)/4(n2)}`$, where the non-universal constant $`\overline{M}`$ is determined by motion of the system during the microscopically early initial stage of relaxation. Indeed, the considered field theory describes the behavior of a ferromagnet only on scales larger than the lattice constant and on times larger than the corresponding microscopic time $`\tau _{\text{mic}}`$. Therefore, the derived power-law behavior, which is usually called the non-linear critical relaxation, does exist only for $`\tau _{\text{mic}}<t<\tau _{\text{mac}}`$. Here $`\tau _{\text{mac}}`$ is the macroscopic time at which there is a crossover to linear relaxation $`M(t)\mathrm{exp}(t/\tau _{\text{mac}})`$. It is determined either by size of the system $`\tau _{\text{mac}}L^z`$ or by temperature $`\tau _{\text{mac}}|\tau |^{z\nu }`$.
Let us show how this picture can be generalized to higher orders. The vertex function is related to the corresponding Green function by the Dyson equation zinn-justin $`_0^{\mathrm{}}𝑑t^{}\mathrm{\Gamma }_q(t,t^{})R_q(t^{},s)(s)=\delta (ts)`$, which is written in the time domain and can be easily check to one loop by means of Eqs. (13) and (19). Using Eq. (11) we obtain
$`\mathrm{\Gamma }_{q=0}(t,s)=A_R^1t^{(z2+\eta )/z}\{[\delta (ts)`$
$`(t/s)^{\overline{\theta }}f^{}(t/s)/s](_s\overline{\theta }/s)+(t/s)^{\overline{\theta }}\psi (t/s)/s^2\},`$
where $`f(v)f_R(0,v)`$ and $`\psi `$ is the solution of Voltera equation $`_1^v𝑑u[f^{}(u)f^{}(v/u)v/uf(v/u)\psi (u)]=0`$ such that $`\psi (\mathrm{})=0`$. This implies that the solution of the equation of motion can be written as $`M(t)=\overline{M}t^{\overline{\theta }}f_M(t)`$ with $`f_M(t)`$ satisfying the condition $`f_M(\mathrm{})=1`$. Taking into account the known asymptotic behavior for $`t\mathrm{}`$ we conclude that the identity $`\overline{\theta }=\beta /\nu z`$ holds to all orders. To find the exact form of $`f_M`$ we have to solve the corresponding integral equation. Recently, the local scale invariance (LSI) was used to predict the exact scaling form of the response function for a system evolving from a completely disordered state henkel-01 . We expect that the transverse fluctuations of the considered model share similar scaling properties. In analogy with Ref. henkel-01 we may expect that $`f(v)=1`$ and consequently $`\psi (v)=0`$ and $`f_M(t)=1`$ to all orders. However, the analysis of two loop graphs suggests that in accordance with Ref. pleimling-05 there should be small corrections to LSI: $`f(v)=1+O(\epsilon ^2)`$. They come from the non-local in time contributions to the vertex function $`\mathrm{\Gamma }_q(t,s)`$ and give rise to very small corrections in the magnetization $`f_M(t)=1+O(\epsilon ^4)`$ that are likely difficult to observe in simulations.
We now discuss the relation between transverse and longitudinal fluctuations. Let us define the longitudinal response function as $`_q(t,s)=d^dxe^{iqx}\sigma (x,t)\widehat{\sigma }(0,s)_\text{c}`$ and analogously the longitudinal correlation function $`𝒞_q(t,s)`$. Expressing $`\sigma `$ and $`\widehat{\sigma }`$ in terms of $`𝝅`$ and $`\widehat{𝝅}`$ we obtain their low temperature expansions. To lowest order the renormalized longitudinal response function is given by
$`_q(t,s)={\displaystyle \frac{n1}{2K_d}}T[\mathrm{ln}{\displaystyle \frac{t}{ts}}+F(\lambda q^2(ts)^2/2t)`$
$`F(\lambda q^2(ts)/2)]R_q^0(t,s)+O(T^2\epsilon ,T^3).`$ (20)
The corresponding expression for the longitudinal correlator is too cumbersome so that for the sake of conciseness we write down only its time derivative at $`q=0`$
$$_s𝒞_{q=0}(t,s)=\frac{n1}{2K_d}T^2\lambda \mathrm{ln}\frac{t+s}{ts}+O(T^2\epsilon ,T^3),$$
(21)
which is needed to compute the longitudinal FDR $`𝒳_q=i\lambda _q(t,s)/_s𝒞_q(t,s)`$. For the latter we obtain $`𝒳_{q=0}^{\mathrm{}}=\frac{1}{2}+O(\epsilon )`$. The computation to order $`\epsilon `$ requires two-loop calculations and is left for future investigations. However, we expect that the scaling ansatzs (11) and (12) are valid also for the longitudinal functions.
After we submitted our preprint to arxiv.org there appeared Refs. annibale-05 and garriga-05 which also study the critical aging from a magnetized state but for infinite range models such as the spherical model, long range ferromagnetic model and Ising model in the limit of large dimension. These works confirm that the critical aging from an ordered state is characterized by $`X^{\mathrm{}}>1/2`$. The $`d2`$ expansions of the response and correlation functions obtained in Ref. annibale-05 for the spherical model are in agreement (up to prefactors) with our results for the longitudinal functions (20) and (21). The difference between the spherical limit ($`n\mathrm{}`$) of the transverse FDR (18) and the FDR of the spherical model annibale-05 suggests that $`𝒳^{\mathrm{}}X^{\mathrm{}}`$ for finite $`n`$ and thus one cannot introduce the observable-independent effective temperature $`T_{\mathrm{eff}}=T/X^{\mathrm{}}`$ cugliandolo-97 . Rewriting the response function of the spherical model in the form (11) we obtain $`f_R^{\mathrm{sp}}(0,v)=1(11/v)^{(d2)/2}`$. Thus in contrast to the transverse response of the $`O(n)`$ system with the scaling function $`f_R(0,v)`$ being finite in the limit $`v\mathrm{}`$, the scaling function of the spherical model decays as $`f_R^{\mathrm{sp}}(0,v)v^1`$. Supposing the same asymptotics for the longitudinal function $`f_{}(0,v)`$ of the $`O(n)`$ system we arrive at $`_{q=0}(t,s)s^{(z2+\eta )/z}(t/s)^{\overline{\theta }1}`$ for $`ts`$. This, of course, has to be checked independently.
In summary, we have studied the nonequilibrium critical behavior of a ferromagnetic systems which evolves from a completely ordered state. We have shown that the universal long-time scaling behavior of magnetization emerges in the macroscopically early initial stage of relaxation and that can be considered as a consequence of the LSI. The two-times functions exhibit aging behavior which in contrast to the aging from a disordered state is characterized by the FDR larger than $`1/2`$. Numerical simulations of Ref. zheng-96 show that the generalized scaling behavior exists in the short-time regime of the critical relaxation even for an initial state with arbitrary magnetization. However, instead of the critical exponent $`\theta `$ there exists a whole universal characteristic function which up to now can be estimated only numerically. We expect that our computation being extended to finite $`h_0`$ can describe this scaling behavior and provide a way to calculate the corresponding characteristic function. We hope that the considered model can also be useful to study the influence of Goldstone modes on phase ordering kinetics at $`T<T_c`$.
We would like to thank G. Schehr, A. Gambassi, M. Henkel, B. Zheng, K.J. Wiese for useful discussions and especially S. Stepanow and M. Pleimling for a critical reading of the manuscript. The support from the DFG (SFB 418) is gratefully acknowledged. |
warning/0507/quant-ph0507065.html | ar5iv | text | # Theory of Photon Blockade by an Optical Cavity with One Trapped Atom
## I Introduction
The phenomenon of photon blockade, first proposed in Ref. imamoglu97 in analogy with Coulomb blockade for electrons fulton87 ; kastner92 ; likharev99 , occurs when the absorption of a first input photon by an optical device blocks the transmission of a second one, thereby leading to nonclassical output photon statistics. Photon blockade has been predicted in many different settings grangier98 ; werner99 ; rebic99 ; rebic02 ; kim99 ; smolyaninov02 , including for a single two-level atom in cavity QED rebic02 ; tian92 ; brecha99 ; hood-thesis . In the latter setting, the blockade is due to the anharmonicity of the Jaynes-Cummings ladder of eigenstates jaynes63 . If an incoming photon resonantly excites the atom-cavity system from its ground state to $`|1,\pm `$ (where $`|n,+()`$ denotes the $`n`$-excitation dressed state with higher (lower) energy), then a second photon at the same frequency will be detuned from either of the next steps up the ladder, i.e. from states $`|2,\pm `$. In the strong coupling regime kimble98 , for which the coherent rate of evolution $`g_0`$ exceeds the dissipative rates $`\kappa `$ and $`\gamma `$, this detuning will be much larger than the excited-state line widths, so that the two-excitation manifold will rarely be populated. This in turn leads to the ordered flow of photons in the transmitted field, which emerge from the cavity one at a time.
We recently reported observations of photon blockade in the light transmitted by an optical cavity containing one atom strongly coupled to the cavity field birnbaum05 . For coherent excitation at the cavity input, the photon statistics for the cavity output displayed both photon antibunching and sub-Poissonian photon statistics. However, as illustrated in Fig. 1, the multiplicity of atomic and cavity states makes our experiment considerably more complex than the simple situation described by the extended Jaynes-Cummings model with damping tian92 ; rebic02 ; brecha99 ; hood-thesis . In our paper birnbaum05 and the accompanying Supplemental Information birnbaum05s , we presented theoretical results from an extended model that described a multistate atom coupled to two cavity modes. The relevant atomic states are the Zeeman states of a particular hyperfine transition in atomic Cesium, namely $`6S_{1/2},F=46P_{3/2},F^{}=5^{}`$ at $`852`$ nm. The relevant modes of the Fabry-Perot cavity are two TEM<sub>00</sub> modes of the same longitudinal order but orthogonal polarizations.
Our purpose in this paper is to provide a more complete discussion of photon blockade than in Ref. birnbaum05 , both by way of particular results from our model calculation and of a general framework for characterization of photon blockade, with which we begin in Section II. In Section III we turn to the details of our actual system and calculate the eigenvalue structure for a two-mode cavity coupled to an atom with multiple internal states. Here, we make explicit the coupling in the model Hamiltonian which is used to determine the eigenvalues displayed in Figure 1(b) of Ref. birnbaum05 . We also incorporate this Hamiltonian into the master equation for the damped, driven system used to compute theoretical results for transmission spectra and photon statistics, as in Figure 2(b) of Ref. birnbaum05 . In Section IV, we gain some perspective on the relevant physical mechanisms by comparing transmission spectra and photon statistics for the case of an external drive that excites the atom (rather than the cavity, as in our experiment). Section V presents an extension to our atom-cavity model which includes the effect of cavity birefringence and FORT-induced ac-Stark shifts in the atomic states. The modified cavity transmission and intensity correlation functions are presented for comparison to previous results. We also display a theoretical result for the time dependence of the intensity correlation function for comparison to our measured result in Ref. birnbaum05 . Finally, in Section VI we offer a discussion of these various results.
## II General considerations
We begin by attempting to address the question “What is photon blockade?”werner99 . Consider the following input electromagnetic field state to some “black box”,
$$|\psi _{in}=a_n|n\text{ .}$$
(1)
Assume that the mapping of input to output by the black box is given by the transmission coefficients $`t_n`$ for each Fock-state $`|n`$. The output state is then of the form
$$|\psi _{out}=t_na_n|n\text{ .}$$
(2)
For a coherent-state input $`|\alpha `$, $`a_n\alpha ^n/\sqrt{n!}`$, so that
$$|\psi _{out}t_n\alpha ^n/\sqrt{n!}|n\text{ .}$$
(3)
A linear transfer function for the black box would be of the form $`t_n(t_1)^n`$, so that
$`|\psi _{out}`$ $``$ $`{\displaystyle (t_1\alpha )^n/\sqrt{n!}|n}`$
$``$ $`|t_1\alpha \text{ .}`$
By contrast, a nonlinear device can modify the photon statistics in a fashion other than $`|t_n||t_1|^n`$. The working criterion that we adopt here for photon blockade is that the transmission coefficients $`|t_n|<|t_1|^n`$ for $`n2`$. For example, an “ideal” photon blockade device that eliminates all Fock states with $`n2`$ ($`t_n=0`$ for $`n2`$) would lead to the following output for a coherent state input:
$$|\psi _{out}|0+t_1\alpha |1\text{ .}$$
(5)
Likewise, a device that produces photon pairs in abundance (and an associated large degree of photon bunching) could be specified by $`|t_{n=2}||t_1|^2`$, with all $`t_{n>2}=0`$.
Of course, even in this simple setting of $`inout`$, the above discussion is incomplete since at least one additional input and output channel is required to preserve unitarity. More generally, the transformation $`inout`$ requires multiple input and output channels (e.g., polarizations for the input and output fields with a continuum of frequencies, relevant quantum degrees of freedom for the material system of the black box, etc.). These additional channels may affect the coherence of the output field, as discussed below in Section VI. Within this more complex setting, however, the conceptual framework that we suggest for identifying photon blockade still rests upon the simple intuition described above. Namely, one of the output channels should have the property that the transmission coefficients $`t_n`$ satisfy $`t_n<|t_1|^n`$ for $`n2`$.
Note that the pioneering work on photon blockade based upon EIT satisfies the above criterion imamoglu97 ; grangier98 ; werner99 ; rebic99 ; rebic02 , as does resonance fluorescence from a single atom kimble77 and the cavity QED schemes considered in Refs. tian92 ; rebic02 ; brecha99 ; hood-thesis . Indeed, by the criterion stated above, most single atomic or molecular emitters function by way of photon blockade. The problem of course is that the efficiency for collecting fluorescence is typically poor. So, in addition to the more fundamental requirement $`t_n<|t_1|^n`$, it seems reasonable to add a second, more practical criterion related to efficiency. Photon blockade is of much less practical significance if the efficiency for the mapping of input to an output channel is negligibly small, but precisely “how small is too small” is hard to quantify and depends upon the particular application. Much of the effort related to photon blockade is directed towards maintaining the “quality” of blockade inherent in single-atom resonance fluorescence while at the same time achieving a sensibly large efficiency, which has led our group to employ an optical cavity within the setting of cavity QED, a system to which we now turn our attention.
## III Eigenvalues of the atom-cavity system
We begin this section, adapted from Ref. birnbaum05s , by considering the eigenvalue structure of the atom-cavity system in the absence of damping, with the model system illustrated in Figure 1. Approximating the atom-cavity coupling as a dipole interaction, we define the atomic dipole transition operators for the $`6S_{1/2},F=46P_{3/2},F^{}=5^{}`$ transition in atomic Cesium as
$$D_q=\underset{m_F=4}{\overset{4}{}}|F=4,m_FF=4,m_F|\mu _q|F^{}=5^{},m_F+qF^{}=5^{},m_F+q|,$$
(6)
where $`q=\{1,0,1\}`$ and $`\mu _q`$ is the dipole operator for $`\{\sigma _{},\pi ,\sigma _+\}`$-polarization, respectively, normalized such that for the cycling transition $`F=4,m_F=4|\mu _1|F^{}=5^{},m_F=5=1`$. The matrix element of the dipole operator $`F=4,m_F|\mu _q|F^{}=5^{},m_F^{}`$ is equal to the Clebsch-Gordan coefficient for adding spin $`1`$ to spin $`4`$ to reach total spin $`5`$, namely $`j_1=4,j_2=1;m_1=m_F,m_2=q|j_{total}=5;m_{total}=m_F^{}`$.
The Hamiltonian of a single atom coupled to a cavity with two degenerate orthogonal linear modes is
$`H_{45^{}}`$ $`=`$ $`\mathrm{}\omega _A{\displaystyle \underset{m_F^{}=5}{\overset{5}{}}}|F^{}=5^{},m_F^{}F^{}=5^{},m_F^{}|+\mathrm{}\omega _{C_1}(a^{}a+b^{}b)`$
$`+\mathrm{}g_0(a^{}D_0+D_0^{}a+b^{}D_y+D_y^{}b)`$
where $`\omega _A`$ is the atomic transition frequency, $`\omega _{C_1}`$ is the cavity resonance frequency, and $`D_y=\frac{i}{\sqrt{2}}(D_1+D_{+1})`$ is the dipole operator for linear polarization along the $`y`$-axis. We are using coordinates where the cavity supports $`\widehat{y}`$ and $`\widehat{z}`$ polarizations and $`\widehat{x}`$ is along the cavity axis. The annihilation operator for the $`\widehat{z}`$ ($`\widehat{y}`$) polarized cavity mode is $`a`$ ($`b`$).
Assuming $`\omega _A=\omega _{C_1}\omega _0`$, we find that the lowest eigenvalues of $`H_{45^{}}`$ have a relatively simple structure. In the manifold of zero excitations, all nine eigenvalues are zero. In manifolds with $`n`$ excitations, the eigenvalues are of the form $`E_{n,k}=n\mathrm{}\omega _0+\mathrm{}g_0\epsilon _k^{(n)}`$, where $`\epsilon _k^{(n)}`$ is a numerical factor and $`k`$ is an index for distinct eigenvalues. There are $`29`$ states in the $`n=1`$ manifold, but due to degeneracy $`k`$ has only $`13`$ distinct values, $`k\{6,\mathrm{}6\}`$; in the $`n=2`$ manifold there are $`49`$ states but $`k\{11,\mathrm{}11\}`$. The total number of states in any manifold can be understood by considering how the excitations can be distributed among the atom and the two cavity modes. For example, in the $`n=1`$ manifold, the atom can be in one of its $`9`$ ground states ($`m_F\{4,\mathrm{}4\}`$) and either cavity mode $`l_y`$ or mode $`l_z`$ can have one photon (giving $`18`$ possible states), or the atom can be in one of its $`11`$ excited states ($`m_F^{}\{5,\mathrm{}5\}`$) while both cavity modes are in the vacuum state, yielding a total of 29 states. Table 1 lists numerical values for $`\epsilon _k^{(1,2)}`$ as well as their respective degeneracies $`\eta _k^{(1,2)}`$. The numerical factors and degeneracies have the symmetries $`\epsilon _k^{(n)}=\epsilon _k^{(n)}`$ and $`\eta _k^{(n)}=\eta _k^{(n)}`$. The resulting eigenvalues $`E_{n,k}`$ for $`n=\{0,1,2\}`$ are displayed in Fig. 2(b).
Although these eigenvalues are certainly not sufficient for understanding the complex dynamics associated with the full master equation, they do provide some insight into some structural aspects of the atom-cavity system. For example, the eigenvalues $`\epsilon _{\pm 6}^{(1)}=\pm 1`$ correspond to the vacuum-Rabi splitting for the states $`|1,\pm `$ for a two-state atom coupled to a single cavity mode \[cf., Fig. 1(a) of Ref. birnbaum05 \]. The one-photon detunings for transitions from the $`n=1n^{}=2`$ manifold are largest for the eigenstates associated with $`\epsilon _{\pm 6}^{(1)}`$. Indeed, just as for the two-state atom with one cavity mode, transitions from the eigenstates at $`\pm g_0`$ have frequency detunings $`\pm (2\sqrt{2})g_0`$ relative to the nearest states in the $`n^{}=2`$ manifold (at $`\epsilon _{\pm 11}^{(2)}=\pm \sqrt{2}`$, respectively). Hence, as a function of probe frequency $`\omega _p`$, the eigenvalue structure in Table 1 suggests that the ratio of two-photon to one-photon excitation would exhibit a minimum around $`\omega _p=\omega _0\pm g_0`$, resulting in reduced values $`g^{(2)}(0)<1`$ g2-define , which the full calculation verifies in Fig. 2(b) of Ref. birnbaum05 .
For excitation to the other eigenstates in the $`n=1`$ manifold, such blockade is not evidenced in Fig. 2(b) of Ref. birnbaum05 . A contributing factor suggested by the structure of eigenvalues in Table 1 is interference of one and two-photon excitation processes. For example, excitation at $`\omega _p\omega _0\pm g_0/4`$ results in two-photon resonance for the eigenstates associated with $`\epsilon _{\pm 1}^{(2)}\pm 0.5`$, and leads to photon bunching with $`g^{(2)}(0)1`$ as confirmed by our full calculation of photon statistics.
Fig. 2(a) provides a global perspective of these various effects. Here, we calculate transmission spectra and intensity correlation functions analogous to those shown in Figure 2(b) of Ref. birnbaum05 , but now with coherent coupling $`g_0`$ much larger than the dissipative rates $`(\kappa ,\gamma )`$ and well beyond what we have achieved in our experiments, $`g_0/\kappa =g_0/\gamma =50`$ T-define . At $`\omega _p=\omega _0\pm g_0`$, $`g_{yz}^{(2)}(0)0.002`$ in evidence of the previously discussed photon blockade suggested by the eigenvalue structure in Table 1. As anticipated, large photon bunching results near $`\omega _p\omega _0\pm g_0/4`$ associated with the two-photon resonance to reach the eigenstates with $`\epsilon _{\pm 1}^{(2)}\pm 0.5`$. Between these two extremes for the eigenvalues with the largest and smallest nonzero magnitudes ($`g_0/4|\omega _p\omega _0|g_0`$), $`g_{yz}^{(2)}(0)`$ displays a complex structure involving multiple excitation pathways through states in the $`n=1`$ manifold to reach states in the $`n^{}=2`$ manifold. The extremely large peak at $`\omega _p=\omega _0`$ is discussed in Refs. carmichael91 ; brecha99 . Similar calculations show that the photon blockade effect described by Fig. 2 is unaffected in its qualitative character if the atomic spontaneous decay rate $`\gamma `$ is made much smaller than the cavity decay rate $`\kappa `$, although we have not set $`\gamma `$ strictly to zero.
## IV Driven Atom
Fig. 2(b) of Ref. birnbaum05 compares the predicted photon statistics when driving the detected cavity mode ($`\widehat{z}`$) with the statistics when driving the other cavity mode ($`\widehat{y}`$). The driven cavity mode has photon statistics which are less strongly sub-Poissonian, an effect we hypothesize to be caused by interference between the atomic dipole radiation and the coherent drive. We will now further explore this hypothesis by considering atom-cavity systems where the driving field is directly coupled to the atom, instead of the cavity mode.
We first study the familiar Jaynes-Cummings system, a two-state atom coupled to a single mode cavity. We assume a coherent drive field made of many photons which we will treat classically. In Fig. 3, we compare the intracavity fields carmichael85 when driving (a) the cavity and (b) the atom for $`g_0/\kappa =g_0/\gamma =50`$. $`T,\stackrel{~}{T}`$ are proportional to the intracavity photon number, with $`T`$ normalized to the empty-cavity on-resonance photon number for the driven cavity and $`\stackrel{~}{T}`$ normalized to half of the saturation parameter $`s`$ for the driven atom. The intensity correlation function $`\stackrel{~}{g}^{(2)}(0)`$ of the system when driving the atom is much lower than $`g^{(2)}(0)`$ (the correlation function of the system when driving the cavity) at $`\omega _p=\omega _0\pm g_0`$, with $`\stackrel{~}{g}^{(2)}(0)0.002`$ and $`g^{(2)}(0)0.02`$. $`\stackrel{~}{g}^{(2)}(0)`$ is super-Poissonian at $`\omega _p=\omega _0\pm g_0/\sqrt{2}`$ due to the two-photon resonance discussed above, but lacks the large peak at $`\omega _p=\omega _0`$ evident in $`g^{(2)}(0)`$. At $`\omega _p=\omega _0`$ we have $`\stackrel{~}{g}^{(2)}(0)1`$.
We next consider a two-mode cavity coupled to the Zeeman states of the $`F=4F^{}=5^{}`$ transition of a single atom, as in Section III. We take the probe field driving the atom to be polarized along $`\widehat{z}`$ and calculate the intracavity photon number $`\stackrel{~}{T}_{zz},\stackrel{~}{T}_{yz}`$ in the $`\widehat{z},\widehat{y}`$ modes, respectively (normalized to the drive strength as above), and the corresponding intensity correlation functions $`\stackrel{~}{g}_{zz}^{(2)}(0),\stackrel{~}{g}_{yz}^{(2)}(0)`$. Results of the calculations are plotted in Fig. 4 for $`g_0/\kappa =g_0/\gamma =50`$. At $`\omega _p=\omega _0\pm g_0`$, both $`\stackrel{~}{g}_{zz}^{(2)}(0)`$ and $`\stackrel{~}{g}_{yz}^{(2)}(0)`$ are close to $`g_{yz}^{(2)}(0)`$ as in Fig. 2. This supports our hypothesis that the somewhat higher value of $`g_{zz}^{(2)}(0)`$ at $`\omega _p=\omega _0\pm g_0`$ is caused by interference with the drive field. Interestingly, though the central peak at $`\omega _p=\omega _0`$ is absent in $`\stackrel{~}{g}_{zz}^{(2)}(0)`$, as we expect in correspondence with the Jaynes-Cummings case and as implied by the usual explanation of the phenomenon as resulting from the interference of the drive field with the atomic dipole radiation, the peak in $`\stackrel{~}{g}_{yz}^{(2)}(0)`$ is even greater than that of $`g_{zz}^{(2)}(0),g_{yz}^{(2)}(0)`$. This effect may be of interest in future studies.
## V Birefringence and Stark Shifts
We now consider the effects of cavity birefringence and $`m_F^{}`$-dependent ac-Stark shifts, expanding our previous treatment from Ref. birnbaum05s . The birefringence and ac-Stark shifts modify the Hamiltonian $`H_{45^{}}`$ in Eq. (III) to
$`H_{full}`$ $`=`$ $`{\displaystyle \underset{m_F^{}=5}{\overset{5}{}}}\mathrm{}\omega _{m_F^{}}|F^{}=5^{},m_F^{}F^{}=5^{},m_F^{}|+\mathrm{}\omega _{C_1^z}a^{}a+\mathrm{}\omega _{C_1^y}b^{}b`$
$`+\mathrm{}g_0(a^{}D_0+D_0^{}a+b^{}D_y+D_y^{}b)`$
The birefringent splitting $`\mathrm{\Delta }\omega _{C_1}`$ is the difference of the resonant frequencies of the two polarization modes, $`\mathrm{\Delta }\omega _{C_1}=\omega _{C_1^z}\omega _{C_1^y}`$. The atomic excited state frequencies are given by $`\omega _{m_F^{}}=\omega _A+U_0\beta _{m_F^{}}`$, where $`\omega _A`$ is the unshifted frequency of the $`F=4F^{}=5^{}`$ transition in free space, $`U_0`$ is the FORT potential, and $`\beta _{m_F^{}}`$ for the FORT wavelength of the experiment is given by $`\{m_F^{},\beta _{m_F^{}}\}`$ $`=`$ $`\{\pm 5,0.18\},`$ $`\{\pm 4,0.06\},`$ $`\{\pm 3,0.03\},`$ $`\{\pm 2,0.10\},`$ $`\{\pm 1,0.14\},`$ $`\{0,0.15\}`$ mckeever03 .
The effect of these corrections to the Hamiltonian on the transmitted field from the steady-state solutions to the master equation are displayed in Fig. 5, with the parameters corresponding to the experimental values from Ref. birnbaum05 . The heights and shapes of the multiplets in $`T_{yz,zz}`$ are modified, but the basic structure is unaffected relative to Fig. 2(b) of Ref. birnbaum05 . The structure of $`g_{yz,zz}^{(2)}(0)`$ is also qualitatively unchanged. The asymmetry of the plots about $`\omega _0`$ is caused by the effective atom-cavity detuning (mostly due to the ac-Stark shifts). The value of $`g_{yz}^{(2)}(0)`$ for $`\omega _p=\omega _0g_0`$ is $`0.02`$ (ignoring the above corrections yields $`g_{yz}^{(2)}(0)0.03`$). These values are consistent with the experimental result of Ref. birnbaum05 , $`g_{yz}^{(2)}(0)=0.13\pm 0.11`$. Note that we have previously reported measurements of $`T_{zz}`$ and made detailed comparisons with the theory described here boca04 .
In Fig. 6, we present the theoretical prediction for $`g_{yz}^{(2)}(\tau )`$ including the effects of Stark shifts and birefringence. We find that $`g_{yz}^{(2)}(\tau )`$ rises to unity at $`\tau 80`$ ns, well above the experimental result of $`\tau 45`$ ns. The theory, however, does not take into account atomic motion, but rather assumes the atom to be fixed at a place of optimal coupling. Since atomic motion clearly has a large effect on $`g_{yz}^{(2)}(\tau )`$, as is evident in Fig. 4(b) of Ref. birnbaum05 , we believe that this is the dominant cause of the discrepancy diedrich87 .
## VI Discussion
From the point of view described in general terms in Section II and elaborated in more detail in Sections III-V, our work in Ref. birnbaum05 satisfies the criteria for photon blockade. As indicated by Fig. 1(a,b) of Ref. birnbaum05 , the Jaynes-Cummings ladder provides a means to achieve the condition $`t_2<|t_1|^2`$, with the consequence that the intensity correlation function $`g^{(2)}(0)<1`$ as presented in Fig. 2 of Ref. birnbaum05 and in the previous sections. Note also that by tuning to a two-photon resonance, our calculations indicate that photon bunching $`g^{(2)}(0)1`$ could be achieved for our atom-cavity system with ($`\omega _p\omega _0\pm g_0/4`$), again in accord with an understanding based upon transmission coefficients $`t_n`$. Photon bunching for the standard Jaynes-Cummings ladder of Fig. 1(a) in Ref. birnbaum05 is likewise achieved by tuning to the two-photon resonance at $`\omega _p\omega _0\pm g_0/\sqrt{2}`$.
Our criteria for photon blockade do not demand the preservation of coherence in the transformation from input to output. One way to express a requirement for coherence in terms of the generic model in Section II is that the output state should be of the form of Eq. (5), and not of the form
$$\rho _{out}|00|+|t_1\alpha |^2|11|\text{ ,}$$
(9)
with the coherent amplitude lost. However, our view is that either Eq (5) or Eq. (9) suffices and qualifies as photon blockade. In fact, in practice the latter case of Eq. (9) might be the more “useful” for the following reasons.
With reference to Fig. 8 in Ref. rebic02 , note that there are various contributions to $`g^{(2)}(0)`$ as a function of the amplitude of the driving field \[Eq. (11b) of Ref. rebic02 \]. The authors point out that “The decomposition shows how the behavior of $`g^{(2)}(0)`$ … can be interpreted as the effect of self-homodyning between the coherent and incoherent components of the intracavity field.” For our initial experiments with relatively modest ratios $`g/(\kappa ,\gamma )`$, our view is that this complex, phase-dependent interplay should be avoided since it makes the blockade effect more “fragile” (less robust) than is the case for Eq. (9). Basically, one wants a situation where there is no need to balance a set of interference terms (as in Eq. (11b) of Ref. rebic02 ), but rather a more “generic” requirement of the sort presented in Section II, namely $`|t_n|<|t_1|^n`$. Note that an interpretation of photon antibunching similar to that expressed in Ref. rebic02 can be given for single-atom resonance fluorescence, as was first analyzed by Carmichael carmichael85 . However, in this case the “miracle” is that the terms always sum to give $`g^{(2)}(0)=0`$ for any drive strength, which is not the case for an atom in a cavity carmichael85 .
Our calculations show that the mean value for the amplitude of the transmitted field with polarization orthogonal to that of the coherent state input is zero. So, there is not a sense in which the coherent amplitude of the input is rotated in polarization by the atom-cavity system in the transformation to the output. Nevertheless, the orthogonally polarized output field can exhibit quantum interference, with an example being the value $`g_{yz}^{(2)}(0)1`$ for $`\omega _p=\omega _0`$ in Figs. 2, 5, which is presumably associated with the quantum state reduction and interference described by Carmichael and coworkers carmichael91 ; brecha99 , with here $`g_{yz}^{(2)}(0)g_{zz}^{(2)}(0)`$ for $`\omega _p=\omega _0`$.
For our experiment with multiple Zeeman states in the ground and excited levels and with two orthogonally polarized cavity modes $`(y,z)`$, the eigenstates that are being driven for excitation along $`y`$ with $`\omega _p=\omega _0\pm g_0`$ are complex superpositions of various atomic Zeeman states and field states for the $`(y,z)`$ polarizations. These eigenstates are entangled, so it is perhaps not surprising that examination of a particular component (e.g., for the field or atomic coherences) results in a mixed state with mean zero. Hence, many coherences for a single degree of freedom may vanish not because of dissipation per se, but rather because of entanglement with other components. Although we have not explored this question in detail, we think that this may explain why the mean amplitude for the transmitted $`z`$ field is zero for excitation along $`y`$, as is the case for our results in Figs. 2, 5.
A final perspective to offer is to extend the discussion from the case with continuous excitation as has been implicit above to the case of pulsed excitation, as in Fig. 4 of Ref. imamoglu97 . Assume a pulse with duration short compared to any time scale associated with the “black box” described in Section II. For resonance fluorescence, consider a $`\pi `$ pulse with duration $`\tau _p\gamma ^1`$, so that there is now no coherent component for the fluorescent light for times $`t>`$ $`\tau _p`$. In this case, resonance fluorescence would no longer satisfy the criterion for the preservation of coherence (as is the case for weak cw excitation mandel-wolf95 ), yet it would still be perfectly antibunched. For the case of our atom-cavity system, we would require $`g_0^1\tau _p\kappa ^1`$, in which case we would presumably obtain single photons on a pulse-by-pulse basis for the transmitted field with polarization orthogonal to that of the drive field, again with no preservation of coherence, which is presumably also true for EIT schemes in the limit of short pulses for the excitation imamoglu97 ; grangier98 ; werner99 ; rebic99 ; rebic02 . Note that in either of these cases, the efficiency for the transformation of the pulsed driving field to a single photon at the output is necessarily small because of the mismatch of bandwidths, $`\tau _p\kappa ^1`$.
## VII Acknowledgements
We gratefully acknowlege stimulating discussions with Atac Imamoḡlu upon which the material in Sections II and VI is based. This research is supported by the National Science Foundation, by the Caltech MURI Center for Quantum Networks, and by the Advanced Research and Development Activity (ARDA). |
warning/0507/hep-th0507258.html | ar5iv | text | # 1 Introduction
## 1 Introduction
There is presently a great deal of interest in the possibility that Lorentz and CPT invariance may be violated in nature. If the fundamental laws of physics do not obey these symmetries, then we would expect to see evidence of this violation even in the low-energy effective theory. Therefore, if small Lorentz or CPT violations were discovered, they would represent crucial clues about the structure of the most basic theory of nature. A general standard model extension (SME), containing possible Lorentz- and CPT-violating corrections to quantum field theory and general relativity has been developed. The SME offers a parameterization of Lorentz violations in low-energy effective field theory, and both its renormalizability and stability have been carefully examined.
The SME provides a theoretical framework for analyzing experimental results. Sensitive tests of Lorentz symmetry have included studies of matter-antimatter asymmetries for trapped charged particles and bound state systems , determinations of muon properties , analyses of the behavior of spin-polarized matter , frequency standard comparisons , measurements of neutral meson oscillations , polarization measurements on the light from distant galaxies , and others. The results of these experiments can be used to set bounds on various SME coefficients. Many coefficients are very strongly constrained, but many others are not.
There are many systems and reaction processes that could potentially be used to set further bounds of the SME coefficients for Lorentz violation. We shall consider a particular process—synchrotron motion and radiation—and examine how it would be impacted by Lorentz violation. Although there have been many analyses of this process in the presence of Lorentz-violating dispersion relations, there is as yet no analysis in terms of the renormalizable operators of the SME.
Analyses of possible Lorentz violation in synchrotron emission have often focused only on changes to particle dispersion relations. One popular approach is that of Myers and Pospelov . Taking a preferred direction $`v^\mu `$ in spacetime, one may add an operator proportional to $`i\varphi ^{}\left(v^\mu _\mu \right)^3\varphi `$ to the Lagrange density for a scalar particle. If $`v^\mu `$ has a time component only, this will add a term proportional to $`E^3`$ to the usual relativistic energy-momentum relation $`E^2=\stackrel{}{p}^2+m^2`$. Of course, the statement that $`v^\mu `$ is purely timelike is not Lorentz invariant, so that assumption must be taken to hold is some particular preferred frame, which is typically the rest frame of the cosmic microwave background. The electromagnetic field is incorporated through the usual minimal coupling procedure. In the presence of this kind of Lorentz violation, the motion of a charged particle in a constant magnetic field is modified, but the projection of the trajectory onto the plane perpendicular to $`\stackrel{}{B}`$ remains circular, and the particle’s speed remains constant. The radiation in the far field can be determined, including information about polarization, and circumstances that could enhance observable effects have been identified .
Stringent bounds on Lorentz violations with modified dispersion relations have been obtained from data from the Crab nebula . These modifications can lead to maximum particle velocities that are less than the speed of light, but the Crab nebula shows evidence of synchrotron emission from electrons with Lorentz factors of $`\gamma =\left(1\stackrel{}{v}^2\right)^{1/2}3\times 10^9`$, or energies of 1500 TeV. So the existence of electrons with velocities this large can be used to constrain models with deformed dispersion relations. For a Lorentz-violating coefficient with a particular sign, the data show that the coefficient must be at least seven orders of magnitude smaller than $`𝒪(E/M_P)`$ Planck-level suppression.
Lorentz violation can also be incorporated into particle physics through the introduction of noncommutative field theory . Synchrotron radiation has also been analyzed within this framework. The minimal coupling between charged matter and the electromagnetic field is modified by the noncommutativity, as is the structure of the free radiation field itself. A discussion in focuses on the particular case in which the magnetic field and the Lorentz-violating noncommutativity parameter are aligned, so that the orbits of charged particles in the plane perpendicular to $`\stackrel{}{B}`$ are again given by circles. It is possible to work out the far fields within this model, at leading order in the noncommutativity, but there are a number a difficulties, including acausality and potential problems with quantization.
However, these analyses ignore some of the most natural Lorentz-violating operators. There is a unique spin-independent, superficially renormalizable SME coupling that is consistent with the gauge invariance of the standard model and which grows in relative importance at high energies. This is a CPT-even two-index tensor $`c^{\nu \mu }`$. We shall look at how the presence of such a constant background tensor (which could arise, for example, as the vacuum expectation value of a dynamical tensor field) will modify synchrotron emission. Using existing data about the nonthermal spectrum of the Crab nebula, we may place a bound of $`6\times 10^{20}`$ on a particular linear combination of the $`c^{\nu \mu }`$ coefficients. The method by which we find this bound is very similar to that used to bound other types of Lorentz violation; however, the $`c^{\nu \mu }`$ interaction is actually much more natural to consider than these, because it is superficially renormalizable. All the analyses so far have been essentially classical in nature, and we shall continue working within the classical framework, although we shall look at when quantum corrections would become important.
## 2 Synchrotron Motion with Lorentz Violation
To study synchrotron motion, we shall consider a theory of fermions interacting with the electromagnetic field. The Lagrange density for this theory is
$``$ $`=`$ $`{\displaystyle \frac{1}{4}}F^{\mu \nu }F_{\mu \nu }+\overline{\psi }[\mathrm{\Gamma }^\mu (i_\mu eA_\mu )m]\psi `$ (1)
$`=`$ $`{\displaystyle \frac{1}{4}}F^{\mu \nu }F_{\mu \nu }+\overline{\psi }[(\gamma ^\mu +c^{\nu \mu }\gamma _\nu )(i_\mu eA_\mu )m]\psi .`$
The $`c^{\nu \mu }`$ interaction is the source of the Lorentz violation. There are other superficially renormalizable couplings contained in the Standard Model Extension, but the $`c`$ couplings are most natural in this context. When considering synchrotron radiation, one is primarily interested in particles with very high energies. Lorentz-violating coefficients that modify the kinetic part of the Lagrangian will grow in relative importance at high energies, so it is natural to consider only these kinetic modifications. There are only two such sets of Lorentz-violating terms that are consistent with the more general standard model gauge couplings—the $`c`$ terms and also a set of $`d^{\nu \mu }`$ terms, which have the same form as the $`c`$ interactions, except for the addition of a $`\gamma _5`$.
However, we shall not consider the $`d`$ interactions here. They are spin-dependent, while the $`c`$ term exists for bosonic (Klein-Gordon) particles as well as fermions. So all our results will apply equally to the motion of spin-zero charged particles. Moreover, spin precession effects will naturally decrease the importance of any $`d`$ terms. For an electron undergoing circular cyclotron motion, with the spin oriented in the plane of the orbit, the spin rotates by $`2\pi \gamma \frac{g2}{2}`$ radians with each orbital revolution. For $`\gamma \alpha ^1`$, the spin will rotate many times during one orbital period, and any effects proportional to the helicity will be diminished by the resultant averaging.
Modifications of the kinetic Lagrangian that are not invariant under the standard model’s $`SU(2)_L`$ gauge symmetry can also exist; however, they can only appear as part of electroweak symmetry breaking, as vacuum expectation values of nonrenormalizable operators. These operators should therefore be further suppressed, and we shall neglect them.
We shall also neglect any Lorentz violation in the photon sector. Modifications of the free electromagnetic Lagrangian will generally change the speed of photon propagation. This leads to the possibility of vacuum Cerenkov radiation , which is not yet fully understood, although threshold analyses can be used to set further limits on Lorentz-violating parameters. Most possible Lorentz-violating terms in the free electromagnetic sector also give rise to photon birefringence, which has been searched for and not seen. The limits on the relevant forms of Lorentz violation are very strong, and we may safely neglect them. The purely electromagnetic terms that do not cause birefringence can be accounted for by adding
$$_F=\frac{1}{4}\left(k_F\right)_{\mu \alpha \nu }^\alpha \left(F^{\rho \mu }F_\rho ^\nu +F^{\mu \rho }F_\rho ^\nu \right).$$
(2)
to $``$. However, a coordinate transformation $`x^\mu x^\mu \frac{1}{2}\left(k_F\right)_{\alpha \nu }^{\alpha \mu }x^\nu `$ will eliminate all the Lorentz violation from the photon sector at leading order . This transformation shifts the Lorentz-violating physics into the charged matter sector, where it manifests itself exactly as a $`c^{\nu \mu }`$ term. We see that consideration of $`c`$ therefore captures all the possible sources of Lorentz violation in a synchrotron process that are not significantly further suppressed. However, the transformation that eliminates $`k_F`$ is frame-dependent, and the new coordinates need not even be rectangular relative to the original ones; so by choosing to consider only this form of Lorentz violation, we are restricting ourselves to working in a very particular and special coordinate system.
We know that the Lorentz-violating coefficients for any physical charged particles are small. Physically, we might expect that the characteristic size for $`c^{\nu \mu }`$ is $`𝒪(m/M_P)`$. However, we shall not make any special use of the fact that $`M_P`$ is the Planck scale. Rather, we may take this size estimate as effectively defining $`M_P`$; $`M_P`$ is whatever large energy scale is needed in order to give the $`c`$ terms the correct magnitude.
The canonical quantization of the fermion field requires some care when the $`c`$ coefficients are nonvanishing. If $`c^{\nu 0}`$ is nonzero, then $``$ will contain nonstandard time derivative terms. In this case, a matrix transformation $`\psi R\psi `$ will be required, to ensure that $`\mathrm{\Gamma }^0=\gamma ^0`$. An explicit power-series expression for the required $`R`$ is given in . For simplicity, we shall assume that any such necessary transformation has already been performed and $`c^{\nu 0}=0`$. However, this will require us to consider the canonical quantization in a single frame only. We may not boost the theory into another frame, because doing so would reintroduce the problematic time derivatives.
In fact, in much of what follows, we shall neglect the $`c^{0\mu }`$ terms as well. While there is no special reason to believe this, we shall assume that the Lorentz violation is purely spacelike in a frame in which $`F^{\mu \nu }`$ contains only a magnetic component. We do this because the problem can then be solved exactly, to all orders in the remaining Lorentz-violating coefficients. However, when we revert to the linearized approximation and derive a limit on the $`c`$ coefficients from the observed properties of the synchrotron spectrum, we shall include the $`c^{0\mu }`$ parts in the calculation.
We shall consider the interaction with the electromagnetic field in two stages. This is standard practice in consideration of cyclotron motion. First, we determine the path traced out by a nonradiating charged particle moving in a spatially homogeneous background magnetic field. Then we evaluate the radiation induced by this periodic motion.
Because we are interested in the synchrotron emission from a single particle, we shall make the natural approximation of treating the Dirac equation as a single-particle wave equation. Standard techniques of relativistic quantum mechanics then apply; however, Lorentz violation will introduce new complexities. The momentum and velocity are not simply related by $`\stackrel{}{v}=\stackrel{}{\pi }/E=\stackrel{}{\pi }/\gamma m`$. The single-particle fermion Hamiltonian derived from $``$ is
$$H=\alpha _j\pi _jc_{lj}\alpha _l\pi _jc_{0j}\pi _j+\beta m,$$
(3)
where $`\alpha _j=\gamma ^0\gamma ^j`$ and $`\beta =\gamma ^0`$ are the usual Dirac matrices, and $`\stackrel{}{\pi }`$ is the mechanical (rather than canonical) three-momentum $`\stackrel{}{\pi }=\stackrel{}{p}e\stackrel{}{A}`$. Time derivatives of operators relating to fermion properties may be found by taking commutators with this Hamiltonian. In particular, the derivative of the particle’s position is
$$\dot{x}_k=\alpha _kc_{lk}\alpha _lc_{0k}.$$
(4)
Similarly, the equation of motion for $`\alpha _k`$ may be written
$$\dot{\alpha }_k=i\left[2\alpha _k(H+c_{0j}\pi _j)+2\pi _k2c_{kj}\pi _j\right],$$
(5)
which has the exact solution
$$\alpha _k(t)=\left(\pi _kc_{kj}\pi _j\right)(H+c_{0j}\pi _j)^1+\left[\alpha _k(0)\left(\pi _kc_{kj}\pi _j\right)(H+c_{0j}\pi _j)^1\right]e^{2i(H+c_{0j}\pi _j)t}.$$
(6)
The second term on the right-hand-side of (6) is matrix-valued and oscillatory. This term describes the particle’s Zitterbewegung, which, for a well-localized wave packet, is generated by interference between positive and negative frequency (i.e. particle and antiparticle) modes. The first term, when combined with (4) gives the bulk velocity
$$v_k=\frac{1}{E+c_{0j}\pi _j}\left(\pi _kc_{kj}\pi _jc_{jk}\pi _j+c_{jk}c_{jl}\pi _l\right)c_{0k}.$$
(7)
This same expression can also be found by calculating the group velocity $`\stackrel{}{v}_g=\stackrel{}{}_\stackrel{}{\pi }E`$. However, as previously stated, we will drop the $`c_{0j}`$ contributions in much of the following and use
$$v_k=\frac{1}{E}\left(\pi _kc_{kj}\pi _jc_{jk}\pi _j+c_{jk}c_{jl}\pi _l\right).$$
(8)
If rotation invariance is unbroken in the inertial frame we are considering, so that $`c_{jk}\delta _{jk}`$, there is merely a rescaling of the velocity. This will lead to fewer interesting effects. Fortunately, even if there is a privileged frame in which $`c_{jk}\delta _{jk}`$, a body emitting synchrotron radiation will not generally be at rest in this frame. We shall therefore assume that there is some breaking of rotation invariance in the rest frame of the source. In general, we shall assume that there is no suppression of rotation invariance violation relative to boost invariance violation only.
The equation of motion for the particle is the unmodified Lorentz force law $`\dot{\stackrel{}{\pi }}=e\dot{\stackrel{}{x}}\times \stackrel{}{B}`$. We shall neglect the Zitterbewegung in $`\dot{\stackrel{}{x}}`$ and consider only
$$\dot{\stackrel{}{\pi }}=e\stackrel{}{v}\times \stackrel{}{B}.$$
(9)
Then, since according to (8) $`\stackrel{}{v}`$ remains a linear function of the momentum, we may solve for the particle’s motion exactly. Again, we emphasize that all these same results for the bulk velocity and equation of motion also apply to Klein-Gordon particles, although the Klein-Gordon equation is even less satisfactory as a single-particle wave equation than is the Dirac equation.
It is advantageous to solve for the time development of the velocity, rather than the momentum (canonical or mechanical). While the velocity is a less fundamental object, the greatest formal problem with it—the Zitterbewegung—has already been neglected. The momentum could be determined with equal ease, but it possesses the unattractive property that a particle could possess zero momentum, yet not be stationary, if a $`c_{0j}`$ term were present. Moreover, even with the Lorentz-violation, the electromagnetic field is coupled directly to the velocity. Because $`c^{\nu 0}=0`$, the electrostatic potential $`\mathrm{\Phi }=A^0`$ is coupled, as usual, to the charge density $`e\psi ^{}\psi `$. Similarly, the vector potential $`\stackrel{}{A}`$ couples to $`e\psi ^{}\dot{\stackrel{}{x}}\psi `$, where $`\dot{\stackrel{}{x}}`$ is given by (4). Neglecting the Zitterbewegung, the coupling is simply to the bulk velocity $`\stackrel{}{v}`$. The fact that that the electromagnetic coupling is standard in this way was already evident in the Lorentz force law (9), and it holds equally in the equations of motion for $`A^\mu `$.
To determine the particle’s motion, we must solve a set of two coupled differential equations. These two equations describe the time evolution of the two components of the velocity in the plane perpendicular to $`\stackrel{}{B}`$; the component of $`\stackrel{}{v}`$ parallel to $`\stackrel{}{B}`$ does not contribute to the $`\stackrel{}{v}\times \stackrel{}{B}`$ force. Let us take $`\stackrel{}{B}`$ to point along the $`z`$-direction, $`\stackrel{}{B}=B\widehat{z}`$. Then the equations of motion for $`\stackrel{}{\pi }`$ are
$`\dot{\pi }_1`$ $`=`$ $`eBv_2`$ (10)
$`\dot{\pi }_2`$ $`=`$ $`eBv_1`$ (11)
$`\dot{\pi }_3`$ $`=`$ $`0.`$ (12)
So $`\pi _3`$ is a constant of the motion, as is $`E=\sqrt{m^2+\left(\pi _kc_{kj}\pi _j\right)\left(\pi _kc_{kl}\pi _l\right)}`$. Differentiating (8) then gives the following equations of motion for $`v_1`$ and $`v_2`$
$`\dot{v}_1`$ $`=`$ $`{\displaystyle \frac{1}{E}}\left[\left(12c_{11}+c_{j1}c_{j1}\right)\dot{\pi }_1+\left(c_{12}c_{21}+c_{j1}c_{j2}\right)\dot{\pi }_2\right]`$ (13)
$`\dot{v}_2`$ $`=`$ $`{\displaystyle \frac{1}{E}}\left[\left(c_{12}c_{21}+c_{j1}c_{j2}\right)\dot{\pi }_1+\left(12c_{22}+c_{j2}c_{j2}\right)\dot{\pi }_2\right].`$ (14)
Combining equations (10), (11), (13), and (14) in matrix form gives
$$\left[\begin{array}{c}\dot{v}_1\\ \dot{v}_2\end{array}\right]=\frac{eB}{E}\left[\begin{array}{cc}\beta & \alpha \\ \gamma & \beta \end{array}\right]\left[\begin{array}{c}v_1\\ v_2\end{array}\right]=\omega _0M\left[\begin{array}{c}v_1\\ v_2\end{array}\right].$$
(15)
The elements of the matrix $`M`$ are $`\alpha =\left(12c_{11}+c_{j1}c_{j1}\right)`$, $`\beta =\left(c_{12}c_{21}+c_{j1}c_{j2}\right)`$, and $`\gamma =\left(12c_{22}+c_{j2}c_{j2}\right)`$, and $`\omega _0=\frac{eB}{E}`$.
The equation (15) is easily solved. Since $`M^2=(\alpha \gamma \beta ^2)I`$ (where $`I`$ is the identity matrix), $`e^{\omega _0Mt}=I\mathrm{cos}\omega t+\frac{\omega _0}{\omega }M\mathrm{sin}\omega t`$, where $`\omega =\omega _0\sqrt{\alpha \gamma \beta ^2}`$. For vanishing $`c^{\nu \mu }`$, $`\omega =\omega _0`$ is the usual synchrotron frequency. If we choose coordinates so that the initial conditions are $`v_1(t=0)=v_{10}`$ and $`v_2(0)=0`$, then
$$\left[\begin{array}{c}v_1(t)\\ v_2(t)\end{array}\right]=e^{Mt}\left[\begin{array}{c}v_{10}\\ 0\end{array}\right]=v_{10}\left[\begin{array}{c}\mathrm{cos}\omega t+\frac{\omega _0}{\omega }\left(c_{12}+c_{21}c_{j1}c_{j2}\right)\mathrm{sin}\omega t\\ \frac{\omega _0}{\omega }\left(12c_{22}+c_{j2}c_{j2}\right)\mathrm{sin}\omega t\end{array}\right].$$
(16)
The velocity in the $`z`$-direction can be found by direct integration of its derivative,
$$\dot{v}_3=\frac{eB}{E}\left[\left(c_{13}c_{31}+c_{j1}c_{j3}\right)v_2+\left(c_{23}+c_{32}c_{j2}c_{j3}\right)v_1\right].$$
(17)
So, if $`v_3(0)=v_{30}`$,
$`v_3(t)`$ $`=`$ $`v_{10}{\displaystyle \frac{\omega _0}{\omega }}\{(c_{23}+c_{32}c_{j2}c_{j3})[\mathrm{sin}\omega t\left({\displaystyle \frac{\omega _0}{\omega }}\right)(c_{12}+c_{21}c_{j1}c_{j2})(\mathrm{cos}\omega t1)]`$ (18)
$`\left({\displaystyle \frac{\omega _0}{\omega }}\right)(c_{13}+c_{31}c_{j1}c_{j3})(12c_{22}+c_{j2}c_{j2})(\mathrm{cos}\omega t1)\}+v_{30}.`$
The particle moves in an elliptical helix; there is a constant drift parallel to $`\stackrel{}{B}`$, superimposed upon an additional periodic motion. If the drift vanishes, then the orbit lies close to, but is not generally in, the plane normal to the magnetic field, because $`v_3`$ does not generally vanish, even if its mean value does.
## 3 Radiation Emission
We shall now move on to the second stage our calculation. We have the particle’s motion prescribed, so we may study the radiation emitted during this motion. For simplicity, we shall consider only the case in which the drift velocity is zero. \[This does not correspond to $`v_{30}=0`$, because there are additional time-independent terms in (18). Instead, the sum of these constant terms must vanish.\] However, since we have now formulated the problem in terms of a particle of prescribed velocity conventionally coupled to the radiation field, normal boosting techniques can be used to generalize these results to a situation in which the time-averaged velocity in the $`z`$-direction is nonvanishing. A crucial quantity to calculate is the speed of the particle, $`|\stackrel{}{v}|`$, which is given by
$`\stackrel{}{v}^2`$ $`=`$ $`{\displaystyle \frac{v_{10}^2}{2}}\left[\left(\eta +\xi \right)+\left(\eta \xi \right)\mathrm{cos}2\omega t+\zeta \mathrm{sin}2\omega t\right]`$ (19)
$`=`$ $`{\displaystyle \frac{v_{10}^2}{2}}\left[\left(\eta +\xi \right)+\sqrt{(\eta \xi )^2+\zeta ^2}\mathrm{cos}\left(2\omega t2\varphi \right)\right],`$ (20)
where $`\mathrm{tan}2\varphi =\zeta /(\xi \eta )`$ and the constants $`\eta `$, $`\xi `$, and $`\zeta `$ are
$`\eta `$ $`=`$ $`1+\left({\displaystyle \frac{\omega _0}{\omega }}\right)^4[(c_{23}+c_{32}c_{j2}c_{j3})(c_{12}+c_{21}c_{j1}c_{j2})`$ (21)
$`+(c_{13}+c_{31}c_{j1}c_{j3})(12c_{22}+c_{j2}c_{j2})]^2`$
$`\xi `$ $`=`$ $`\left({\displaystyle \frac{\omega _0}{\omega }}\right)^2\left[(c_{12}+c_{21}c_{j1}c_{j2})^2+(12c_{22}+c_{j2}c_{j2})^2+(c_{23}+c_{32}c_{j2}c_{j3})^2\right]`$ (22)
$`\zeta `$ $`=`$ $`2{\displaystyle \frac{\omega _0}{\omega }}(c_{12}+c_{21}c_{j1}c_{j2})2\left({\displaystyle \frac{\omega _0}{\omega }}\right)^3(c_{23}+c_{32}c_{j2}c_{j3})`$
$`\times [(c_{23}+c_{32}c_{j2}c_{j3})(c_{12}+c_{21}c_{j1}c_{j2})+(c_{13}+c_{31}c_{j1}c_{j3})(12c_{22}+c_{j2}c_{j2})].`$
Thus far, our results have been exact, except that we have neglected the radiative reaction force. Henceforth, we shall be making use of the standard, Lorentz-invariant results on the power radiated by a particle undergoing synchrotron motion . However, the standard methodology for evaluating synchrotron emission involves a number of approximations. One often neglects any effect suppressed by a positive power of the Lorentz factor $`\gamma `$, and we shall follow this prescription. Among the things we may therefore neglect is the radiation due to the component of the acceleration parallel to the velocity; this contribution to the emission is small in comparison with that arising from the perpendicular component of the acceleration. We may also ignore the angular width of the radiation beam. All the emitted energy is beamed into a narrow pencil of angles centered around the instantaneous direction of the velocity. The range of angles covered is $`𝒪(\gamma ^1)`$, but we shall neglect this spread, instead assuming that all radiation is emitted along a ray tangent to the particle’s path.
We shall also neglect the Lorentz violation as a source of angular deviation. Although the exact orbit is neither circular nor in the plane normal to $`\stackrel{}{B}`$, the deviations from the conventional trajectory are small, of $`𝒪(c)`$. It would not be feasible to measure changes in the angular distribution of the emitted radiation induced by the presence of the Lorentz violation. We shall therefore neglect the changes in the orbital shape. All effects we shall consider will therefore be related to the modification of $`|\stackrel{}{v}|`$ (20). (This is similar to the approach adopted in , where the magnitude of the velocity was also taken as the central quantity.) As the velocity changes around the particle’s nearly circular path, the rate at which radiation is emitted will vary. The most sensitive tests of $`c`$-type Lorentz violation in synchrotron radiation could come from comparing the power output in different directions. (Unfortunately, such measurements are obviously not possible for single astrophysical sources.)
The phase $`\varphi `$ represents the angular position of the particle in its orbit at the time when the velocity is a maximum. At the antipodal point of the orbit, the velocity is also maximal. The greatest amount of radiation is then emitted along the tangent rays at these two points and propagates in the directions given by the azimuthal angles $`\varphi \pm \frac{\pi }{2}`$. Similarly, the smallest radiated power is in the directions $`\varphi `$ and $`\varphi +\pi `$. The presence of this effect is of course dependent on the existence of rotation invariance violation.
Neglecting radiation due to the component of the acceleration parallel to the velocity \[which is smaller by a factor of $`𝒪(\gamma ^2)`$\], the intensity spectrum per unit spectral frequency $`\omega _s`$ is
$$\frac{dI}{d\omega _s}=\sqrt{3}e^2\gamma \frac{\omega _s}{\omega _c}_{\omega _s/\omega _c}^{\mathrm{}}𝑑xK_{5/3}(x).$$
(24)
The critical frequency is $`\omega _c=\frac{3}{2}\gamma ^3\rho ^1`$, and $`\rho `$ is the instantaneous radius of curvature of the orbit $`\rho =\stackrel{}{v}^2/|\stackrel{}{a}^{}|`$, where $`\stackrel{}{a}^{}`$ is the component of the acceleration perpendicular to $`\stackrel{}{v}`$, $`\stackrel{}{a}^{}=\dot{\stackrel{}{v}}\frac{\dot{\stackrel{}{v}}\stackrel{}{v}}{\stackrel{}{v}^2}\stackrel{}{v}`$. Neglecting the Lorentz-violating corrections, $`\rho `$ is approximately $`E/|e|B`$. $`K_{5/3}(x)`$ is a modified Bessel function of the second kind. One could go further and calculate the radiation fields in the far field explicitly. However, we shall not do this, because the polarization structure of the emitted radiation is not substantially effected by the Lorentz violation.
For ultrarelativistic particles, for which $`1|\stackrel{}{v}|1`$, the Lorentz factor is roughly $`\gamma 1/\sqrt{2(1|\stackrel{}{v}|)}`$, and this is a rapidly increasing function of the speed—$`d\gamma /d|\stackrel{}{v}|=|\stackrel{}{v}|\gamma ^3\gamma ^3`$. The description of the Lorentz violation through an effective field theory containing only $`c^{\nu \mu }`$ terms will break down if the modifications of the velocity due to the presence of $`c`$ can render the speed superluminal. According to (8), this can occur when $`|\stackrel{}{\pi }|/E1|c|`$, where $`|c|`$ is a characteristic size for the Lorentz-violating coefficients. This gives us an estimate of the maximum value of $`\gamma `$ that can be achieved before new physics must come into play if some form of causality is to be preserved: $`\gamma _{\mathrm{max}}1/\sqrt{|c|}`$. This corresponds to an energy scale $`E_{\mathrm{max}}\sqrt{mM_P}`$.
## 4 Prospects for Observability
It is still unclear whether the changes we have described in the emission will be observable, and we shall now turn our attention to this issue. The total radiated synchrotron power in the ultrarelativistic regime is proportional to $`\gamma ^4`$. Therefore, the change in the radiated power as the velocity varies around the orbit is given by
$$\frac{\mathrm{\Delta }P}{P}2\frac{\left[d\left(\gamma ^4\right)/d|\stackrel{}{v}|\right]|c|}{\gamma ^4}8\gamma ^2|c|.$$
(25)
The factor of 2 comes from the fact that the deviations in $`|\stackrel{}{v}|`$ range over both positive and negative values. The characteristic size $`|c|`$ used in this calculation should be essentially the same as that used in the determination of $`\gamma _{\mathrm{max}}`$, because in both cases, $`|c|`$ measures the magnitude of the contribution that $`c^{\nu \mu }`$ can make to the velocity. Inserting $`\gamma =\gamma _{\mathrm{max}}`$ into (25) gives a result that is greater than unity. This means that the fractional change in the emitted power can be of order one in the regime in which the theory is valid; we do not have to go to an energy scale so high that new physics must emerge in order to see changes in the emission. However, we do need to get comparatively close to the scale at which the theory breaks down in order to observe deviations in the spectrum.
In fact, for $`|c|10^{19}`$, we will find a $`\frac{\mathrm{\Delta }P}{P}`$ of one percent at $`\gamma 10^8`$. For the lightest charged particle, the electron, this corresponds to an energy of roughly 50 TeV, orders of magnitude beyond anything one could create in the laboratory. We conclude that these effects are unobservable for Earth-based sources.
The only sources of synchrotron radiation that are high enough in energy to give observable results of the type we are considering are astrophysical. However, as each astrophysical source can only be observed from a single direction, more than one source would be required in order to make the kind of directional observations that could constrain $`c`$ most strongly. Ideally, we would want to have two or more very clean sources of synchrotron radiation, for which the spectra due to the motion of multiple species of particles (e.g., both electrons and protons) could be resolved. Then we could look for systematic differences between the emission profiles for the species. This would mean effectively using the proton spectra, for example, as local magnetometers and looking to see whether the electron spectra are consistent with the measured fields. One could then set bounds on a combination of the $`c`$ coefficients for the electron and the proton.
Although observations of distant synchrotron sources (such as far-off radio galaxies) are ideal for constraining Lorentz violation in the photon sector, they are not so helpful here. A long line of sight will magnify small effects that modify the propagation structure of radiation. However, large distances do nothing to assist measurements of Lorentz violation in the charged emitters themselves. A nearby, accurately understood source is better than a distant one.
The best-understood synchrotron source is the Crab nebula, but its spectrum still appears too complicated for the kind of procedure we have suggested to be at all feasible. (For a good review of the Crab nebula’s nonthermal emission spectrum, see .) For example, the spectrum contains two different electron synchrotron components, with significantly different characteristics. Any observed proton synchrotron radiation would probably fail as a sensitive magnetometer, because it would be impossible to associate it uniquely and in a model-independent fashion with either one or the other electron population. Based on measurements of the entire spectrum, the average strength of the magnetic field in the x-ray production region is known to be in the tenths of mG, but it is not known to high accuracy. There are also large relative uncertainties in the radiation rates in some regimes, particularly the highest energy.
However, we still can get a strong constraint on $`c`$ from the Crab nebula data. This constraint, like the one derived in , is based upon the fact that Lorentz violation may give rise to a maximum particle velocity. The existence of electrons with large velocities then constrains the Lorentz-violating parameters. For a particle moving in the direction of a unit vector $`\stackrel{}{e}`$, the maximum allowed velocity is (to leading order in $`c`$) $`1c_{jk}e_je_kc_{0j}e_j`$. We observe via the Crab nebula synchrotron spectrum electrons with Lorentz factors as large as $`3\times 10^9`$. This means that the maximum velocity in the Crab-to-Earth direction is greater than $`16\times 10^{20}`$, hence $`c_{jk}e_je_k+c_{0j}e_j<6\times 10^{20}`$. As in , this is a one-sided limit; one sign of this combination of coefficients leads to a maximum velocity in the relevant direction, but the other does not.
The direction $`\stackrel{}{e}`$ can be transformed into the standard sun-centered celestial equatorial coordinates used in the study of Lorentz violation . The location of the Crab nebula is right ascension 5h 34m 32s, declination 22 0 52<sup>′′</sup>, lying close to the ecliptic plane. So the unit vector pointing from the nebula to the Earth has components $`e_X=0.10`$, $`e_Y=0.92`$, and $`e_Z=0.37`$. This gives us the particular elements of the $`c`$ tensor that are constrained by this measurement. The specific constraint is
$`[0.01c_{XX}+0.85c_{YY}+0.14c_{ZZ}+0.09c_{(XY)}+0.04c_{(XZ)}`$ (26)
$`+0.34c_{(YZ)}0.10c_{0X}0.92c_{0Y}0.37c_{0Z}]`$ $`<`$ $`6\times 10^{20},`$
where $`c_{(jk)}`$ is the symmetric sum $`c_{jk}+c_{kj}`$. Similar constraints could be obtained for other well-resolved synchrotron sources; this would provide further constraints on the symmetric part of $`c_{jk}`$ and on the $`c_{0j}`$. \[At leading order, the antisymmetric part of $`c_{jk}`$ just represents a change in the representation of the Dirac matrices, and it is already evident from (7) that it will not contribute to the velocity.\]
In order for these constraints to be valid, we must know that there are no other effects that will interfere with our result. In particular, we would like to address the question of whether quantum corrections would affect the emission before Lorentz-violating corrections become important. The leading order quantum corrections to the standard synchrotron formulas may be found by making the replacement $`\omega _s\omega _s\left(1+\frac{\omega _c}{E}\right)`$ in $`\frac{1}{\omega _s}\frac{dI}{d\omega _s}`$ . The corrections are negligible if $`\omega _cE`$, or equivalently if $`\gamma \frac{m^2}{|e|B}`$. This is $`\gamma \left(3\times 10^{13}\right)B^1`$ if the particle is an electron and $`B`$ is measured in Gauss. If a typical field strength is that within the Crab nebula, $`B0.2`$$`0.3`$ mG, the maximum values of $`\gamma `$ are extremely high. So our classical treatment could apply up to scales well above those at which we would expect to start seeing marked deviations from the conventional results.
Synchrotron radiation has already been used to set strong limits on nonrenormalizable Lorentz-violating modifications of quantum electrodynamics. Lorentz violation in synchrotron radiation is also theoretically interesting, and there have been a number of prior analyses of the emission spectrum in specific Lorentz-violating models. In this paper, we have looked at the impact of the renormalizable SME coefficient $`c^{\nu \mu }`$ on synchrotron processes. Although we have used a number of standard approximations to simply our analysis of the radiation, no approximations relating to the Lorentz violation were required; the expressions (16) and (18) are exact, valid to all orders in $`c`$. The $`c^{\nu \mu }`$ coefficients for electrons, particular the diagonal coefficients, can be difficult to bound experimentally . Since only a single fermion is involved in synchrotron radiation, this is a process in which it is relatively easy to isolate electron-specific effects, and the kind of constraints we have obtained here should prove useful.
## Acknowledgments
The author is grateful to V. A. Kostelecký for helpful discussions. This work is supported in part by funds provided by the U. S. Department of Energy (D.O.E.) under cooperative research agreement DE-FG02-91ER40661. |
warning/0507/nucl-th0507070.html | ar5iv | text | # Introduction to statistical models and non-extensive statistics
## I Introduction
Statistical models in general apply to phenomena which appear in a sufficiently large number and can be repeated independently and (in principle) indefinitely. Whether high energy physics and in particular heavy-ion collisions provide such phenomena, is a question to be discussed on its own right. We shall try to contribute to an answer to this question here by reviewing some basic concepts of statistics, kinetic theory and thermodynamics, with the aim to clarify the limits of the application of statistical models.
First we review the central limit theorem of statisticsCENTRAL , as the basic law governing the result of very many independent influences. Then kinetic theory, designed to describe the way towards equilibrium and its maintenance, is discussed. By doing so we point out quite a few modern applications of statistical methods describing a stationary state, which obeys non-extensive thermodynamical rulesNEXT-THERMO . A review of non-extensive rules, equilibrium one-particle distributions and entropy density formulas follows.
An underlying particular application we have in mind is a relativistic heavy ion collisionRHIC . The first touch physics is probably dominated by nonlinear field dynamics and parton collisions. While some very energetic partons may escape in form of jets, in a central heavy-ion event most of the beam energy is transformed into a compression of the nuclear matter and production of relativistic quarks and gluons. How a quark-gluon plasma in a thermal state is formed in these events and with which properties, is still an objective of the contemporary research.
Assuming i) such a quark matter formation in an intermediate state, ii) ergodization of energy between many newly produced particles and iii) a relatively fast hadronization, the experimentally measured specific hadron spectra may reflect statistical, presumably even equilibrium thermal, properties of the precursor matterSTAT-MOD . Strong final state interaction on the other hand, a so called hadronic afterburner, may wash out spectral characteristics of earlier stages of the evolution. Whether this is the case, can in principle be studied by checking quark coalescence rules or comparing hadron- and lepton-spectraALCOR . Finally the late resonance decayRESONANCES during free streaming of hadrons changes the hadrochemical composition. Fortunately, some properties (e.g. the transverse momentum spectra) are influenced only partially (at their low end) by this. Since the relativistic energy is given by $`E=m_T\mathrm{cosh}y`$ with transverse mass $`m_T=\sqrt{p_T^2+m^2}`$ and rapidity $`y`$ for a particle with mass $`m`$, the best way to study statistical equilibrium distribution of hadrons is the comparison of $`m_T`$-spectra at rapidity $`y=0`$ for different particles. A universal behavior ($`m_T`$-scaling)MT-SCALING indicates that the one-particle distributions depend on the energy only and not on all momentum components: a basic feature of generalized and conventional thermal distributions.
In this lecture we review general basic ideas and mathematical formulas related to statistical models. By doing so emphasis is given to particular cases when the conventional picture of earlier textbook statistical physics does not apply, instead some generalized concepts have to be considered. Starting with the central limit theorem of statistics we continue with a discussion of the Brownian motion in the framework of the Langevin and Fokker-Planck equations. Here the very level of generality is identified that leads to a Tsallis distribution instead of the Gibbs-Boltzmann oneNEXT-MT-SCALING . The potential for further generalization is pointed out, too. Then a discussion of the Boltzmann equation follows introducing to modern generalizations capable to describe non-exponential equilibrium one-particle distributionsNLBE ; NEBE . We close with some remarks about how certain a thermodynamical state may be reconstructed from the observed one-particle energy spectra.
## II Statistics: the law of big numbers
There is a mathematical property behind the applicability of statistical physics: many ”normal” distributions of probability fold to a Gaussian with a width scaling down with increasing number of individual constituents. This property is expressed nicely in the central limit theoremCENTRAL . An enlightening simple example of its action is given by the distribution of the sum of uniform random variables in a finite interval. Finally in this section an example not falling under the reign of the central limit theorem, the Lorentzian distribution, is presentedALMOST-EXP ; LEVY .
Let $`x_i`$ be random distributed according to $`w_i(x_i)`$. We are interested in the distribution of a scaled sum of $`n`$ such variables:
$$P_n(x)=\underset{i=1}{\overset{n}{}}w_i(x_i)\delta \left(xa_n\underset{k=1}{\overset{n}{}}x_k\right).$$
(1)
Here we assumed that the joint probability of having $`n`$ values is a product of the individual probabilities; this is the requirement of statistical independency. With this assumption the Fourier transform of the seeked probability is a product of properly scaled Fourier transforms of individual probabilities:
$$\stackrel{~}{P}_n(k)=𝑑xe^{ikx}P_n(x)=\underset{i=1}{\overset{n}{}}\stackrel{~}{w}_i(a_nk).$$
(2)
From the Taylor expansion of $`\mathrm{ln}\stackrel{~}{P}(k)`$ around $`k=0`$ one obtains the central moments (correlations). For the $`\mathrm{}`$-th moment the scaling law, $`\sigma _n^{(\mathrm{})}=a_n^{\mathrm{}}n\overline{\sigma }^{(\mathrm{})}`$ applies with $`\overline{\sigma }`$ being the finite average of the given central moment of the individual distributions. Whenever $`\sigma _i^{(\mathrm{})}=0`$ for all $`\mathrm{}<\mathrm{}_0`$ values and $`\sigma _i^{(\mathrm{}_0)}`$ is finite, all the higher moments can be scaled down in the folded distribution by choosing $`a_nn^{1/\mathrm{}_0}`$. The resulted scaling, $`\sigma _n^{(\mathrm{})}=n^{1\mathrm{}/\mathrm{}_0}\overline{\sigma }^{(\mathrm{})}`$, in the $`n\mathrm{}`$ limit leaves us with only one nonzero central moment, the $`\mathrm{}_0`$-th one.
Usually the $`\mathrm{}=0`$ central moment is zero due to the normalization of the probability: $`\mathrm{ln}\stackrel{~}{P}(0)=\mathrm{ln}1=0`$. The first moment ($`\mathrm{}=1`$) can be made zero due to symmetry or by a simple shift in the variables $`x_i`$. The second central moment is then the first nonzero value. As a consequence the resulted distribution of the scaled sum has a second moment with all higher moments vanishing, therefore $`\mathrm{ln}\stackrel{~}{P}(k)`$ is quadratic in $`k`$, so $`\stackrel{~}{P}(k)`$ and with that $`P(x)`$ is Gaussian.
A nice, simple example is given for $`x_i`$-s uniformly distributed in the interval $`(1,+1)`$. The distribution of
$$x=\sqrt{\frac{3}{n}}\underset{i=1}{\overset{n}{}}x_i$$
(3)
approaches the Gaussian:
$$\underset{n\mathrm{}}{lim}\stackrel{~}{P}_n(k)=\underset{n\mathrm{}}{lim}\left(\frac{\mathrm{sin}(k\sqrt{3/n})}{k\sqrt{3/n}}\right)^n=\mathrm{exp}(k^2/2).$$
(4)
A counterexample is given by the Lorentzian distribution, having a Fourier transform $`\stackrel{~}{w}_i(k)=\mathrm{exp}(|k|)`$. Now there is a non-Gaussian limiting distribution with an altered scaling for
$$x=\frac{1}{n}\underset{i=1}{\overset{n}{}}x_i.$$
(5)
It is itself a Lorentzian:
$$\underset{n\mathrm{}}{lim}\stackrel{~}{P}_n(k)=\underset{n\mathrm{}}{lim}\left(e^{|k|/n}\right)^n=\mathrm{exp}(|k|).$$
(6)
This is a special case of the Lévy distributionLEVY .
## III Kinetic theory
In a brief review of concepts distilled from the kinetic theory approach to thermodynamics we shall rely on the notion of ”noise” heavily. The sum of random influences (forces) is itself a random variable. It shows relationship to the sum of random numbers; its distribution under quite general circumstances can be considered as a Gaussian distribution. The independency of individual influences is assumed first of all in time: such a view deals with uncorrelated stochastic changes in the parameters describing a physical system.
The simplest physical theory of such a system is that of the Brownian motionBROWN , considering a free, massive particle under the influence of forces acting as an uncorrelated, Gaussian (white) noise. A description of the motion of such a particle is given by the classical Langevin equation, or equivalently a statistics over possible such motions is described by the Fokker-Planck equation. A balance between damping and accelerating forces leads to a stationary state, with vivid microscopical dynamics, but macroscopically (on the average over many particles) presenting a thermodynamical equilibrium state. General statements about this balance are comprised in the fluctuation-dissipation theorem. In the next section we review generalizations of the Boltzmann equation, a somewhat more complicated kinetic theory. The effects of such generalizations on the the concept and definition of entropy and on equilibrium distributions will be an issue of a further section.
### III.1 General Langevin problem
Let us consider a simple, one degree of freedom motion. The change of momentum $`p`$ in time is given by a force depending on this momentum and on a noise variable $`z`$:
$$\dot{p}=F(p,z).$$
(7)
Following the method pioneered by Ornstein and UhlenbeckUHLENBECK a distribution of many possible values of $`p`$ at a time $`t`$ is considered. This $`f(p,t)`$ distribution governs average values of a smooth but otherwise arbitrary test function $`R(p)`$. The same integral over $`p`$ can be expressed at the time $`t+dt`$ assuming a statistical average over the noise $`z`$ in the time interval passed since $`t`$:
$$𝑑pR(p)f(p,t+dt)=𝑑pR(p+dtF(p,z))f(p,t).$$
(8)
One assumes that the averaging of the force $`F`$ over the noise $`z`$ gives the result:
$`F`$ $`=`$ $`G(p),`$
$`FFFF`$ $`=`$ $`2D(p)/dt.`$ (9)
The above scaling of the correlation with $`dt`$ follows from the Gaussian nature of the noise $`z`$. Expanding now the equation(8) up to terms linear in $`dt`$ one arrives at:
$$𝑑pR(p)\frac{f}{t}(p,t)=𝑑p\left[G(p)R^{}(p)+D(p)R^{\prime \prime }(p)\right]f(p,t).$$
(10)
After partial integration and considering arbitrary $`R(p)`$ one gets the Fokker-Planck equation:
$$\frac{f}{t}=\frac{}{p}\left(Gf\right)+\frac{^2}{p^2}\left(Df\right).$$
(11)
### III.2 Particular Langevin problem
The above Langevin and Fokker-Planck problem is still quite general. Damping and diffusion coefficients depend on the momentum $`p`$ in a general way. Ergodization in phase space is achieved on the other hand when constant energy surfaces are covered. In such a situation the distribution $`f`$, as well as the coefficients $`G`$ and $`D`$ (the latter related to the noise), depends on the energy $`E(p)`$ only. Note, however, that they are not constant.
Such a particular Langevin equation is given bySLIDING-SLOPE
$$\dot{p}=zG(E)\frac{E}{p}$$
(12)
containing an energy dependent damping proportional to the general velocity $`v=E/p`$, and a zero-average noise, $`z(t)=0`$, with a correlation
$$z(t)z(t^{})=2D(E)\delta (tt^{}).$$
(13)
The Fokker-Planck equation contains in this case the factors $`D(p)=D(E)`$ and $`G(p)=G(E)E/p`$. Its stationary solution is given by
$$f(p)=\frac{A}{D(E)}\mathrm{exp}\left(\frac{G(E)}{D(E)}𝑑E\right)=A\mathrm{exp}\left(\frac{dE}{𝒯(E)}\right).$$
(14)
This result is not readily the Boltzmann-Gibbs distribution, $`\mathrm{exp}(E/T)`$, only in the case of energy-independent damping and noise coefficients. A general equilibrium is able to feature almost any other distribution of the energy of a single degree of freedom picked out of its environment. Instead of a constant temperature, $`T`$, in the general case a sliding inverse logarithmic slope is characteristic to such states. From $`1/𝒯(E)=d\mathrm{ln}f(E)/dE`$ its relation to the damping and diffusion coefficients follows:
$$𝒯(E)=\frac{D(E)}{G(E)+D^{}(E)}.$$
(15)
The low-energy limit of this expression, pretending as $`G(E)`$ and $`D(E)`$ were constant, leads to an experimentally feasible definition of the Gibbs temperature, $`T_{Gibbs}=D(0)/G(0)`$. From the viewpoint of the Brownian motion another temperature may be used, the Einstein temperature $`T_{Einstein}=lim_E\mathrm{}D(E)/G(E)`$. None of these two approximations are, however, coincident with the sliding slope of particle spectra given by eq.(15).
An important and historically mostly considered particular case is given by the constant slope distribution, the Boltzmann-Gibbs distribution: $`T(E)=T`$. A modern, non-classical distribution, the Tsallis distribution seems to be the next simplest, having a linear inverse slope – energy relation:
$$𝒯(E)=T/q+E(11/q).$$
(16)
The corresponding equilibrium distribution in this case turns to be an exponential of a logarithm, which is a power-law:
$$f(p)=\frac{1}{Z}\left(1+(q1)\frac{E}{T}\right)^{\frac{q}{q1}}$$
(17)
It has the interesting property, that the parameter $`T`$ is the fixed point of the sliding (linear) slope: $`𝒯(T)=T`$. This parameter shall be referred to as the Tsallis temperature.
### III.3 Fluctuation-dissipation theorem
In a realistic system there are many microscopic degrees of freedom to be considered. Denoting a point in the $`6N`$-dimensional phase space by $`p_i`$, the Langevin equation can be implemented in the form
$$\dot{p}_i=(S_{ij}G_{ij})_jE+z_i.$$
(18)
The symplectic coefficient, $`S_{ij}`$, does not change the total energy of the system, $`E(p_i)`$, it causes conservative motion inside a given energy shell only. For the sake of energy distribution it is therefore not interesting and shall be omitted in the followings. The damping and dissipation terms with the respective symmetric coefficients $`G_{ij}`$ and $`D_{ij}`$ keep balance on the long term. An ergodized equilibrium distribution can be a single function of the energy $`E`$ only, therefore these coefficient matrices have to be connected by a single function of energy, too. This gives rise to a general fluctuation dissipation theorem:
$$D_{ij}(E)=𝒯(E)\left(G_{ij}(E)+D_{ij}^{}(E)\right).$$
(19)
It is highly nontrivial that two high-dimensional matrix functions of the phase space coordinates $`p_i`$ would be related by a single scalar function of energy only! Recalling that $`𝒯(E)`$ is the inverse logarithmic slope of the equilibrium distribution, the fluctuation dissipation relation can be expressed using $`f(E)`$, too:
$$D_{ij}(E)=\frac{1}{f(E)}\underset{E}{\overset{\mathrm{}}{}}G_{ij}(x)f(x)𝑑x.$$
(20)
Still, quite general diffusion and damping coefficient matrices are allowed, but $`D_{ij}`$ is connected to $`G_{ij}`$ via an energy dependent scalar, the equilibrium energy distribution, $`f(E)`$. Particular cases of this relation are i) the use of the Gibbs distribution, $`f(E)\mathrm{exp}(E/T)`$, leading to $`D_{ij}=TG_{ij}`$ with constant matrices (the usual textbook case), or ii) the use of a Tsallis distribution (17) giving rise to $`D_{ij}(E)=\left(T+(q1)E\right)G_{ij}`$ with constant $`G_{ij}`$ but linearly energy dependent diffusion coefficient, $`D_{ij}`$. A further interesting case is presented by assuming that both the damping and the diffusion coefficient matrix is a linear function of the energy, but their ratio (the Einstein temperature) is constant. This assumption is typical to field theory calculations. Due to eq.(19) the sliding slope is, however, not constant; it rather interpolates between a linear rise at low energy and a constant at high energy. From $`D=TG=\gamma T(1+E/E_c)`$ it follows $`1/𝒯(E)=\mathrm{\hspace{0.25em}1}/T+1/(E+E_c)`$.
### III.4 Additive and multiplicative noise
Another approach to reach a non-exponential stationary energy distribution from a kinetics described by the Langevin equation considers a stochastic damping coefficient. The Langevin equation is kept linear,
$$\dot{p}=\zeta \gamma p,$$
(21)
but, more general than in the classical approach, both $`\zeta `$ and $`\gamma `$ are stochastic variablesMULT-NOISE . With constant mean values, $`\zeta =F`$ and $`\gamma =G`$, and white-noise correlations,
$$\gamma (t)\gamma (t^{})=2C\delta (tt^{}),\zeta (t)\gamma (t^{})=2B\delta (tt^{}),\zeta (t)\zeta (t^{})=2D\delta (tt^{}),$$
(22)
the equivalent Fokker-Planck equation (11) contains a damping factor $`G(p)=GpF`$ and a diffusion factor $`D(p)=Cp^22Bp+D`$. For $`B0`$ the two noisy coefficients are cross-correlated. For a single degree of freedom the stationary distribution can be obtained analytically:
$$f(p)=A\left(\frac{D}{D(p)}\right)^v\mathrm{exp}\left(\frac{\alpha }{\theta }\mathrm{atan}(\frac{p\theta }{DBp})\right)$$
(23)
with the power $`v=1+G/2C`$, the exponent factor $`\alpha =GB/CF`$ and the variable $`\theta =\sqrt{DCB^2}`$. For $`F=0`$ a characteristic momentum scale in this distribution is given by $`p_c^2=D/C`$, the ratio of the additive and multiplicative noise widths. For $`F=0`$ and $`B=0`$, i.e. for two independent noises, the Tsallis distribution arises as a stationary solution:
$$f(p)=A\left(1+\frac{C}{D}p^2\right)^v.$$
(24)
Utilizing the energy formula for a free, massive, non-relativistic particle it reads as
$$f(p)=A\left(1+(q1)\frac{E}{T}\right)^{\frac{q}{q1}}$$
(25)
with the Tsallis index $`q=1+2C/G`$ and the Tsallis temperature coinciding with that of the classical Brownian motion, $`T=D/mG`$. In the small momentum limit, $`pp_c`$ the Tsallis distribution is nearly Gaussian, independent from the properties of the multiplicative noise, $`f=A\mathrm{exp}(Gp^2/2D)`$ (or expressed with the energy a Gibbs distribution, $`f=A\mathrm{exp}(E/T)`$). In the opposite limit at high energy it is a power-law distribution, $`f=A(p/p_c)^{2v}=A(E/E_c)^v`$. It is interesting to note, that a definite relation arises between the energy scale, $`E_c`$, the temperature parameter, $`T`$, and the tail power, $`v`$:
$$v=1+E_c/T.$$
(26)
This relation can be experimentally tested.
## IV Non-Extensive Boltzmann Equation
The heart of kinetic theory is the classical Boltzmann equation. It does not only describe dynamical evolution of large systems microscopically, resulting in stationary distributions on which thermodynamics can be based on, but it also offers a microscopical foundation to the key quantity entropy. It has, of course, also quite a few assumptions built in the theory; dropping one or other of them may lead to a generalization of the classical approach. The most recognized assumption, the micro-reversibility of the transition probability, establishes H-theorem and the definition of entropy; it should not be dropped. Less explicit assumptions, like taking the two-particle probability as a product of one-particle probabilities, or taking the total energy of a colliding pair as the sum of the respective one-particle energies for freely moving (asymptotic) particles, can be more readily generalized.
We review these two generalizations of the Boltzmann equation in this section: the generalization of the product rule for probabilities (dropping statistical independency) leads to a non-linear Boltzmann equation (NLBENLBE ), while considering two-particle energies composed by an extended addition rule mounds in the non-extensive Boltzmann equation (NEBENEBE ). Resulting stationary distributions, the H-theorem and the main characteristics of generalized thermodynamicses following from this will be presented.
### IV.1 NLBE: generalized product
The general structure of the Boltzmann equation describes the evolution of the one-particle phase space density (interpreted as finding probability of a particle or a microstate) via an integral over all possible transitions to and from other states:
$$\dot{f}_1=\underset{234}{}w_{1234}\left(f_{34,12}f_{12,34}\right).$$
(27)
Here the dot denotes a total time derivative (Vlasov operator) comprising the essential evolution of the one-particle phase space density, $`f_1`$. The indices $`1234`$ refer to two particles before and after a microcollision. The transition probability, $`w_{1234}`$ reflects microreversibility and partner symmetry in the permutation of these indices. It is positive and it contains some conditions on conserving physical quantities; at least momentum and energy.
$$w_{1234}=M_{1234}^2\delta ((\stackrel{}{p}_1+\stackrel{}{p}_2)(\stackrel{}{p}_3+\stackrel{}{p}_4))\delta (E_{12}E_{34}),$$
(28)
with $`E_{12}`$ total two-particle energy before and $`E_{34}`$ after the collision. The particle density factors, $`f_{12,34}`$ and $`f_{34,12}`$ weight the transition yields for a $`3+41+2`$ and for a $`1+23+4`$ process, respectively.
In the NLBE approach the traditional simple additivity of energy is kept,
$$E_{12}=E_1+E_2,$$
(29)
but the classical product formula of Boltzmann, $`f_{12,34}=f_1f_2`$, or the supplement with blocking factors due to Uehling and Uhlenbeck $`f_{12,34}=f_1f_2(1\pm f_3)(1\pm f_4)`$ is generalized. The generalization still reflects particle separation property (particle $`1`$ goes into particle $`3`$ and particle $`2`$ goes into particle $`4`$ classically), but abandons the linearity:
$$f_{12,34}=\gamma (f_1,f_3)\gamma (f_2,f_4).$$
(30)
This intends to simulate statistical correlations between initial and final states and a nonlinearity of transition yields. Further requirement is that the phase space density factor, $`\gamma `$ factorizes to a production factor $`a`$, to a blocking factor $`b`$ (depending on the respective initial and final phase space densities only) and to a factor symmetric in both:
$$\gamma (x,y)=a(x)b(y)c(x,y)$$
(31)
with $`c(x,y)=c(y,x)`$.
The stationary state of the NLBE is governed by the ratio $`\kappa (x)=a(x)/b(x)`$. It is easy to see from the following derivation of the generalized H-theorem. We seek for a quantity called entropy in form of a one-particle phase space integral, $`S=\underset{1}{}\sigma (f_1)`$. The question is what $`\sigma (f)`$ functional form guarantees macro-irreversibility, i.e. $`\dot{S}0`$. Using the general Boltzmann equation (27) one writes
$$\dot{S}=\underset{1}{}\dot{f}_1\sigma ^{}(f_1)=\underset{1234}{}w_{1234}c_{13}c_{24}\sigma _1^{}(a_3b_1a_4b_2a_1b_3a_2b_4),$$
(32)
with the corresponding indices referring to arguments of the functions $`a`$, $`b`$ and $`c`$. Now we explore particle permutation symmetries of this expression. After the exchange of particle $`1`$ with $`2`$ and simultaneously particle $`3`$ with $`4`$ we describe the same process. An exchange of the initial with the final state is done by $`13`$ and $`24`$ (micro-reversibility). It amounts to a relative minus sign in the total rate by exchanging gain and loss terms. Finally the combined operation of both also contributes with a minus sign in total. Dressing now the transition rate with factors carrying the same symmetry we use $`\stackrel{~}{w}_{1234}=w_{1234}c_{13}c_{24}b_1b_2b_3b_4/4`$. With the notation $`\kappa _i=a_i/b_i`$ for the production to blocking ratios we arrive at
$$\dot{S}=\underset{1234}{}\stackrel{~}{w}_{1234}\left(\sigma _1^{}+\sigma _2^{}\sigma _3^{}\sigma _4^{}\right)\left(\kappa _3\kappa _4\kappa _1\kappa _2\right).$$
(33)
This already reminds to the structure of the H-theorem result of the classical Boltzmann equation. The correct entropy density can be read off as satisfying
$$\sigma ^{}(f)=\mathrm{ln}\kappa (f).$$
(34)
The stable equilibrium, where – due to ergodization – $`f(p)`$ can be expressed as a solely function of the corresponding energy, satisfies a product rule for the $`\kappa (f_i)`$-s while the energy is additive. The only solution is that the canonical equilibrium of the NLBE is given by
$$\kappa (f)=\frac{1}{Z}\mathrm{exp}(E/T).$$
(35)
As particular cases the well-known Boltzmann-Gibbs, Fermi-Dirac or Bose-Einstein distributions are recovered, but this result is far more general. The corresponding general entropy formula can be constructed founding a generalized thermodynamics with (in general) non-extensive entropy composition rules when merging large subsystems.
### IV.2 NEBE: generalized sum
Another, recently pursued way is to keep the statistical independency, $`f_{12,34}=f_1f_2`$, but to generalize the energy addition formula to a nontrivial composition rule
$$E_{12}=h(E_1,E_2).$$
(36)
Rules not being a simple sum, $`h(x,y)x+y`$, present a non-extensive energy composition. Latest in the thermodynamical limit, only associative rules are physical: $`h(h(x,y),z)=h(x,h(y,z))`$. Due to a mathematical theoremMATH for functional equations the general solution of the associativity requirement is a strict monotonous mapping to the simple addition:
$$X(h)=X(x)+X(y).$$
(37)
This is unique up to a constant factor.
As a consequence for any microcollision $`X(E_1)+X(E_2)=X(E_3)+X(E_4)`$ holds. The stationary solution of NEBE is hence given by
$$f(p)=\frac{1}{Z}\mathrm{exp}(X(E)/T)$$
(38)
allowing again for a general, non-exponential energy dependence. It is noteworthy that to each associative non-extensive composition rule $`h(x,y)`$ there exist a mapping function $`X(E)`$. An in-medium dispersion relation, or quasi-energy is obtained this way: the total sum,
$$X(E_{tot})=\underset{i}{}X(E_i)$$
(39)
is conserved by the NEBE.
In order to connect the approaches NLBE and NEBE (and both with the non-extensive thermodynamics) the non-extensive composition rule of the energy has to be extended to composition rules of other, traditionally ”extensive” quantities. Most important is the entropy composition rule. And, as putting a corner stone into the right place, the scaling law
$$X(E)/T=X_s(E/T)$$
(40)
allows us to relate the equilibrium distributions by a one-variable functional relation,
$$Zf_{eq}(E)=\mathrm{exp}(X_s(E/T))=\kappa ^1(\mathrm{exp}(E/T)).$$
(41)
At the same time the generalized entropy density satisfies
$$\sigma ^{}(f)=\mathrm{ln}\kappa (Zf)=X_s^1(\mathrm{ln}Zf)$$
(42)
in equilibrium. This gives rise to an interpretation of the H-theorem for NEBE where the never-decreasing total entropy is given by the same strict monotonic back-mapping from the original Boltzmann-entropy:
$$S_B=X_s(S_{tot})=f\mathrm{ln}f.$$
(43)
The use of a non-extensive entropy composition rule, $`h_s(x,y)`$, mapped to additivity by $`X_s(t)`$ leads to the following mapping of the individual entropy-density: $`s(f)=fX_s^1(\mathrm{ln}f)`$. The usual non-extensive entropy is then defined by $`S_T=s(f)`$.
## V Thermodynamicses
It is enlightening to list some particular cases of non-extensive composition laws and corresponding thermodynamicsesNEBE+NLBE . The trivial law, $`h(x,y)=x+y`$, is mapped by the identity, $`X(E)=E`$, and in (canonical) equilibrium the Boltzmann-Gibbs distribution, $`\mathrm{exp}(E/T)`$, emerges. The entropy density is given by $`s(f)=f\mathrm{ln}f`$.
Tsallis’ non-extensive thermodynamics relies on the composition rule, $`h(x,y)=x+y+axy`$, connected to the mapping
$$X_s(S)=\frac{1}{aT}\mathrm{ln}\left(1+aTS\right),$$
(44)
first presented by AbeABE . The canonical equilibrium distribution,
$$f_{eq}(E)=\frac{1}{Z}\left(1+aE\right)^{1/aT},$$
(45)
can easily be connected to the Tsallis distribution (17) by $`q=1aT`$. Finally an entropy density following Tsallis’ original suggestion (encountered also earlier by several authorsOTHERS ): $`s(f)=\frac{f^qf}{q1}`$, can be obtained.
An interesting endeavor is to consider the composition rule, $`h(x,y)=\left(x^b+y^b\right)^{1/b}`$. Mapping the energy to its power, $`X(E)=(aE)^b/a`$, it gives rise in equilibrium to a Lévy distributionLEVY , a.o. known from anomalous diffusion problems:
$$f_{eq}(E)=\frac{1}{Z}\mathrm{exp}\left((aE)^b/aT\right).$$
(46)
The corresponding expression for the entropy is an incomplete gamma function.
The pure multiplicative law, $`h(x,y)=axy`$, which may be considered as the high-energy limit of Tsallis’ construction with fixed energy scale $`E_c=1/a`$, is mapped by $`X(E)=\frac{1}{a}\mathrm{ln}(aE)`$ and leads to a pure power-law behavior, $`Zf_{eq}=(aE)^{1/aT}`$. The entropy formula due to $`X_s(t)=(1/aT)\mathrm{ln}(aTt)`$, $`s(f)=\frac{1}{1q}f^q`$ reminds to the original suggestion by RényiRENYI . In fact defining $`S_R=X_s(S_T)`$ the (Tsallis) entropy is mapped to the extensive quantity:
$$S_R=\frac{1}{1q}\mathrm{ln}f^q.$$
(47)
### V.1 Canonical and Extended Equilibrium
Summarizing this part it seems to be useful to spell out the general conditions leading to the traditional and to the generalized treatment of canonical equilibrium. Three conditions are important:
1. In the equilibrium state the energy shell is covered uniformly (ergodic assumption),
2. The average energy exchange between the system under study and its environment is zero (not a driven system),
3. The fluctuations in the system’s energy are negligible compared to the average value of its energy (canonical limit).
The third requirement has its physical basis in the fact that the characteristic range of forces acting in the system is much less then the system size. So in the thermodynamical limit of large systems, $`VA\mathrm{}`$, with a finite interaction range, $`\mathrm{}`$, all surface effects (proportional to $`A`$) in a large volume, $`V`$, become negligible. In this case the entropy is quite generally additive, $`S_{12}=h(S_1,S_2)=S_1+S_2`$, and the equilibrium canonical distribution is of Boltzmann-Gibbs type, $`Zf=\mathrm{exp}(E/T)`$. Note that already the quantum statistical distributions violate this requirement: the exchange interaction (Pauli blocking or Bose enhancement) is long-ranged.
In an extended equilibrium state, non-extensive thermodynamics and generalized kinetic descriptions deal with, the third requirement is not fulfilled. Either long range forces, or large fluctuations (related to each other) keep the non-extensivity parameter $`aA\mathrm{}/V`$ finite. Applying this idea to the time direction, also long-term memory effects spoil the traditional picture of canonical equilibrium and textbook thermodynamics. The unusual large fluctuations are treated in the framework of superstatisticsSUPER-STAT , where the traditionally intensive thermodynamical parameters, like temperature, have a probability distribution instead of a fixed value. This view, of course, also transforms the Boltzmann-Gibbs distribution into an arbitrary one:
$$f_{eq}(E)=𝑑\beta 𝒫(\beta )\frac{1}{Z(\beta )}e^{\beta E}.$$
(48)
For example, the Tsallis distribution can be obtained by considering a Gamma-distributed inverse temperature, $`\beta =1/T`$. Some preview about a numerical simulation of NEBE for the Tsallis distribution can be obtained from a poster presented by G. Purcsel at this conference.
## VI Hadron Statistics
Statistical models have been often applied to hadron physics. Starting with Rolf Hagedorn’s statistical model of meson resonancesHAGEDORN , several attempts occured to describe hadron multiplicities in elementary collisions by means of statistical distributions. The very idea of a phase transition between confined and deconfined quark matter relies on traditional equilibrium thermodynamics. With the dawn of relativistic heavy ion physics a search for the nuclear and quark matter equation of state began assuming a local thermal equilibrium in an otherwise exploding fireball. Such models and their predictions for experimentally observable properties will be reviewed in other lecturesTOR-LECT ; CSER-LECT .
### VI.1 Particle spectra and equation of state
Here, following the above introduction to non-extensive thermodynamics, just some basic ideas about the interpretation of the power-law tails of nearly exponential spectra shall be discussed. Both exponential and power-law regions are observed in hadron spectra stemming from pp or heavy-ion collisions. The $`m_T`$-scaling property, which reveals in particular after taking into account Doppler blue shift like corrections due to a transverse flow, indicates a thermodynamical origin. Although the traditional approach explains the power-law tail at very high $`p_T`$ values by pQCD calculationspQCD , the non-extensive statistics provides a unified view for the whole spectrum. Certainly, jets also contribute to the hard part of light hadron spectra, but their angular distribution is far from uniform. In principle, these contributions may be filtered out. The question is, whether the rest still shows a (now statistical) power-law behavior.
Another point can be made by inspecting the minimum bias pion $`p_T`$-spectrum form RHIC AuAu collisions at 200 GeV (Fig.1 in Ref. SLIDING-SLOPE ). To this a non-exponential fit can already be made at the $`p_T`$-region between $`1`$ and $`4`$ GeV. The extrapolation of this fit almost coincides with the fit to the whole observed range between $`1`$ and $`12`$ GeV. So one concludes that the power-law behavior is not only a very hard scale physics. Furthermore in the non-extensive statistical approach there is a connection between the soft properties (temperature $`T`$) and the hard properties (power $`v`$, transition scale $`E_c`$) cf. eq.(26).
The dilemma between the statistical and pQCD explanation can be well expressed in the following formula for the particle yield:
$$\frac{(2\pi \mathrm{})^3}{V}\frac{d^3N}{dk^3}=\underset{0}{\overset{\mathrm{}}{}}𝑑\omega \omega \rho (\omega ,\stackrel{}{k})f(\omega /T).$$
(49)
The observed spectra, even not considering a collective flow or an extended source, are already convolutions of a spectral function, $`\rho (\omega ,\stackrel{}{k})`$, and statistical-enviromental effects, $`f(\omega /T)`$. Both these factors may embrance non-trivial effects, distorting this way a naively expected Boltzmann-Gibbs exponential. Both may contain a further energy scale; the spectral function $`\mathrm{\Lambda }_{QCD}`$, the thermodynamical weight $`E_c`$. Even in the case, when $`\rho `$ would contain a very sharp peak at a given dispersion relation, $`\omega _k=\omega (\stackrel{}{k})`$, only, the result is a composite function, $`\mathrm{exp}(\omega _k/T)`$, undistinguishable from a deformed thermodynamical equilibrium, $`\mathrm{exp}(X(|k|)/T)`$. Quite generally in this quasiparticle picture, whenever the effective dispersion relation, $`\omega (\stackrel{}{k})`$, contains parameters depending e.g. on the temperature, $`T`$, an assumption of an exponential equilibrium distribution like $`\mathrm{exp}(\omega (\stackrel{}{k},T)/T)`$ questions the physical meaning of the temperature. If it is an enviromental parameter, the inverse slope is not necessarily equal to this, if it is an inverse slope, it is not necessarily the enviromental parameter governing an effective dispersion relation.
Finally we note, that non-exponential distributions can also be interpreted in the framework of so called superstatisticsSUPER-STAT . Here a distribution of the intensive thermodynamical parameters is assumed instead of a fixed value. The physical reason behind maybe the finiteness of the observed subsystem, as well as extraordinarily large fluctuations not scaling down with the system size alike the ones governed by the Gaussian distribution would do. The Tsallis distribution, for example, can be given by a continous distribution of the inverse temperature according to a Gamma distribution:
$$(1+x/c)^{(c+1)}=\frac{1}{\mathrm{\Gamma }(c+1)}\underset{0}{\overset{\mathrm{}}{}}𝑑tt^ce^te^{xt/c}.$$
(50)
There are several possibilities to intrepret such a distribution. Among others a heat conduction equation with multiplicative noiseHEAT-COND or taking into account an energy imbalance in two-body collisions due to a presence of further participant agentsFLUCT-ENER leads to the required result.
### VI.2 Limiting temperature with Tsallis distribution
Another probe for the Tsallis distribution than the one-particle $`p_T`$-spectra may be given by considering the spectrum of heavy hadronic resonances. An exponentially growing mass spectrum, originally proposed by Hagedorn and recently checked again latest experimental data in Ref.RES-SPECT , with its famous consequence of having a limiting (or Hagedorn-) temperature for such a system, can be reconstructed on the basis of Tsallis distributed quark constituents. This approachANDRE assumes that the Tsallis distribution of the quarks and antiquarks is folded into mesonic and baryonic distributions of the conserved total energy satisfying $`X(E)=_iX(E_i)`$. In general for an $`N`$-fold convolution of massless constituents with $`d`$ dimensional momenta one easily obtains that the average energy satisfies
$$X(E)=N\underset{j=1}{\overset{d}{}}\frac{TE_c}{E_cjT}.$$
(51)
This expression diverges first as the temperature reaches the limiting value, $`T_H=E_c/d`$ starting from zero (for positive $`E_c=1/a`$, i.e. for repulsive modifications of the extensive energy composition rule). This way Hagedorn hadrons emerge from Tsallis partons. This model also explains naturally, why the baryonic and mesonic mass spectrum seems to have a different rise: they contain different polynomial coefficients in front of three exponential factors.
## VII Conclusion
In conclusion we have reviewed basic ideas of statistical physics and kinetic theory which may lead to a non-extensive thermodynamics. In particular the role of a multiplicative noise in the linear Fokker-Planck and Langevin approach has been emphasized. Generalizations of the Boltzmann equation towards non-factorizing yield factors or non-additive energy composition rules were then shown to lead again to non-extensive entropy definitions and thermodynamics. Finally these ideas have been related to hadron spectra observed in relativistic heavy ion collisions. We pointed out that not only the omnipresence and $`m_T`$-scaling of power-law tails in particle spectra, but also a possible interpretation of the Hagedorn spectrum of heavy resonances may support the presence of a non-exponential equilibrium distribution in hot quark matter.
Acknowledgment
Enlightening discussions and common work with Drs. Géza Györgyi at Eötvös University and Antal Jakovác at the Technical University Budapest are hereby gratefully acknowledged. Further collaborations with Berndt Müller, André Peshier and Gábor Purcsel contributed essentially to the development of ideas presented above. This work has been supported by the Hungarian National Science Fund OTKA (T037689) and the Deutsche Forschungsgemeinschaft. |
warning/0507/hep-th0507282.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The recent past has seen a new interest on exactly marginal deformations of $`𝒩=4`$ SYM theory preserving $`𝒩=1`$ supersymmetry , in particular after the supergravity duals of the so–called $`\beta `$-deformations of $`𝒩=4`$ $`SU(N)`$ SYM theory have been found by Lunin and Maldacena in .
Now new results start to emerge also from the field theory side: in various properties of composite operators of the deformed theory have been investigated at the perturbative level (see also ). The outcome is that the deformed theory shares some properties with the undeformed $`𝒩=4`$ theory, but new features emerge, as for example the finite corrections of the two- and three-point functions of protected operators .
The chiral ring of the theory was identified in for generic values of the deformation parameter. It is given by the operators $`\mathrm{Tr}(\mathrm{\Phi }_i^J)`$, $`i=1,2,3`$, and $`\mathrm{Tr}(\mathrm{\Phi }_1^J\mathrm{\Phi }_2^J\mathrm{\Phi }_3^J)`$. In it was shown that also the operator $`\mathrm{Tr}(\mathrm{\Phi }_1\mathrm{\Phi }_2)`$ does not acquire anomalous dimension.
In this paper we will focus on the operators $`\mathrm{Tr}(\mathrm{\Phi }_i^J\mathrm{\Phi }_j)`$, $`ij`$. As opposed to what happens in the undeformed $`𝒩=4`$ case, these operators are not protected and their anomalous dimension was computed at one-loop in . Our interest in these operators is motivated by the fact that in the large $`J`$ limit they resemble the BMN operators of the $`𝒩=4`$ theory . Indeed a perturbative superfield analysis performed at low orders shows that the supergraphs contributing to their anomalous dimension are exactly the same as the ones of the BMN case. Then we apply the derivation of to this class of operators and compute their exact anomalous dimension in the planar limit.
The consistency between the perturbative supergraph approach and the result obtained by using the method of suggests that the one-loop superconformal invariance condition remains valid in the planar limit to all orders of perturbation theory (at least for real values of the $`\beta `$ parameter of the $`\beta `$-deformation which is the case we consider here). We confirm this result by exploiting the formal analogy between $`\beta `$-deformed and non-commutative field theories.
The paper is organized as follows. In Section 2 we present the setup for the perturbative calculation that we perform in Section 3, where we compute the anomalous dimensions of the operators $`\mathrm{Tr}(\mathrm{\Phi }_i^J\mathrm{\Phi }_j)`$, $`ij`$ in the large $`N`$ limit up to order $`g^4`$. The contributing graphs are the same that one would encounter in the calculation of the anomalous dimension of BMN operators. In Section 4 we apply to our operators the procedure introduced in and give an all-order result for their anomalous dimensions in the large $`J`$ and large $`N`$ limit. Then in Section 5 we prove that the one-loop condition of superconformal invariance remains valid to all orders in the planar limit.
## 2 Generalities
We consider the following deformation of the $`𝒩=4`$ SYM theory
$`S[j,\overline{j}]`$ $`=`$ $`{\displaystyle d^8z\mathrm{Tr}\left(e^{gV}\overline{\mathrm{\Phi }}_ie^{gV}\mathrm{\Phi }^i\right)}+{\displaystyle \frac{1}{2g^2}}{\displaystyle d^6z\mathrm{Tr}W^\alpha W_\alpha }`$ (2.1)
$`+ih{\displaystyle d^6z\mathrm{Tr}(q\mathrm{\Phi }_1\mathrm{\Phi }_2\mathrm{\Phi }_3\frac{1}{q}\mathrm{\Phi }_1\mathrm{\Phi }_3\mathrm{\Phi }_2)}i\overline{h}{\displaystyle d^6\overline{z}\mathrm{Tr}(\overline{q}\overline{\mathrm{\Phi }}_1\overline{\mathrm{\Phi }}_3\overline{\mathrm{\Phi }}_2\frac{1}{\overline{q}}\overline{\mathrm{\Phi }}_1\overline{\mathrm{\Phi }}_2\overline{\mathrm{\Phi }}_3)}`$
$`+{\displaystyle d^6zj𝒪}+{\displaystyle d^6\overline{z}\overline{j}\overline{𝒪}}`$
where we have set $`qe^{i\pi \beta }`$ and in the following we choose $`\beta `$ real so that $`q\overline{q}=1`$. We have added to the classical action source terms for composite chiral operators generically denoted by $`𝒪`$ with $`j`$ ($`\overline{j}`$) (anti)chiral sources. The superfield strength $`W_\alpha =i\overline{D}^2(e^{gV}D_\alpha e^{gV})`$ is given in terms of a real prepotential $`V`$ and $`\mathrm{\Phi }_{1,2,3}`$ contain the six scalars of the original $`𝒩=4`$ SYM theory organized into the $`\mathrm{𝟑}\times \overline{\mathrm{𝟑}}`$ of $`SU(3)SU(4)`$. We write $`V=V^aT_a`$, $`\mathrm{\Phi }_i=\mathrm{\Phi }_i^aT_a`$ where $`T_a`$ are $`SU(N)`$ matrices in the fundamental representation. In the following we use the notation $`\mathrm{Tr}(T^aT^bT^c\mathrm{})(abc\mathrm{})`$.
The $`\beta `$–deformation breaks the original $`SU(4)`$ R–symmetry to $`U(1)_R`$. However, besides the $`Z_3`$ symmetry associated to cyclic permutations of $`(\mathrm{\Phi }_1,\mathrm{\Phi }_2,\mathrm{\Phi }_3)`$, two extra non–R–symmetry $`U(1)`$’s survive. Applying the $`a`$–maximization procedure and the conditions of vanishing ABJ anomalies it turns out that $`U(1)_R`$ is the one which assigns the same R–charge $`\omega `$ to the three elementary superfields, whereas the charges respect to the two non–R–symmetries $`U(1)_1\times U(1)_2`$ can be chosen to be $`(\mathrm{\Phi }_1,\mathrm{\Phi }_2,\mathrm{\Phi }_3)(0,1,1)`$ and $`(1,1,0)`$, respectively.
As discussed in , at the quantum level the theory is superconformal invariant (and then finite) up to two loops if the coupling constants satisfy the following condition (vanishing of the beta functions)
$$h\overline{h}\left[1\frac{1}{N^2}\left|q\frac{1}{q}\right|^2\right]g^2=0$$
(2.2)
In the large $`N`$ limit this condition reduces simply to $`g^2=h\overline{h}`$, independently of the value of $`q`$. In the three loop correction to (2.2) has been also evaluated. Since it turns out to be suppressed for $`N\mathrm{}`$, the condition $`g^2=h\overline{h}`$ is the correct condition for superconformal invariance in the planar limit up to three loops.
At the superconformal fixed point of the theory we compute the anomalous dimensions for the class of non–protected operators <sup>1</sup><sup>1</sup>1The choice of $`\mathrm{\Phi }_1`$ and $`\mathrm{\Phi }_2`$ superfields is totally arbitrary and we expect the operators $`\mathrm{Tr}(\mathrm{\Phi }_i^J\mathrm{\Phi }_k)`$, for any $`i,k`$ with $`ik`$ to have similar quantum properties. We will comment on this point later on.
$$𝒪_J=\mathrm{Tr}(\mathrm{\Phi }_1^J\mathrm{\Phi }_2)$$
(2.3)
They are charged under $`U(1)_1\times U(1)_2`$ with charges $`(1,1J)`$.
Using the equations of motion from the action (2.1) with $`j=\overline{j}=0`$ (from now on we neglect factors of $`e^{\pm gV}`$ since they are not relevant to our purposes)
$$\overline{D}^2\overline{\mathrm{\Phi }}_3^a=ih\mathrm{\Phi }_1^b\mathrm{\Phi }_2^c[q(abc)\frac{1}{q}(acb)]$$
(2.4)
it is easy to see that
$$𝒪_J=\frac{i}{h[q\frac{1}{q}]}\overline{D}^2\mathrm{Tr}(\mathrm{\Phi }_1^{J1}\overline{\mathrm{\Phi }}_3)+\frac{1}{N}\mathrm{Tr}(\mathrm{\Phi }_1^{J1})\mathrm{Tr}(\mathrm{\Phi }_1\mathrm{\Phi }_2)$$
(2.5)
As long as $`J>1`$, in the large $`N`$ limit the operator $`𝒪_J`$ becomes descendent of the primary $`\mathrm{Tr}(\mathrm{\Phi }_1^{J1}\overline{\mathrm{\Phi }}_3)`$, whereas for finite $`N`$ the combination $`𝒪_J\frac{1}{N}\mathrm{Tr}(\mathrm{\Phi }_1^{J1})\mathrm{Tr}(\mathrm{\Phi }_1\mathrm{\Phi }_2)`$ is descendent. The exceptional case $`J=1`$ corresponds to the chiral primary operator whose protection has been proven perturbatively in .
In the next Sections we will concentrate on the evaluation of the anomalous dimensions for the $`𝒪_J`$ operators.
## 3 The perburbative calculation
In this Section we compute the anomalous dimension of $`𝒪_J`$ in (2.3) perturbatively, up to two loops. For generic values of $`J`$ we perform the calculation in the large $`N`$ limit in order to avoid dealing with mixing with multitrace operators.
In order to compute anomalous dimensions we evaluate one–point correlators $`𝒪_Je^{S_{int}}`$ where $`S_{int}`$ is the sum of the interaction terms in (2.1). Divergent contributions proportional to the operator itself are removed by a multiplicative renormalization which in dimensional regularization reads
$$𝒪_J^{(bare)}𝒪_J\left(1+\underset{k=0}{\overset{\mathrm{}}{}}\frac{a_k(\lambda ,q,N)}{ϵ^k}\right)Z𝒪_J$$
(3.1)
where we have introduced the ’t Hooft coupling $`\lambda =\frac{g^2N}{4\pi ^2}`$. There is no dependence on the $`h`$ coupling since we are at the superconformal point where $`h`$ can be expressed in terms of the other couplings through the condition of vanishing beta functions.
The anomalous dimension is then given by
$$\gamma 2\lambda \frac{da_1(\lambda ,q,N)}{d\lambda }$$
(3.2)
Therefore, at any order it is easily read from the simple pole divergence.
We perform perturbative calculations in a superspace setup by following closely the procedure used in (we refer to those papers for conventions and technical details). After D–algebra supergraphs are reduced to ordinary Feynman diagrams which we evaluate in momentum space. We work in dimensional regularization, $`d=42ϵ`$, and minimal subtraction scheme.
From the action (2.1) quantized in the Feynman gauge we read the superfield propagators ($`z(x,\theta ,\overline{\theta })`$)
$`V^a(z_1)V^b(z_2)=\delta ^{ab}{\displaystyle \frac{1}{(x_1x_2)^2}}\delta ^{(4)}(\theta _1\theta _2)`$
$`\mathrm{\Phi }_i^a(z_1)\overline{\mathrm{\Phi }}_j^b(z_2)=\delta _{ij}\delta ^{ab}{\displaystyle \frac{1}{(x_1x_2)^2}}\delta ^{(4)}(\theta _1\theta _2)`$ (3.3)
and the three–point vertices
$`(\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi })_{\mathrm{vertex}}ih\mathrm{\Phi }_1^a\mathrm{\Phi }_2^b\mathrm{\Phi }_3^c[q(abc){\displaystyle \frac{1}{q}}(acb)]`$
$`(\overline{\mathrm{\Phi }}\overline{\mathrm{\Phi }}\overline{\mathrm{\Phi }})_{\mathrm{vertex}}i\overline{h}\overline{\mathrm{\Phi }}_1^a\overline{\mathrm{\Phi }}_2^b\overline{\mathrm{\Phi }}_3^c\left[\overline{q}(acb){\displaystyle \frac{1}{\overline{q}}}(abc)\right]`$
$`(\overline{\mathrm{\Phi }}V\mathrm{\Phi })_{\mathrm{vertex}}g\overline{\mathrm{\Phi }}_i^aV^b\mathrm{\Phi }_i^c[(abc)(acb)]`$ (3.4)
At the lowest order the only contribution to the one–point function for the operator $`𝒪_J`$ is the one given in Fig. 1 where, using the notation introduced in , the horizontal bold line indicates the spacetime point where the operator is inserted.
Figure 1: One–loop contribution to the $`𝒪_J`$ anomalous dimension
The corresponding contribution is proportional to the self–energy integral
$$I_1d^nk\frac{1}{k^2(pk)^2}\frac{1}{(4\pi )^2}\frac{1}{ϵ}$$
(3.5)
Evaluating the color factor, the combinatorics and taking into account a minus sign from D–algebra we obtain
$$\mathrm{Diagram}1\frac{1}{ϵ}|q\frac{1}{q}|^2\frac{|h|^2N}{(4\pi )^2}$$
(3.6)
Using the one–loop superconformal condition in the planar limit ($`g^2=|h|^2`$) and the definition (3.2) we immediately find the one–loop anomalous dimension
$$\gamma ^{(1)}=\frac{1}{2}\left|q\frac{1}{q}\right|^2\lambda $$
(3.7)
At two loops (order $`\lambda ^2`$) the diagrammatic contributions are drawn in Fig. 2.
Figure 2: Two–loop contributions to the $`𝒪_J`$ anomalous dimension
Performing the D–algebra we reduce all the diagrams to ordinary Feynman diagrams containing the loop structure as in Fig. 3.
Figure 3: The two–loop Feynman diagram
The associated momentum integral is
$$I_2d^nk_1d^nk_2\frac{1}{k_1^2(p_1k_1k_2)^2k_2^2(p_1+p_2k_2)^2}$$
(3.8)
As long as we are only concerned with UV divergences we can safely set one of the external momenta to zero. Thus the graph is easily evaluated being proportional to two nested self–energies. We obtain (in the G-scheme)
$$I_2\frac{1}{(4\pi )^4}\frac{1}{2ϵ^2}(1+5ϵ)\frac{1}{(p^2)^{2ϵ}}$$
(3.9)
where we have kept only divergent terms. Performing the subtraction of the subdivergence we finally have
$$\left[I_2\right]_{sub}\frac{1}{(4\pi )^4}\left[\frac{1}{2ϵ^2}+\frac{1}{2ϵ}\right]$$
(3.10)
Computing the combinatorics, the color factors and taking into account minus signs from the vector propagator we find that the factors in front of (3.10) for the various diagrams are
$`(2a)2(q{\displaystyle \frac{1}{q}})(\overline{q}{\displaystyle \frac{1}{\overline{q}}})g^2|h|^2N^2`$
$`(2b)2(q{\displaystyle \frac{1}{q}})(\overline{q}{\displaystyle \frac{1}{\overline{q}}})g^2|h|^2N^2`$
$`(2c)2(q{\displaystyle \frac{1}{q}})(\overline{q}{\displaystyle \frac{1}{\overline{q}}})g^2|h|^2N^2`$
$`(2d)(q{\displaystyle \frac{1}{q}})(\overline{q}{\displaystyle \frac{1}{\overline{q}}})\left({\displaystyle \frac{q}{\overline{q}}}+{\displaystyle \frac{\overline{q}}{q}}\right)|h|^4N^2`$ (3.11)
Summing all the contributions, using the planar superconformal condition $`|h|^2=g^2`$ and the definition (3.2), we find
$$\gamma ^{(2)}=\frac{1}{8}\left|q\frac{1}{q}\right|^4\lambda ^2$$
(3.12)
We observe that the diagrams contributing to the anomalous dimensions for our operators are exactly the same as the ones for BMN operators in $`𝒩=4`$ SYM in the planar limit . In fact, up to this order the calculation is exactly the same under the formal identification $`|q\frac{1}{q}|^2(e^{i\varphi }+e^{i\varphi }2)`$, where $`\varphi `$ is the phase of BMN operators . We expect that the same pattern will persist at any order in perturbation theory. In particular, as in the BMN case, the graphs relevant for the calculation are only the ones where the interactions are close to the “impurity” $`\mathrm{\Phi }_2`$: at $`L`$–loop order the interactions may involve at most the $`\mathrm{\Phi }_1`$ lines which are $`L`$–steps far away from the impurity. As an important consequence, in the large $`J`$ limit the anomalous dimensions do not grow with $`J`$.
To close this Section we note that the result we have found for the anomalous dimensions of the operators $`\mathrm{Tr}(\mathrm{\Phi }_1^J\mathrm{\Phi }_2)`$ at large $`N`$ are actually valid for any operator of the form $`\mathrm{Tr}(\mathrm{\Phi }_i^J\mathrm{\Phi }_k)`$ with $`ik`$. In fact the superpotential is invariant under cyclic permutation of $`(\mathrm{\Phi }_1,\mathrm{\Phi }_2,\mathrm{\Phi }_3)`$, and in addition it becomes invariant if non–cyclic exchanges of fields are accompanied by
$$q\frac{1}{q}$$
(3.13)
Since the anomalous dimensions are proportional to powers of the effective coupling $`\alpha \lambda \left|q\frac{1}{q}\right|^2`$ which is invariant under (3.13) we conclude that the result is valid for any operator of the form $`\mathrm{Tr}(\mathrm{\Phi }_i^J\mathrm{\Phi }_k)`$, $`ik`$.
## 4 The exact anomalous dimensions
Motivated by the formal correspondence of the previous calculation with the BMN case, in this Section we are going to compute the exact anomalous dimensions in the large $`N`$, large $`J`$ limit by using the procedure introduced in for BMN operators. In the context of $`\beta `$–deformed theories the same procedure has been applied to the class of BMN operators .
We concentrate on the operator $`𝒪_{J+1}`$ which, as follows from eq. (2.5), in the planar limit satisfies
$$\overline{D}^2𝒰_J=ih[q\frac{1}{q}]𝒪_{J+1}$$
(4.1)
where we have defined
$$𝒰_J\mathrm{Tr}(\mathrm{\Phi }_1^J\overline{\mathrm{\Phi }}_3)$$
(4.2)
As already noticed, this shows that the $`𝒪_{J+1}`$ operators are descendants of the $`𝒰_J`$ ones. Being part of the same superconformal multiplet they share the renormalization properties, i.e. they will have the same scaling dimension and the same perturbative corrections to their overall normalization. Moreover since $`𝒰_J`$ is not a Konishi-like operator it is not affected by the Konishi anomaly.
As discussed in details in , in any $`N=1`$ superconformal field theory the two–point function for a primary operator $`𝒜_{s,\overline{s}}`$ is fixed and given by ($`z(x,\theta ,\overline{\theta })`$)
$`<𝒜_{(s,\overline{s})}(z)\overline{𝒜}_{(s,\overline{s})}(z^{})>=f_𝒜(g^2,N,h,\overline{h})\{{\displaystyle \frac{1}{2}}D^\alpha \overline{D}^2D_\alpha +{\displaystyle \frac{w}{4(\mathrm{\Delta }_0+\gamma )}}[D^\alpha ,\overline{D}^{\dot{\alpha }}]i_{\alpha \dot{\alpha }}`$
$`+{\displaystyle \frac{(\mathrm{\Delta }_0+\gamma )^2+w^22(\mathrm{\Delta }_0+\gamma )}{4(\mathrm{\Delta }_0+\gamma )(\mathrm{\Delta }_0+\gamma 1)}}\mathrm{}\}{\displaystyle \frac{\delta ^4(\theta \theta ^{})}{|xx^{}|^{2(\mathrm{\Delta }_0+\gamma )}}}`$ (4.3)
where $`\mathrm{\Delta }_0=s+\overline{s}`$ is the tree–level dimension of the operator, $`\omega =s\overline{s}`$ is its R–symmetry charge <sup>2</sup><sup>2</sup>2We assume $`\omega `$ not to renormalize. In fact, once the R–symmetry of the elementary fields is fixed by requiring the exact R–symmetry of the superpotential, any composite operator has a fixed charge given by the sum of the charges of its elementary constituents. and $`\gamma `$ is the exact anomalous dimension.
The relation (4.3) can be straightforwardly applied to our primary operators $`𝒰_J`$. The analysis of the two-point correlator for the $`𝒪_J`$’s is somewhat subtler since, as we see from eq. (4.1), these chiral operators are not primaries and in principle the relation (4.3) cannot be applied to their correlators. However, as we are going to show, in the large $`J`$ limit these operators turn out to behave as CPO’s and (4.3) can be safely used.
To this end we remind that in general given a chiral operator $`𝒜`$, the condition for the operator to be non–protected (anomalous dimension acquired) is equivalent to the condition that its chiral nature is not maintained under superconformal transformations, i.e. $`\overline{D}(\delta _{\overline{S}}𝒜)\{\overline{S},\overline{D}\}𝒜0`$ (see for instance ). In fact, writing schematically the superconformal algebra relation for a scalar operator as $`\{\overline{S},\overline{D}\}=\mathrm{\Delta }\omega `$, we have
$$\{\overline{S},\overline{D}\}𝒜=(\mathrm{\Delta }\omega )𝒜=\left[(\mathrm{\Delta }_0+\gamma )\omega \right]𝒜=\gamma 𝒜$$
(4.4)
where we have used $`\mathrm{\Delta }=\mathrm{\Delta }_0+\gamma `$ and for a chiral operator $`\omega =\mathrm{\Delta }_0`$. Therefore if $`\gamma 0`$, $`\overline{S}𝒜`$ is not chiral anymore. Viceversa, if $`\{\overline{S},\overline{D}\}𝒜=0`$, then $`\overline{S}𝒜`$ is still chiral and the dimension is protected by the well–known condition $`\mathrm{\Delta }=\omega `$.
An alternative proof goes through the simple observation that the conditions
$$s+\overline{s}=\mathrm{\Delta }_0+\gamma s\overline{s}=\mathrm{\Delta }_0$$
(4.5)
imply
$$s=\mathrm{\Delta }_0+\frac{\gamma }{2}\overline{s}=\frac{\gamma }{2}$$
(4.6)
The appearance of $`\overline{s}0`$ signals the lack of chirality of the quantum operator.
We now apply the previous argument to our operators $`O_J`$ to prove that in the large $`N`$, large $`J`$ limit the violation of chirality is suppressed and they behave as CPO’s. In the limit of large R–symmetry $`\omega =J`$ it is more natural to consider
$$\frac{1}{J}\{\overline{S},\overline{D}\}𝒪_J=\left(\frac{\mathrm{\Delta }_0+\gamma }{J}1\right)𝒪_J=\frac{\gamma }{J}𝒪_J$$
(4.7)
As discussed in the previous Section, at any fixed order in perturbation theory the anomalous dimension $`\gamma `$ does not grow with $`J`$. It follows that in the large $`J`$ limit the r.h.s. of eq. (4.7) is suppressed and the operator behaves as a chiral primary. In particular, in this limit it is consistent to apply eq. (4.3) for the evaluation of its two–point function.
Supported by these considerations we can now proceed exactly as in and find
$`<\overline{D}^2𝒰_J(z)D^2\overline{𝒰}_J(z^{})>=`$
$`={\displaystyle \frac{N^{J+1}}{(4\pi ^2)^{J+1}}}f\overline{D}^2\{{\displaystyle \frac{1}{2}}D^\alpha \overline{D}^2D_\alpha +{\displaystyle \frac{J1}{4(J+1+\gamma )}}[D^\alpha ,\overline{D}^{\dot{\alpha }}]i_{\alpha \dot{\alpha }}`$
$`+{\displaystyle \frac{(J+1+\gamma )^2+(J1)^22(J+1+\gamma )}{4(J+1+\gamma )(J+\gamma )}}\mathrm{}\}D^2{\displaystyle \frac{\delta ^4(\theta \theta ^{})}{|xx^{}|^{2(J+1+\gamma )}}}`$
$`={\displaystyle \frac{N^{J+1}}{(4\pi ^2)^{J+1}}}f(\gamma ^2+2\gamma )\overline{D}^2D^2{\displaystyle \frac{\delta ^4(\theta \theta ^{})}{|xx^{}|^{2(J+2+\gamma )}}}`$ (4.8)
and
$$<𝒪_{J+1}(z)\overline{𝒪}_{J+1}(z^{})>=\frac{N^{J+2}}{(4\pi ^2)^{J+2}}f\overline{D}^2D^2\frac{\delta ^4(\theta \theta ^{})}{|xx^{}|^{2(J+2+\gamma )}}$$
(4.9)
where $`f`$ is the common normalization function not fixed by superconformal invariance.
From the relation (4.1) the two correlators are related by
$$\overline{D}^2𝒰_J(z)D^2\overline{𝒰}_J(z^{})=|h|^2\left|q\frac{1}{q}\right|^2𝒪_{J+1}(z)\overline{𝒪}_{J+1}(z^{})$$
(4.10)
Therefore, inserting in (4.10) the explicit expressions (4.8) and (4.9) we end up with an algebraic equation
$$\gamma ^2+2\gamma =|h|^2\left|q\frac{1}{q}\right|^2\frac{N}{4\pi ^2}$$
(4.11)
which allows to find the exact expression for the anomalous dimensions
$`\gamma `$ $`=`$ $`1+\sqrt{1+|h|^2\left|q{\displaystyle \frac{1}{q}}\right|^2{\displaystyle \frac{N}{4\pi ^2}}}`$ (4.12)
$`=`$ $`{\displaystyle \frac{1}{2}}|h|^2\left|q{\displaystyle \frac{1}{q}}\right|^2{\displaystyle \frac{N}{4\pi ^2}}{\displaystyle \frac{1}{8}}|h|^4\left|q{\displaystyle \frac{1}{q}}\right|^4{\displaystyle \frac{N^2}{(4\pi ^2)^2}}+\mathrm{}`$
Up to the second order this expression coincides with the perturbative results obtained in the previous Section.
We note that our operators $`𝒪_J`$ can be thought as dual to the 0–modes of the BMN sector considered in . Formula (4.12) is in agreement with the results presented in those papers for the spectrum of the 0–modes.
## 5 The superconformal condition at large $`N`$
In the previous Section, exploiting the superconformal invariance of the theory and its equations of motion we have shown that the exact anomalous dimension for the $`𝒪_J`$ operator can be written as
$$\gamma =1+\sqrt{1+2\gamma ^{(1)}}$$
(5.1)
where $`\gamma ^{(1)}`$ is the one–loop anomalous dimension. A direct calculation provides an expression for $`\gamma ^{(1)}`$ proportional to $`|h|^2`$. As discussed in Section 3, at this order we are allowed to use the planar superconformal condition $`|h|^2=g^2`$ to re-express $`\gamma ^{(1)}`$ in terms of $`g^2`$ only (see eq. (3.7)). Now if in (5.1) we use $`\gamma ^{(1)}`$ given in terms of $`g^2`$ and expand the square root, we obtain a perturbative formula for $`\gamma `$ which agrees with the actual perturbative calculation only if the condition $`|h|^2=g^2`$ is valid at any order.
Motivated by this observation we are led to conjecture that in the large $`N`$ limit the condition $`|h|^2=g^2`$ is indeed the correct condition for superconformal invariance at any order in perturbation theory. Direct confirmations of this conjecture can be found in the literature up to order $`g^6`$ . Now we give an argument to prove that this is true to all orders.
We remind that in $`𝒩=1`$ supersymmetric theories the superconformal invariance condition (i.e. vanishing of beta functions) can be expressed as the vanishing of the anomalous dimensions of the elementary superfields . Therefore, in order to study superconformal invariance, it is sufficient to focus on the divergent corrections to the propagators of the elementary fields.
In the $`\beta `$–deformed theory we consider a generic $`L`$–loop diagram contributing to the propagator of the $`\mathrm{\Phi }_i`$ superfield. The crucial observation is the following: If we prove that at the planar level, as long as $`q\overline{q}=1`$, this diagram does not depend on $`q`$, then we are sure that $`|h|^2=g^2`$ is the exact solution of the superconformal invariance equations. In fact, if it is independent of $`q`$, the corresponding perturbative contribution is the same for any deformed theory, independently of the choice of the $`q`$–deformation. In particular, it is the same for any deformed theory ($`q1`$) and for the underformed one ($`q=1`$). Focusing on the undeformed case we can conclude that $`|h|^2=g^2`$ is the exact condition for the planar superconformal invariance, since $`q=1`$ and $`|h|^2=g^2`$ bring us back to the $`𝒩=4`$ case which is known to be exactly superconformal. The independence of the perturbative corrections on $`q`$ allows to extend this statement to any deformed theory.
To conclude the proof we need to show that the contribution from a generic self–energy planar diagram never depends on $`q`$. We can focus on diagrams containing only matter vertices because adding vector propagators cannot introduce any $`q`$-dependence. We exploit the formal analogy between the deformed theory and noncommutative (nc) field theory. As observed in the deformed potential can be written as
$$ihd^6z\mathrm{Tr}(\mathrm{\Phi }_1[\mathrm{\Phi }_2,\mathrm{\Phi }_3]_{})+\mathrm{h}.\mathrm{c}.$$
(5.2)
where
$$fg=e^{i\pi \beta Q_f^{(i)}M_{ij}Q_g^{(j)}}fg,$$
(5.3)
$`Q^{(i)}`$, $`i=1,2`$ being the non–R–symmetry $`U(1)_1\times U(1)_2`$ charges and $`M`$ the antisymmetric matrix $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$. When drawing a Feynman diagram we can consider the flow of the charges inside the diagram. Observing that the charges are conserved at any vertex and propagate through the straight lines we can formally identify them with the ordinary momenta in noncommutative diagrams. A known property of planar diagrams in nc field theory is that the star product phase factors dependent on the loop momenta cancel out (for a proof see ) and only an overall phase depending on the external momenta survives. In our case, exploiting the formal identification of charges with momenta, we can use the same arguments to conclude that any planar diagram will have a phase factor from (5.3) depending only on the configuration of the external charges. In the particular case of self–energy diagrams the overall phase is zero since $`\mathrm{\Phi }_i`$ and $`\overline{\mathrm{\Phi }}_i`$ have equal (but opposite) charges. In other words, any self–energy planar diagram always contains an equal number of $`q=\frac{1}{\overline{q}}`$ and $`\overline{q}=\frac{1}{q}`$ vertices. This concludes the proof of the $`q`$ independence of perturbative self–energy corrections.
We state that in the $`N\mathrm{}`$ limit the exact condition for superconformal invariance is simply $`|h|^2=g^2`$. Therefore, the theory described by the action
$`S={\displaystyle d^8z\mathrm{Tr}\left(e^{gV}\overline{\mathrm{\Phi }}_ie^{gV}\mathrm{\Phi }^i\right)}+{\displaystyle \frac{1}{2g^2}}{\displaystyle d^6z\mathrm{Tr}W^\alpha W_\alpha }`$
$`+ig{\displaystyle d^6z\mathrm{Tr}(e^{i\pi \beta }\mathrm{\Phi }_1\mathrm{\Phi }_2\mathrm{\Phi }_3e^{i\pi \beta }\mathrm{\Phi }_1\mathrm{\Phi }_3\mathrm{\Phi }_2)}+ig{\displaystyle d^6\overline{z}\mathrm{Tr}(e^{i\pi \beta }\overline{\mathrm{\Phi }}_1\overline{\mathrm{\Phi }}_2\overline{\mathrm{\Phi }}_3e^{i\pi \beta }\overline{\mathrm{\Phi }}_1\overline{\mathrm{\Phi }}_3\overline{\mathrm{\Phi }}_2)}`$
(5.4)
represents a $`𝒩=1`$ superconformal invariant theory for any value of $`\beta `$ real.
## 6 Conclusions
For the $`SU(N)`$, $`\beta `$–deformed $`𝒩=4`$ SYM theory we have considered the particular class of operators $`𝒪_J=\mathrm{Tr}(\mathrm{\Phi }_1^J\mathrm{\Phi }_2)`$. We have computed perturbatively their anomalous dimensions up to two loops. The calculation has been performed in the large $`N`$ limit in order to avoid mixing with multi-trace operators. Exploiting the techniques introduced in for the BMN operators we have evaluated their exact anomalous dimensions in the large $`N`$, large $`J`$ limit. In this limit the exact expression for the anomalous dimension depends on the deformation parameter $`q`$ only through the combination $`|q\frac{1}{q}|^2`$ and it is then invariant under $`q\frac{1}{q}`$. Observing that in $`\beta `$–deformed theory exchanging $`q\frac{1}{q}`$ amounts to exchange $`\mathrm{\Phi }_i\mathrm{\Phi }_j`$, for any $`i=1,2,3`$, $`ij`$, we may conclude that in the $`𝒪_J`$ sector, in the large $`N`$, large $`J`$ limit there is an enhancement of the $`SU(3)`$ symmetry and all the operators of the form $`\mathrm{Tr}(\mathrm{\Phi }_i^J\mathrm{\Phi }_k)`$, $`ik`$ renormalize in the same way.
A comparison between the exact result and the perturbative calculation suggests that the condition $`|h|^2=g^2`$, which up to three–loops guarantees the superconformal invariance of the theory in the planar limit, is actually sufficient for the exact invariance for $`N\mathrm{}`$. Indeed, we have given a direct proof at any order in perturbation theory. The main result of our paper is that the action in (5.4) is superconformal invariant at the quantum level without additional conditions on the couplings. In the context of the AdS/CFT correspondence this is the theory whose strong coupling phase is described by the supergravity dual found in . The $`𝒪_J`$ sector of this theory for large $`J`$ shares many similarities with the BMN sector of the $`𝒩=4`$ theory in the pp–wave limit. This opens the possibility for these operators to be dual to superstring states in some particular sector of the theory.
It is interesting to consider the extension of our calculations to the case of $`\beta `$ complex ($`q\overline{q}1`$). In this case the condition for superconformal invariance up to three loops, in the planar limit becomes
$$\frac{1}{2}|h|^2\left(q\overline{q}+\frac{1}{q\overline{q}}\right)=g^2$$
(6.1)
When (6.1) holds it is easy to see that the perturbative anomalous dimension still coincides with the expansion of the exact result (4.12). As for the case of $`\beta `$ real consistency of the perturbative calculation with the exact result would suggest that the condition (6.1) should be valid at any order. However, the proof we have presented in Section 5 makes repeated use of the requirement $`q\overline{q}=1`$ and cannot be immediately extended to the more general case. A different procedure should be found to prove or disprove that the one–loop condition (6.1) is sufficient to insure the exact conformal invariance even in the case of $`\beta `$ complex.
Generalizations of the present results to other deformed theories are presently under investigation and will be reported in .
## Acknowledgements
This work has been supported in part by INFN, PRIN prot. 2003023852\_008 and the European Commission RTN program MRTN–CT–2004–005104. |
warning/0507/quant-ph0507008.html | ar5iv | text | # Explicit Spin Coordinates
## 1 The Coordinate Representation of Spin
The nature of the electron remains an enigma ; however (and notwithstanding the Pauli and Dirac representations of spin ) the closest physical model is a spherically symmetric distribution of charge rotating around an axis \[3, p.55\]. This internal rotation is known as the electron’s spin; its quantized angular momentum component is known to be $`\pm \mathrm{}/2`$ (associated with a magnetic moment $`\mathrm{}e/2mc`$).
In the coordinate representation of spin, the direction of the electron’s axis of rotation may be specified by two spherical polar angles: $`\theta _s,\varphi _s`$; the subscript, $`s`$, distinguishes these angles from the $`\theta ,\varphi `$ that describe the angular position of the electron relative to the nucleus in the Schrödinger theory of the hydrogen atom; the $`\theta ,\varphi `$ coordinate system has its origin at the nucleus, whereas the $`\theta _s,\varphi _s`$ coordinate system has its origin at the electron.
The recent discovery of spherical harmonics, $`Y_{\mathrm{}}^m(\theta ,\varphi )`$, for half-odd-integer values of $`\mathrm{}`$ and $`m`$ , provides the basis for a coordinate representation of electron spin, the coordinates being the two spherical polar angles $`\theta _s,\varphi _s`$.
In this representation the operators for the square of the total spin angular momentum, $`𝐒^\mathrm{𝟐}`$, and its $`z`$-component, $`𝐒_𝐳`$, are expressed in terms of the spherical polar angles \[5, p.207\], \[6, p.95\] by:
$`𝐒^\mathrm{𝟐}=\mathrm{}^2\left\{{\displaystyle \frac{1}{\mathrm{sin}\theta _s}}{\displaystyle \frac{}{\theta _s}}\left(\mathrm{sin}\theta _s{\displaystyle \frac{}{\theta _s}}\right)+{\displaystyle \frac{1}{\mathrm{sin}^2\theta _s}}{\displaystyle \frac{^2}{\varphi _s^2}}\right\}𝐒_𝐳={\displaystyle \frac{\mathrm{}}{\mathrm{i}}}{\displaystyle \frac{}{\varphi _s}}`$ (1)
The eigenfunctions of these operators, $`\alpha (\theta _s,\varphi _s)`$, $`\beta (\theta _s,\varphi _s)`$, corresponding to the known spin-states of an electron are LABEL:JPhysA:
$`\alpha (\theta _s,\varphi _s)={\displaystyle \frac{1}{\pi }}\sqrt{\mathrm{sin}\theta _s}e^{i(+\varphi _s/2)}\mathrm{for}m_s=+1/2`$ (2)
$`\beta (\theta _s,\varphi _s)={\displaystyle \frac{1}{\pi }}\sqrt{\mathrm{sin}\theta _s}e^{i(\varphi _s/2)}\mathrm{for}m_s=1/2`$
where $`m_s`$ is the well-known spin quantum number, and the factor of $`1/\pi `$ normalizes the functions with respect to integration over $`\theta _s`$ and $`\varphi _s`$ (see below).
Application of the operators (1) to the functions (2) shows that both of these functions are eigenfunctions:
* of $`S^2`$ with eigenvalue $`3\mathrm{}^2/4`$, and
* of $`S_z`$ with eigenvalues of $`m_s\mathrm{}=\pm \mathrm{}/2`$.
The functions (2) are normalized in the sense that:
$$_0^{2\pi }𝑑\varphi _s_0^\pi Y^{}(\theta _s,\varphi _s)Y(\theta _s,\varphi _s)\mathrm{sin}(\theta _s)𝑑\theta _s=1$$
(3)
in which denotes the complex conjugate, and $`Y(\theta _s,\varphi _s)`$ is one of $`\alpha (\theta _s,\varphi _s)`$ or $`\beta (\theta _s,\varphi _s)`$. They are also orthogonal in the sense that:
$$_0^{2\pi }\alpha ^{}(\theta _s,\varphi _s)\beta (\theta _s,\varphi _s)𝑑\varphi _s=_0^{2\pi }\alpha (\theta _s,\varphi _s)\beta ^{}(\theta _s,\varphi _s)𝑑\varphi _s=0$$
(4)
The normalization (3) is the one usually used for the spherical harmonics. However, it should be noted that the Fermion functions (2) are only single-valued when $`\varphi `$ is considered to have the range of a double circle: $`0`$$``$$`\varphi _s`$$``$$`4\pi `$; this $`\mathrm{\hspace{0.25em}4}\pi `$ symmetry is a well known property of fermion wavefunctions \[7, p.21 & p.138\]. Alternatively, if the range of $`\varphi _s`$ is regarded as that of a single circle, $`0`$$``$$`\varphi _s`$$``$$`2\pi `$, they are double-valued functions in the sense that:
$`Y(\theta _s,\varphi _s+2\pi )=Y(\theta _s,\varphi _s)`$
The double-valuedness is admissible because the probability density,
$`Y^{}(\theta _s,\varphi _s)Y(\theta _s,\varphi _s)`$ is single-valued and positive. The alternative of basing the normalization upon $`0`$$``$$`\varphi _s`$$``$$`4\pi `$ would simply introduce a factor of $`\sqrt{2}`$ into the normalization constant.
## 2 Previous Representation of Spin
The history of intrinsic (spin) angular momentum is a fascinating story . The spin of an electron has been represented most simply by eigenfunctions, $`\alpha `$ and $`\beta `$, which are abstract in the sense that they are not functions of any coordinates; they are defined to be eigenfunctions of $`𝐒^\mathrm{𝟐}`$ and $`𝐒_𝐳`$ with eigenvalues of $`\mathrm{}^2s(s+1)`$ and $`\mathrm{}m_s`$, with $`s=\frac{1}{2}`$ and $`m_s=\pm \frac{1}{2}`$ :
$`𝐒^\mathrm{𝟐}\left\{\begin{array}{c}\alpha \\ \beta \end{array}\right\}=\mathrm{}^2\frac{1}{2}(\frac{1}{2}+1)\left\{\begin{array}{c}\alpha \\ \beta \end{array}\right\}𝐒_𝐳\left\{\begin{array}{c}\alpha \\ \beta \end{array}\right\}=\left\{\begin{array}{c}+\\ \end{array}\right\}\frac{\mathrm{}}{2}\left\{\begin{array}{c}\alpha \\ \beta \end{array}\right\}`$ (15)
Different authors handle the abstract nature of the spin eigenfunctions, $`\alpha `$ and $`\beta `$, differently. In presenting the orthonormality relations:
$`<\alpha |\alpha >=<\beta |\beta >=1`$ (16)
$`<\alpha |\beta >=<\beta |\alpha >=0`$ (17)
Dykstra \[9, p.185\] avoids explicit consideration of the functional dependence of the “abstract functions” $`\alpha `$ and $`\beta `$ on the “abstract spin coordinate”, and likewise avoids explicit consideration of the integration “over the spin coordinate” implicit in (16) and (17).
In this purely abstract interpretation of $`\alpha `$ and $`\beta `$ the simplest concept is to define them by the eigenvalue equations (15) with normalization defined by (16) (their orthogonality, (4), follows from (15) because they are eigenfunctions of $`S_z`$ with different eigenvalues).
Expectation values such as $`<\alpha |𝐒^\mathrm{𝟐}|\alpha >`$, $`<\beta |𝐒^\mathrm{𝟐}|\beta >`$, $`<\alpha |S_z|\alpha >`$, and $`<\beta |S_z|\beta >`$, and (16) are likewise abstract - they do not have to refer to “integration over the spin coordinate”.
Other authors \[10, p.15\], \[11, p.14\], \[12, p.253\], \[2, p.233\], invoke the Pauli 2$`\times `$2 matrix representation of electron spin to infer that the spin coordinate/variable spans a discrete space of just 2 values, which are $`\{1,0\}`$ for $`\alpha `$ and $`\{0,1\}`$ (or $`\{0,1\}`$ ) for $`\beta `$. Calais then correctly says that:
> “integration over spin space …is a misnomer for summation over spin space”
In the Pauli matrix representation this summation (in (16,17)) is the scalar product of the transpose of $`\alpha ,\beta `$ ($`\{1,0\}`$, $`\{0,1\}`$) with the 2$`\times `$1 column vectors. Identifying the 2 elements of the column-vector eigenfunctions of the Pauli representation as a space of 2 points spanned by the “spin coordinate”, is a dubious concept, for as Schiff shows \[2, pp.144-147\] the number of points of the spin-space increases with the total angular momentum quantum number, $`j`$: it is $`2j+1`$ for $`j=1/2,1,3/2`$, etc.
In many electron wavefunctions one must designate which electron is in a spin eigenstate, $`\alpha ,\beta `$. This is commonly indicated by $`\alpha (s_1),\beta (s_2)`$ etc, or simply $`\alpha (1),\beta (2)`$ etc. This notation indicates that $`\alpha (s_1)`$ denotes an $`\alpha `$ spin function for electron $`1`$, and some authors (naturally) refer to $`s_1,s_2`$, etc, as the “spin coordinate” of electron 1, 2, etc. In this context these electron designators are really specifying labels rather than coordinates. Similar designators will be needed in our spin-coordinate representation of many-electron wavefunctions; different angles, $`\theta _{s_j}`$, $`\varphi _{s_j}`$, ($`j=1,2,\mathrm{}`$) must be used for each electron.
## 3 Discussion
Heretofore no scalar representation in terms of the explicit (angular) coordinates of the spin motion has been used, and yet such positional coordinates are believed to have a privileged role in the interpretation of quantum mechanics .
The logical steps in the derivation of the Fermion Spherical Harmonics parallel the arguments used to derive the abstract eigenfunctions and eigenvalues of angular momentum, which (as is well-known) lead to both integer and half-odd-integer values of $`\mathrm{}`$ and $`m`$ \[14, §5.4,p.115\]. It was the belief that the wavefunction (rather than the probability) must be single-valued that inhibited discovery of these functions for so many years.
We hope that the explicit coordinates, $`\theta _s`$, $`\varphi _s`$, for the orientation of the spin vector of each electron will be adopted in the teaching of quantum chemistry, for they will make “integration over spin coordinates” a true integration entirely comparable with integration over the positional coordinates of each electron (such as $`r,\theta ,\varphi `$ for an atomic electron), and the spin eigenfunctions $`\alpha `$, $`\beta `$ will have an explicit functional dependence upon the angles $`\theta _s,\varphi _s`$ as shown explicitly in (2).
Apart from their pedagogical value, these angular dependent spin eigenfunctions will facilitate the construction of wavefunctions in which the orientation of a particular spin vector is neither parallel to, nor anti-parallel to, another spin vector or to an external magnetic field.
One possible application is to the theory of Nuclear Magnetic Resonance (NMR), where the spins of different nuclei will, in general, be neither parallel nor anti-parallel.
Another example is the Einstein-Podolski-Rosen gedanken experiment, in which two electrons go away from a source in opposite directions towards spin-orientation detectors which are inclined at any angle, $`\theta `$ (when parallel, $`\theta `$$`=`$$`0`$, or anti-parallel, $`\theta `$$`=`$$`\pi `$). The experiment measures correlations between detections at the two detectors as a function of their angle of inclination, $`\theta `$; further detail is beyond the scope of this article \[15, Ch.4,p.35\].
## Acknowledgements
This work was supported by the Natural Sciences and Engineering Research Council of Canada. |
warning/0507/cond-mat0507672.html | ar5iv | text | # Chaotic Hypothesis, Fluctuation Theorem and Singularities
## Introduction
There is a quite strong interest in stationary states of systems subject to the action of non conservative forces. These forces perform work on the system while by suitable mechanisms heat is extracted, so that the system can stay in a statistically stationary state. Theoretical and experimental works are steadily becoming avalaible on the matter. Theoretical work implement the heat extraction in several ways introducing “thermostat models”, which can be stochastic or deterministic forces.
A strong idealization of a system in a nonequilibrium steady state subject to deterministic forces is provided by the Anosov systems: their motion can be considered to be paradigm of chaotic behavior, playing in chaotic dynamics the role that harmonic motions play in regular dynamics. The chaoticity of the motions is immediately apparent from the definition of Anosov systems: locally around each point it has to be possible to draw three coordinate surfaces $`W_s,W_u,W`$ such that segments of curves on $`W_s,W_u`$ contract exponentially as time grows to $`+\mathrm{}`$ or, respectively, recedes to $`\mathrm{}`$ while segments on $`W`$, the one dimensional “neutral” flow direction, neither expand nor contract. They change their length but keep it of the same order of the initial one. If the dynamics is described by a map the neutral direction is omitted in the definition.
The chaotic hypothesis, Gallavotti and Cohen (1995a); Ruelle (1999), see below, proposes that chaotic systems should be considered as Anosov systems “for practical purposes”. This has several consequences: in particular about fluctuations in time reversible models, where the hypothesis leads to severe constraints through the Fluctuation Theorem. This is a mathematical property of the large deviations function of the phase space contraction of a time reversible Anosov map $`S`$. However, some obvious restrictions, analogous to the ones that are (often tacitly) assumed when one says that the “pendulum is isochronous” or that phonons in a crystal correspond to “harmonic excitations”, have to be taken into account when applying the hypothesis to realistic systems, which are not strictly Anosov systems.
The prediction has been tested in several simulations and we summarize the precise statement of it below: usually the results have been positive. However, there have been, in the literature, a few claims of failure of the chaotic hypothesis based on the apparent failure of the predictions of the fluctuation theorem. Here we concentrate on one such attempt, which studies systems violating the Anosov property because singularities of the interparticle potentials play an important role in the dynamics, Evans et al. (2005): a situation considered, correctly, in the literature as not important for most physical properties but which requires care if the fluctuation relation is specifically tested on such systems (in the same way care has to be used if isochrony is tested on a pendulum or harmonicity is tested in a crystal model).
Here we show that even in singular systems the chaotic hypothesis and the fluctuation theorem are not in contradiction: we develop a theory that extends the fluctuation theorem to singular systems continuing ideas that were introduced to study a special Gaussian noise thermostat, Van Zon and Cohen (2003), (this is a “random thermostat” not to be confused with the thermostats satisfying Gauss’ principle for some non holonomic constraint, like the isokinetic constraint).
The structure of the paper is the following: in section I we recall the basic notations and statements, and some alternative formulations of the fluctuation relation. We discuss its (trivial) form in equilibrium and how one can take a meaningful and non-trivial equilibrium limit. In section II we discuss the application of the chaotic hypothesis to singular systems. First we present a very simple example which shows that the effect of singularities is very important. Then we discuss how one can obtain quantitative predictions on the modification of the fluctuation relation due to the presence of singularities. Finally we discuss a prescription to remove singularities that follows from a careful examination of the proof of the fluctuation theorem for Anosov flows. The results are compared with recent numerical simulations. In section III we draw the conclusions and compare our interpretation with the one of Evans et al. (2005).
## I The fluctuation relation
We shall denote by $`\mathrm{\Omega }`$ the phase space (a smooth compact boundaryless Riemannian manifold), by $`S:\mathrm{\Omega }\mathrm{\Omega }`$ an invertible map on $`\mathrm{\Omega }`$ and by $`\sigma (x)`$ the volume contraction
$$\sigma (x)=\mathrm{log}|det_xS(x)|$$
(1)
Time reversal is defined as an isometry $`I:\mathrm{\Omega }\mathrm{\Omega }`$ with
$$IS=S^1I,\sigma (Ix)=\sigma (x)$$
(2)
If $`S`$ is an Anosov maps, existence of a unique invariant probability distribution $`\mu `$, called the SRB distribution and describing the long–time statistics of the motions whose initial data are chosen randomly with respect to the volume measure, is established, Ruelle (1995); Gallavotti et al. (2004). It has the property that, with the exception of points $`x\mathrm{\Omega }`$ in a set of $`0`$–volume, we have
$$\underset{\tau \mathrm{}}{lim}\frac{1}{\tau }\underset{t=0}{\overset{\tau 1}{}}F(S^tx)\stackrel{def}{=}F=_\mathrm{\Omega }F(y)\mu (dy)$$
(3)
for all smooth observables $`F`$ defined on phase space.
It is intuitive that “phase space cannot expand”; this is expressed by the following result of Ruelle \[Ruelle, 1996\]:
If $`\sigma _+\stackrel{def}{=}\sigma `$ it is $`\sigma _+0`$
Clearly if $`S`$ is volume preserving $`\sigma _+=0`$. If $`\sigma _+>0`$ the system does not admit any stationary distribution of the form $`\mu (dx)=\rho (x)dx`$, with density with respect to the volume measure $`dx`$ (often called absolutely continuous with respect to the volume).
This motivates calling systems for which $`\sigma >0`$ dissipative and conservative the others.
For Anosov systems which are transitive (i.e. with a dense orbit), reversible and dissipative one can define the dimensionless phase space contraction, a quantity often related to entropy creation rate (see \[Gallavotti, 2004a\]), averaged over a time interval of size $`\tau `$. This is
$$p(x)=\frac{1}{\sigma _+\tau }\underset{k=\tau /2}{\overset{\tau /21}{}}\sigma (S^kx)$$
(4)
provided of course $`\sigma _+>0`$.
Then for such systems the probability with respect to the stationary state, i.e. to the SRB distribution $`\mu `$, that the variable $`p(x)`$ takes values in $`\mathrm{\Delta }=[p,p+\delta p]`$ can be written as $`\mathrm{\Pi }_\tau (\mathrm{\Delta })=e^{\tau \mathrm{max}_{p\mathrm{\Delta }}\zeta (p)+O(1)}`$, where $`\zeta (p)`$ is a suitable function and, for any fixed choice of $`\mathrm{\Delta }`$ contained in an open interval $`(p^{},p^{})`$, $`p^{}1`$, the correction term at the exponent is $`O(1)`$ with respect to $`\tau ^1`$, as $`\tau \mathrm{}`$ (this is often informally expressed as $`lim_\tau \mathrm{}\frac{1}{\tau }\mathrm{log}\mathrm{\Pi }_\tau (p)=\zeta (p)`$ for $`p^{}<p<p^{}`$). The function $`\zeta (p)`$ is called in probability theory the rate function for the large deviations of $`p`$.
The function $`\zeta (p)`$ is analytic in $`p`$ and convex in the interval of definition $`(p^{},p^{})`$. Analyticity and convexity of large deviation rates are general properties, established by Sinai and valid for the SRB-averages of smooth observables (in Anosov systems), Sinai (1972, 1977); Gallavotti et al. (2004). In fact more can be said for the specific case of the large deviation rate of the observable $`p`$, and one can prove the following fluctuation theorem:
In transitive time reversible dissipative Anosov systems the rate function $`\zeta (p)`$ for the dimensionless phase space contraction $`p(x)`$ defined in (4) is analytic and strictly convex in an interval $`(p^{},p^{})`$ with $`+\mathrm{}>p^{}1`$ and $`\zeta (p)=\mathrm{}`$ for $`|p|>p^{}`$. Furthermore
$$\zeta (p)=\zeta (p)p\sigma _+,\mathrm{for}|p|<p^{}$$
(5)
which is called the “fluctuation relation” (FR).
Strict convexity follows from a theorem of Griffiths and Ruelle which shows that the only way strict convexity could fail is if $`\sigma (x)=\phi (Sx)\phi (x)+c`$ where $`\phi (x)`$ is a smooth function (typically a Lipschitz continuous function) and $`c`$ is a constant, see propositions (6.4.2) and (6.4.3) in \[Gallavotti et al., 2004\]. The constant $`c`$ vanishes if time reversal holds and $`\sigma (x)=\phi (Sx)\phi (x)`$ contradicts the assumption that $`\sigma _+>0`$, because $`\tau ^1_{\tau /2}^{\tau /21}\sigma (S^kx)=\tau ^1\left[\phi (S^{\tau /21}x)\phi (S^{\tau /2}x)\right]0`$ as $`\tau \mathrm{}`$.
The value of $`p^{}`$ must be $`p^{}1`$ otherwise the average of $`p`$ could not be $`1`$ (as it is by its very definition): it is defined, adopting the natural convention that $`\zeta (p)=\mathrm{}`$ for the values of $`p`$ whose probability goes to $`0`$ with $`\tau `$ faster than exponentially, as the infimum of the $`p>0`$ for which $`\zeta (p)=\mathrm{}`$. Alternatively $`\pm p^{}`$ are the asymptotic slopes as $`\lambda \pm \mathrm{}`$ of the Laplace transform $`\mathrm{log}e^{\lambda p}_{SRB}`$ Gallavotti (1995).
The fluctuation relation was discovered in a numerical experiment, \[Evans et al., 1993\], dealing with a non smooth system (hence not Anosov). The formulation and proof of the above proposition is in \[Gallavotti and Cohen, 1995a\] and in the context of Anosov systems the relation (5) is properly called the fluctuation theorem. The difference between this theorem and other fluctuation relations proposed in the literature has been clarified in Cohen and G.Gallavotti (1999). The theorem can be extended to Anosov flows (i.e. to systems evolving in continuous time), Gentile (1998).
### Alternative formulations
Sometimes, e.g. in Searles and Evans (2000); Evans et al. (2005), rather than the above $`p`$ the quantity $`a=\tau ^1_{j=\tau /2}^{\tau /21}\sigma (S^jx)`$ is considered and eq.(5) becomes
$$\stackrel{~}{\zeta }(a)=\stackrel{~}{\zeta }(a)a,\mathrm{for}|a|<a^{}p^{}\sigma _+$$
(6)
where $`\stackrel{~}{\zeta }(a)`$ is trivially related to $`\zeta (p)`$. This form dangerously suggests that in systems with $`\sigma _+=0`$ the distribution of $`a`$ is asymmetric (because the extra condition $`|a|<p^{}\sigma _+`$ might be forgotten, see \[Gallavotti, 2004b\]).
Note that $`p^{}`$ is certainly $`<+\mathrm{}`$ because the variable $`\sigma (x)`$ is bounded (being continuous on the bounded manifold on which the Anosov map is defined).
However no confusion should be made between $`p^{}\sigma _+`$ and $`\sigma _{max}\stackrel{def}{=}\mathrm{max}|\sigma (x)|`$: unlike $`\sigma _{max}`$ the quantity $`p^{}`$ is a non trivial dynamical quantity, independent on the metric used on phase space to measure distances, hence volume. This point has not been always understood and confusion has appeared in the published literature with unexpected consequences. In fact it is very easy to build examples of Anosov systems in which $`p^{}\sigma _+<\sigma _{max}`$: still, this does not mean that fluctuation relation is violated for such systems. Some explicit examples are discussed in next Section.
### Conservative systems and the equilibrium limit
Considering more closely the cases $`\sigma _+=0`$ it follows that $`\sigma (x)=\phi (Sx)\phi (x)`$ (again by the above mentioned result of Griffiths and Ruelle), with $`\phi `$ a smooth function of phase space. Hence the variable
$$a=\frac{1}{\tau }\underset{j=\tau /2}{\overset{\tau /21}{}}\sigma (S^jx)\frac{\phi (S^{\frac{\tau }{2}}x)\phi (S^{\frac{\tau }{2}}x)}{\tau }$$
(7)
is bounded and tends to $`0`$ uniformly. One could repeat the theory developed for $`p`$ when $`\sigma _+>0`$ but one would reach the conclusion that $`\stackrel{~}{\zeta }(a)=\mathrm{}`$ for $`|a|>0`$ and we see that the result is trivial. In fact in this case it follows that the system admits an absolutely continuous SRB distribution. The distribution of $`a`$ is symmetric (trivially by time reversal symmetry) and becomes a delta function around $`0`$ as $`\tau \mathrm{}`$.
Nevertheless the fluctuation relation is non trivial in cases in which the map $`S`$ depends on parameters $`\underset{¯}{E}=(E_1,\mathrm{},E_n)`$ and becomes volume preserving (“conservative”) as $`\underset{¯}{E}\underset{¯}{0}`$: in this case $`\sigma _+0`$ as $`\underset{¯}{E}\underset{¯}{0}`$ and one has to rewrite the fluctuation relation in an appropriate way to take a meaningful limit.
The result is that the limit as $`\underset{¯}{E}\underset{¯}{0}`$ of the fluctuation relation in which both sides are divided by $`\underset{¯}{E}^2`$ makes sense and yields (in the case considered here of transitive Anosov dynamical systems) relations which are non trivial and that can be interpreted as giving Green–Kubo formulae and Onsager reciprocity for transport coefficients, \[Gallavotti, 1996a; Gallavotti and Ruelle, 1997\].
In fact the very definition of the duality between currents and fluxes so familiar in nonequilibrium thermodynamics since Onsager can be set up in such systems using as generating function the $`\sigma _+`$ regarded as a function of $`\underset{¯}{E}`$. Note that the fluxes are usually “currents” divided by the temperature: therefore via the above interpretation one can try to define the temperature even in nonequilibrium situations, \[Gallavotti and Cohen, 2004; Gallavotti, 2004a; Zamponi et al., 2005\].
## II Singular systems
The fluctuation relation has been proved only for Anosov systems. However, a Chaotic Hypothesis has been proposed, which states that, for the purpose of studying the physically interesting observables, a chaotic dynamical system can be considered as an Anosov system, Gallavotti and Cohen (1995a, b); Gallavotti (2000).
In applying the chaotic hypothesis to singular systems, e.g. a system of particles interacting via a Lennard-Jones potential (which is infinite in the origin), one might encounter apparent difficulties. We will discuss them in the following.
### The effect of a change of metric for Anosov flows
The simplest example (out of many) is provided by the simplest conservative system which is strictly an Anosov transitive system and which has therefore an SRB distribution: this is the geodesic flow $`S_t`$ on a surface of constant negative curvature, \[Bonetto et al., 2000\]. We discuss here an evolution in continuous time because the matter is considered in the literature for such systems, \[Gallavotti, 2004b\] (even simpler examples are possible for time evolution maps).
The phase space $`M`$ is compact, time reversal is just momentum reversal and the natural metric, induced by the Lobatchevsky metric $`g_{ij}(q)`$ on the surface, is time reversal invariant: the SRB distribution is the Liouville distribution and $`\sigma (x)0`$. However one can introduce a function $`\mathrm{\Phi }(x)`$ on $`M`$ which is very large in a small vicinity of a point $`x_0`$, arbitrarily selected, constant outside a slightly larger vicinity of $`x_0`$ and positive everywhere. A new metric could be defined as $`g_{new}(x)=(\mathrm{\Phi }(x)+\mathrm{\Phi }(Ix))g(x)\stackrel{def}{=}e^{F(x)}g(x)`$: it is still time reversal invariant but its volume elements will no longer be invariant under the time evolution $`S_t`$ associated with the geodesic flow with respect to the Lobatchevsky metric. The rate of change of phase space volume in the new metric will be $`\sigma _{new}(x)=\frac{d}{dt}F(x)`$. Then the phase space contraction $`\sigma _{new}(x)`$ takes values that not only are not identically $`0`$ but which can in general be arbitrarily large, depending on the specific choice of $`\mathrm{\Phi }(x)`$. The distribution of $`a=\frac{1}{\tau }_0^\tau \sigma _{new}(S_tx)𝑑t=\tau ^1\left[F(S_\tau x)F(x)\right]`$, at any finite time, will violate (6), simply because it is symmetric around $`0`$, by time reversal.
In the limit $`\tau \mathrm{}`$, as long as $`F(x)`$ is bounded, $`a=\tau ^1\left[F(S_\tau x)F(x)\right]\begin{array}{c}\tau \mathrm{}\hfill \end{array}\mathrm{\hspace{0.17em}0}`$ uniformly in $`x`$, as in the corresponding map case, and the SRB distribution of $`a`$ will tend to a delta function centered in $`0`$ (hence $`\stackrel{~}{\zeta }(a)=\mathrm{}`$ for $`a0`$). However, if $`F(x)`$ is not bounded (e.g. if it is allowed to become infinite in $`x_0`$) this is not the case in general, as we shall discuss in detail in next section.
### The effect of singular boundary terms
One can realize that terms of the form $`\tau ^1\left[F(S_\tau x)F(x)\right]`$ with $`F(x)`$ not bounded can affect the large fluctuations of $`\sigma (x)`$, at least if the probability of an arbitrarily large value of $`F`$ is not too small, i.e. if asymptotically for big values of $`F`$ it is exponentially small in $`F`$ (or larger), e.g. it is of the form $`e^{\kappa F}`$, for some constant $`\kappa >0`$. This is a valuable and interesting remark brought up for the first time, and correctly interpreted, already in Van Zon and Cohen (2003) and in the following papers Van Zon and Cohen (2004); Van Zon et al. (2004). The analysis of Van Zon and Cohen (2003, 2004); Van Zon et al. (2004) applies to cases where the unbounded fluctuations are driven by an external white noise. In the following we extend the theoretical analysis in Van Zon and Cohen (2004); Van Zon et al. (2004) to cases in which the unbounded fluctuations do not arise from a Gaussian noise but from a deterministic evolution like the ones in Searles and Evans (2000); Evans et al. (2005): this is a simple extension of the main idea and method of Van Zon and Cohen (2004) and provides an alternative interpretation to the analysis in Searles and Evans (2000); Evans et al. (2005).
Our analysis can be applied to the example of the Anosov flow with singular metric considered above and to more realistic systems: among them systems of particles interacting via an unbounded potential (like a Lennard–Jones (LJ) or a Weeks–Chandler–Andersen (WCA) potential), driven by an external field and subject to an isokinetic or a Nosé–Hoover thermostat. To be definite one can consider system of $`N`$ particles in $`d`$ dimensions, described by evolution equations $`\dot{𝐩}_i=𝐄_{𝐪_i}\mathrm{\Phi }\alpha 𝐩_i`$, $`\dot{𝐪_i}=𝐩_i`$. For an isokinetic Gaussian thermostat, $`\alpha `$ is a function of $`𝐩_i`$, chosen so to keep the total kinetic energy fixed to $`_i𝐩_i^2=Nd\beta ^1`$. For a Nosé–Hoover thermostat $`\alpha (t)`$ is a variable independent of $`𝐪_i(t),𝐩_i(t)`$ and satisfying the evolution equation $`\dot{\alpha }=\frac{1}{Q}\left[_i𝐩_i^2Nd\beta ^1\right]`$, with $`Q,\beta >0`$ parameters.
In both cases the phase space contraction $`\sigma (x)`$ has the form $`\sigma _0(x)\beta \frac{d}{dt}V(x)`$, where $`\beta `$ has the interpretation of inverse temperature. In the isokinetic case, $`\sigma _0(x)`$ is bounded, and $`V=\mathrm{\Phi }`$. In the Nosé–Hoover case $`\sigma _0(x)`$ has, in the SRB distribution, a fast decaying tail (Gaussian at equilibrium, and likely to remain such in presence of external forcing) and $`V=_i\frac{𝐩_i^2}{2}+\mathrm{\Phi }(𝐪)+Q\frac{\alpha ^2}{2}`$, Nosé (1984); Hoover (1985).
In both cases, in equilibrium, the SRB probability of $`V`$ has an exponential tail $`e^{\beta V}`$ (possibly with power-law corrections). For the purpose of illustration we assume, from now on, that the same happens in presence of the force $`𝐄`$. This is an essential and far from obvious assumption useful, as discussed below, to understand the possible role of the singularities, but it should not be assumed lightly as it is well known that the SRB distributions may have very peculiar $`E`$ dependence and, at the moment, a not intuitive character, Derrida et al. (2002); Bertini et al. (2001). Nevertheless, in preliminary numerical simulations, it seems approximately correct, at least within the accuracy of the numerical data and for $`|𝐄|`$ not too large; furthermore the analysis that follows can be naturally adapted to more general assumptions on the tails.
In such cases the non normalized variable $`a`$ (introduced before Eq. (6)) has the form $`a_0+\frac{\beta }{\tau }(V_iV_f)`$ where $`V_i,V_f`$ are the values of $`V(x)`$ at the initial and final instants of the time interval of size $`\tau `$ on which $`a`$ is defined, and $`a_0\stackrel{def}{=}\frac{1}{\tau }_0^\tau \sigma _0(S_tx)𝑑t`$:
$$a=\frac{1}{\tau }_0^\tau \sigma (S_tx)𝑑ta_0+\frac{\beta }{\tau }(V_iV_f)$$
(8)
If the system is chaotic and $`\tau `$ is large, the variables $`a_0,V_i,V_f`$ can be regarded as independently distributed, because $`a_0`$ depends essentially only on the length $`\tau `$ of the time interval, while $`V_i`$ and $`V_f`$ depend on the precise locations of the extremes of the interval. Moreover the distribution of $`V=V_i`$ or $`V=V_f`$ is essentially $`e^{\beta V}dV`$ to leading order as $`V\mathrm{}`$, as discussed above. Therefore the rate function of the variable $`a`$ can be computed as
$$\begin{array}{cc}& \underset{\tau \mathrm{}}{lim}\frac{1}{\tau }\mathrm{log}_{p^{}\sigma _+}^{p^{}\sigma _+}da_0_0^{\mathrm{}}dV_i_0^{\mathrm{}}dV_f\hfill \\ & e^{\tau \stackrel{~}{\zeta }_0(a_0)\beta V_i\beta V_f}\delta [\tau (aa_0)+\beta V_i\beta V_f]\hfill \\ & =\underset{\tau \mathrm{}}{lim}\frac{1}{\tau }\mathrm{log}_{p^{}\sigma _+}^{p^{}\sigma _+}𝑑a_0e^{\tau \stackrel{~}{\zeta }_0(a_0)\tau |aa_0|}\hfill \end{array}$$
(9)
where $`\stackrel{~}{\zeta }_0(a_0)`$ is the rate function of $`a_0`$; thus
$$\stackrel{~}{\zeta }(a)=\underset{a_0[p^{}\sigma _+,p^{}\sigma _+]}{\mathrm{max}}\left[\stackrel{~}{\zeta }_0(a_0)|aa_0|\right]$$
(10)
Defining $`a_{}`$ by $`\stackrel{~}{\zeta }_0^{}(a_{})=\pm 1`$, by the strict convexity of $`\stackrel{~}{\zeta }_0(a_0)`$ it follows
$$\stackrel{~}{\zeta }(a)=\{\begin{array}{cc}\stackrel{~}{\zeta }_0(a_{})a_{}+a,\hfill & a<a_{}\hfill \\ \stackrel{~}{\zeta }_0(a),\hfill & a[a_{},a_+]\hfill \\ \stackrel{~}{\zeta }_0(a_+)+a_+a,\hfill & a>a_+\hfill \end{array}$$
(11)
If we assume that $`\stackrel{~}{\zeta }_0(a_0)`$ satisfies FR (as expected from the chaotic hypothesis, see below), then $`\stackrel{~}{\zeta }_0(a_0)=\stackrel{~}{\zeta }_0(a_0)+a_0`$ and by differentiation it follows that $`a_{}=\sigma _+`$, where $`\sigma _+`$ is the location of the maximum of $`\stackrel{~}{\zeta }_0`$, i.e. is the average of $`a`$, and that $`\stackrel{~}{\zeta }_0(a_{})=\stackrel{~}{\zeta }_0(\sigma _+)=\stackrel{~}{\zeta }_0(\sigma _+)\sigma _+=\sigma _+`$. Moreover it is clear that $`a_+>\sigma _+`$ because $`\stackrel{~}{\zeta }_0^{}(a_+)<0`$. Using these informations one can show that, for $`a0`$:
$$\stackrel{~}{\zeta }(a)\stackrel{~}{\zeta }(a)=\{\begin{array}{cc}a,\hfill & a<\sigma _+\hfill \\ \stackrel{~}{\zeta }_0(a)+a,\hfill & \sigma _+aa_+\hfill \\ \stackrel{~}{\zeta }_0(a_+)+a_+,\hfill & a>a_+\hfill \end{array}$$
(12)
It follows that, if $`\stackrel{~}{\zeta }_0(a_0)`$ satisfies FR up to $`a=p^{}\sigma _+`$, then $`\stackrel{~}{\zeta }(a)`$ satisfies FR only in the interval $`|a|<|a_{}|=\sigma _+`$. Outside this interval $`\stackrel{~}{\zeta }(a)`$ does not satisfy the FR and in particular for $`aa_+`$ it is $`\stackrel{~}{\zeta }(a)\stackrel{~}{\zeta }(a)=const.`$, as already described in Van Zon and Cohen (2004). Eq. (12) is the generalization of the result of Van Zon and Cohen (2004) to the case where $`\stackrel{~}{\zeta }_0(a_0)`$ is not Gaussian.
Translated into the normalized variables $`p_0=a_0/\sigma _+`$ and $`p=a/\sigma _+`$, this means that, even if the rate function of $`p_0`$ satisfies FR up to $`p^{}>1`$, the rate function of $`p`$ verifies FR only for $`|p|1`$. This is the effect due to the presence of the singular boundary term. Note that the scenario above applies only to the case in which $`V_i,V_f`$ are unbounded and have exponential tails. A repetition of the discussion above in the case that $`V_i,V_f`$ are unbounded but with tails faster than exponential would lead to the conclusion that $`\stackrel{~}{\zeta }(a)=\stackrel{~}{\zeta }_0(a)`$. In particular if $`V_i,V_f`$ are assumed to be bounded $`\stackrel{~}{\zeta }(a)=\stackrel{~}{\zeta }_0(a)`$. Of course in these cases the times of convergence of $`\stackrel{~}{\zeta }(a)`$ to $`\stackrel{~}{\zeta }_0(a)`$ will depend on the details of the tails of $`V_i,V_f`$ (for instance if $`V_i,V_f`$ are bounded by a constant $`B`$, the times of convergence will grow with $`B`$). Note also that the result above does not depend on the details of the distribution of $`V_i,V_f`$ for small $`V`$ (in particular it does not depend on the lower cutoff $`V=0`$ assumed in Eq. 9).
An example of $`\stackrel{~}{\zeta }(a)`$ is reported in Fig. 1: it is a simple stochastic model for the FT (taken from Sect. 5 in Bonetto et al. (1997), see also the extensions in Lebowitz and Spohn (1999); Maes (1999)). The example is the Ising model without interaction in a field $`h`$, i.e. a Bernoulli scheme with symbols $`\pm `$ with probabilities $`p_\pm =\frac{e^{\pm h}}{2\mathrm{cosh}h}`$. Defining $`a_0=\frac{1}{\tau }_{i=0}^{\tau 1}2h\sigma _i`$, so that $`\sigma _+=a_0=2h\mathrm{tanh}h`$, and setting $`x\stackrel{def}{=}\frac{1+a_0/(2h)}{2}`$, and $`s(x)=x\mathrm{log}x(1x)\mathrm{log}(1x)`$, one computes $`\stackrel{~}{\zeta }_0(a_0)=s(x)+\frac{1}{2}a_0+const`$ which is not Gaussian and it is defined in the interval $`[a^{},a^{}]`$ with $`a^{}=2h`$. In this case the large deviation function $`\stackrel{~}{\zeta }_0(a_0)`$ satisfies FR for $`|a_0|a^{}`$. If a singular term $`V=\mathrm{log}(_{i=0}^{\mathrm{}}2^{i1}\frac{\sigma _i+1}{2})`$ is added to $`a_0`$, defining $`a=a_0+\beta (V_iV_f)`$ (with $`\beta =\mathrm{log}_2(1+e^{2h})`$ so that the probability distribution of $`V`$ is $`e^{\beta V}`$ for large $`V`$), the resulting $`\stackrel{~}{\zeta }(a)`$ does not verify FR for $`a>a=2h\mathrm{tanh}h`$. In particular, for $`h0`$, the interval in which the FR is satisfied vanishes.
### How to remove singularities
From the discussion above it turns out that singular terms which are proportional to total derivatives of unbounded functions (like the term $`\frac{dV}{dt}`$ that appears in the phase space contraction rate of thermostatted systems) can induce “undesired” (or “unphysical”) modifications of the large deviations function $`\zeta (p)`$.
On heuristic grounds, when dealing with singular systems, one could follow the prescription that unbounded terms in $`\sigma (x)`$ which are proportional to total derivatives should be subtracted from the phase space contraction rate. If the resulting $`\sigma _0(x)`$ is bounded (as it is e.g. for the Gaussian isokinetic thermostat models considered) or at least if the tails of its distribution decay faster than exponentially, then its large deviations function should verify the FR for $`|p|p^{}`$, $`p^{}`$ being the intrinsic dynamic quantity defined above.
Note that after the subtraction of the divergent terms the remaining contraction, in the considered cases, is bounded for isokinetic thermostats or has a tail decaying faster than exponential in the case of Nosé–Hoover thermostats. In the following for definiteness we will assume $`\sigma _0`$ bounded but the same discussion is valid for $`\sigma _0`$ unbounded with tails decaying faster than exponential.
If the singular terms are not subtracted, the FR will appear to be valid only for $`|p|1`$ even if $`p^{}>1`$. This seems to have generated statements that the Chaotic Hypothesis does not apply to isokinetic systems, see Evans et al. (2005).
The heuristic prescription above can be motivated by a careful analysis of the proof of the fluctuation theorem for Anosov flows. In the following let us call again $`a`$ the integral of the total phase space contraction rate $`\sigma (x)`$ (which includes singular terms) and $`a_0`$ the integral of the bounded variable $`\sigma _0(x)`$ from which singular total derivatives have been removed.
The fluctuation theorem was proved in Gallavotti and Cohen (1995a); Gallavotti (1996b); Ruelle (1999) for Anosov maps and only later it has been extended in Gentile (1998) to Anosov flows. Very sketchily, the extension of the fluctuation theorem to Anosov flows in Gentile (1998) is proved as follows. One reduces the Anosov flow on $`\mathrm{\Omega }`$ to a map via a Poincaré’s section, associated with surfaces on $`\mathrm{\Omega }`$ transversal to the flow. The passage of the flow through any one of such surfaces is called a timing event. The map between two consecutive timing events is called a “Poincaré’s map”. The union $`\mathrm{\Omega }_P`$ of the surfaces represents the phase space of the Poincaré’s map. The surfaces in $`\mathrm{\Omega }_P`$ can be suitably chosen, in such a way that the Poincaré’s map is a chaotic map which although not smooth, hence not an Anosov map, has (a non trivial fact Gentile (1998)), all the properties necessary to prove the fluctuation theorem (which therefore applies to systems more general than the Anosov maps, although there is not a general characterization of the systems which are not Anosov and to which it applies). So, for such a map the fluctuation theorem holds and this in turn leads to a FR for the flow by the theory in Gentile (1998) under the assumption that the variable $`\sigma (x)`$ is bounded.
If, as in the case under analysis, $`\sigma (x)`$ is not bounded, we can interpret the chaotic hypothesis as applying to the map associated with a Poincaré’s section which avoids the singularities of the potential, a very natural prescription which allows us to apply the theory in Gentile (1998) and derive a FR for both the map and the flow. For instance, we can choose as timing events the instants in which either the potential energy or the Nosé’s “extended Hamiltonian” exceed some fixed value $`\overline{V}`$. If we make this choice, the (discrete) average $`\widehat{a}`$ of $`\sigma (x)`$ over a sequence of iterations of the Poincaré’s map will coincide with the (discrete) average $`\widehat{a}_0`$ of $`\sigma _0(x)`$ along the same sequence: this simply follows from the remark that by construction the total increment of $`\sigma (x)\sigma _0(x)`$ between two timing events, given by $`\beta (V_fV_i)`$, is $`0`$ (by construction $`\mathrm{\Omega }_P`$ is chosen as a subset of $`\{x\mathrm{\Omega }:V(x)=\overline{V}\}`$ where $`V_f=V_i`$). Then, by the same argument in Gentile (1998), the fact that the rate function of $`\widehat{a}_0`$ satisfies a FR and that $`\sigma _0(x)`$ is bounded implies that the rate function of the continuous average $`a_0`$ of $`\sigma _0(x)`$ along a trajectory of the flow will satisfy the fluctuation theorem.
Therefore the distribution of $`a_0`$ will satisfy the FR (by the chaotic hypothesis) for $`|a_0|<p^{}\sigma _+`$. By the above maximum argument, the distribution of $`a`$ will also verify, as a consequence, the FR but only for $`|a|\sigma _+`$, i.e. in the form (12).
Then the (natural) prescription to study FR for chaotic flows is to reduce the problem to a chaotic map considering only Poincaré’s sections which do not pass through a singularity of $`\sigma (x)`$. The sum of $`\sigma (x)`$ over a large number of timing events on such sections is equal to the time integral of $`\sigma _0(x)`$ plus a bounded term which can be neglected. Thus the prescription on the choice of Poincaré’s sections is equivalent to the heuristic prescription of removing from $`\sigma (x)`$ all the unbounded total derivatives.
It follows that the chaotic hypothesis leads to a clear prediction on the outcome of possible numerical simulations of particle systems interacting via unbounded potentials and subject to the isokinetic or the Nosé–Hoover thermostat: the FR will hold for all $`|a|\sigma _+`$ and, once the term $`\frac{dV}{dt}`$ is removed, for all $`|a_0|<p^{}\sigma _+`$ with $`p^{}1`$. Note that in the cases under analysis $`a_0`$ coincides with the dissipation function of Evans and Searles that was in fact predicted to satisfy FR Searles and Evans (2000); Evans et al. (2005), even though for different reasons. We believe that the correct interpretation of the fact that FR for $`\stackrel{~}{\zeta }_0(a_0)`$ holds for all $`|a_0|<p^{}\sigma _+`$ is the one given above.
The numerical results of Dolowschiák and Kovács (2005); Zamponi et al. (2004); Giuliani et al. (2005) agree with the prediction that FR for the rate function of $`a_0`$ is valid even beyond $`a_0=\sigma _+`$. The prediction that (at least near equilibrium) the rate function of $`a`$ should satisfy FR only up to $`a=\sigma _+`$ and that should become linear for $`aa_+`$ at the moment has been experimentally confirmed only in Gaussian cases Van Zon and Cohen (2003, 2004); Van Zon et al. (2004). It would be very interesting to investigate in detail the structure of $`\stackrel{~}{\zeta }(a)`$ even in non Gaussian cases. Note that this is far from being an easy task (in particular the analysis in Giuliani et al. (2005) was not sophisticated enough to study this problem). In fact, as discussed in detail in Zamponi et al. (2004), the presence in the definition of $`\sigma `$ of a total derivative of an unbounded function may enlarge of $`2`$ orders of magnitudes the times needed for the probability distribution of $`a`$ to reach its asymptotic shape: even in the Gaussian region (small fluctuations of $`a`$ around $`\sigma _+`$) the convergence times for $`\stackrel{~}{\zeta }(a)`$ are found to be of order $`1000`$ decorrelation times, versus a time of order $`10`$ decorrelation times needed for $`\stackrel{~}{\zeta }_0(a_0)`$ to converge to its asymptotic shape Zamponi et al. (2004). Clearly, for times of order $`1000`$ decorrelation times, it is very hard to observe fluctuations of $`a`$ larger than $`a_+\sigma _+`$. In order to verify the prediction for the shape of $`\stackrel{~}{\zeta }(a)`$ beyond $`a=a_+`$ an experiment specifically designed for this purpose would be needed, together with a detailed investigation of the finite time corrections to $`\stackrel{~}{\zeta }(a)`$, along the lines in Giuliani et al. (2005).
## III Conclusions and remarks
We showed that the Chaotic Hypothesis can be applied even to singular chaotic systems (in particular even to Gaussian isokinetic or Nosé–Hoover thermostatted systems), by identifying their macroscopic behavior with that of reversible Anosov systems with singular metric. Reversible Anosov systems with singular metric are systems to which the mathematical analysis usually leading to FR can be (rigorously) repeated to lead to a modified FR, illustrated by Eq.s (11), (12). Note that for $`\sigma _+0`$ Eq. (11) tends to the distribution $`\stackrel{~}{\zeta }(a)=|a|`$ violating the usual FR (simply because the limiting distribution is symmetric, by time reversal). For Anosov systems with singular metric, the prescription to avoid oddities (i.e. to avoid a modified FR) is to subtract from $`\sigma (x)`$ a total derivative $`dV/dt`$, in such a way that the variable $`\sigma _0=\sigma dV/dt`$ is bounded or has faster than exponential tails. The distribution of $`\sigma _0`$ will verify the FR also for $`|p|>1`$. This prescription is equivalent to the very reasonable prescription that the Poincaré’s section used for mapping the flow into a map does not pass through a singularity of $`\sigma (x)`$. Accepting the Chaotic Hypothesis, we propose to apply the same prescription to remove singularities to singular chaotic systems. Our prescription coincides with other prescriptions proposed earlier (for different reasons) in the literature.
The analysis in Sect. 2 above applies as well to understand how to apply the FR to systems with Gaussian (or unbounded) noise and the compatibility between the general theory of Kurchan (1998); Lebowitz and Spohn (1999); Maes (1999) with the works Van Zon and Cohen (2004) and Van Zon et al. (2004); Garnier and Ciliberto (2005).
Note that the picture we propose is different from the interpretation of the apparent violations to FR in singular systems proposed recently in Evans et al. (2005), where in particular it is argued that FR and CH do not hold for thermostatted systems near equilibrium. We conclude by comparing more closely our discussion with the corresponding discussion in Evans et al. (2005).
(1) As stressed above it is dangerous (and wrong) to consider (6) without the restriction $`|a|p^{}\sigma _+`$ as the prediction of fluctuation theorem. In Evans et al. (2005) the authors, after having correctly pointed out this point, seem (quite surprisingly!) to forget about this condition in the following. For instance, when studying the problem of approach to equilibrium, in order to show a contradiction between FR and GK relations, they assume that the relation Eq.(6) without the condition $`|a|<p^{}\sigma _+`$ “is correct both at equilibrium and near equilibrium” and they proceed to infer from this a contradiction. Of course such assumption is wrong and the fact that from this contradictions follow is not an argument against CH or FR.
(2) An argument in Evans et al. (2005) is supposed to prove that the relation $`\stackrel{~}{\zeta }(a)=\stackrel{~}{\zeta }(a)a`$ (without the condition $`|a|<p^{}\sigma _+`$) holds for reversible Anosov systems for all $`a`$’s, also for $`\sigma _+=0`$ (in particular they say that “the division by $`\sigma _+`$ does not seem to be necessary for the proof in ”). This is not the case: at equilibrium as well as near equilibrium, as remarked above and as illustrated also by Bonetto et al. (2000), there are examples of systems for which the proof of FR can be rigorously repeated step by step but for which the correct conclusion of the proof is that the relation $`\stackrel{~}{\zeta }(a)=\stackrel{~}{\zeta }(a)a`$ is violated for $`a>p^{}\sigma _+`$. For instance, this is the case for a conservative Anosov flow with singular metric (in which the relation above is violated trivially by time reversal). These counterexamples show that the assumption $`\sigma _+>0`$ is, instead, essential for the proof of fluctuation theorem. The necessity of the assumption $`\sigma _+>0`$ is stressed in the early paper Gallavotti (1995) which the Authors of Evans et al. (2005) quote; it is stressed also in the paper Ruelle (1999) which also makes clear that $`\stackrel{~}{\zeta }(a)=\stackrel{~}{\zeta }(a)a`$ can only hold under the assumption that $`|a|`$ does not exceed a maximum value.
(3) The analysis in Sect. 2 above shows that the probability distribution describing isokinetic systems near equilibrium are SRB distributions (contrary to what is claimed in Evans et al. (2005)): this is mathematically obvious by the very definition of SRB distribution in the case of Anosov systems (even if isokinetic or in general with singular metric, see e.g. the geodesic flow discussed above) and it appears to be true also in non Anosov systems that have so far been considered.
(4) In the case of the thermostatted particle systems considered in Evans et al. (2005) the unbounded derivative $`\frac{dV}{dt}`$ is also the contraction rate of the volume in equilibrium, i.e. with $`𝐄=0`$. Thus, for $`𝐄\mathrm{𝟎}`$, one can remove the total derivative from $`\sigma (x)`$ simply considering the contraction with respect to the equilibrium invariant distribution $`e^{\beta V}`$, as stated in Evans et al. (2005); Searles and Evans (2000). However, this observation does not provide a general prescription to remove the singular part from the phase space contraction rate because it rests on the very special fact that the singularities of the function $`V(x)`$ (i.e. of the potential $`\mathrm{\Phi }`$) do not depend on $`𝐄`$. The prescription that the phase space contraction should be computed on non singular Poincaré’s sections, instead, does not require any other assumption. In general the two prescriptions and the corresponding predictions differ and we believe that in general the prescription of computing $`\sigma `$ with respect to the equilibrium invariant distribution has not the desired effect of removing all singularities (then in general $`\sigma `$ with respect to the equilibrium invariant distribution could violate FR for $`\sigma _+<a<p^{}\sigma _+`$).
###### Acknowledgements.
We are indebted to E.G.D. Cohen, and R. Van Zon for many enlightening discussions. F.Z. wish to thank G.Ruocco for many useful discussions. G.G is indebted to Rutgers University, I.H.E.S and École Normale Supérieure, where he spent periods of leave while working on this project.
e-mail: bonetto@math.gatech.edu
giovanni.gallavotti@roma1.infn.it
alessandro.giuliani@roma1.infn.it
francesco.zamponi@phys.uniroma1.it
web: http://ipparco.roma1.infn.it
Dip. Fisica, U. Roma 1,
00185, Roma, Italia
REVTEX |
warning/0507/cond-mat0507242.html | ar5iv | text | # Exact transformation for spin-charge separation of spin-1/2 Fermions without constraints
## Transformation to quasiparticle operators.
Let $`n_r`$ be the particle number at site $`r`$, and $`n_{}`$ and $`n_{}`$ be respectively the number of up and down electrons: $`n_r=n_{,r}+n_{,r}`$. Since the Hubbard interaction obeys $`n_{,r}n_{,r}=\frac{1}{2}(n_r1+(n_r1)^2)`$, by choosing the chemical potential appropriately we can replace the usual Hubbard interaction with the particle-hole symmetric interaction $`\frac{1}{2}U(n_r1)^2`$. We thus write the Hubbard model in the particle-hole symmetric form $`H=t_{rr^{}}(c_{s,r}^{}c_{s,r^{}}^{}+CC)+U_r\frac{1}{2}(n_r1)^2`$ where the sum is over all sites and nearest neighbor bonds counted once.
We now demonstrate an exact transformation of the local operator algebra of the Hubbard model to new Fermionic $`CP^1`$ “quasicharge” $`\widehat{c}_r^{}`$ and $`SU(2)`$ “quasispin” operators $`q_r^i`$ that obey, respectively, Fermi and Bose statistics. The operators are given by
$`\widehat{c}_r^{}`$ $`=`$ $`c_{,r}^{}(1n_{,r})+(1)^rc_{,r}^{}n_{,r}`$ (1)
$`q_r^+`$ $`=`$ $`(c_{,r}^{}(1)^rc_{,r}^{})c_{,r}^{}`$ (2)
$`q_r^{}`$ $`=`$ $`(q_r^+)^{}`$
$`q_r^z`$ $`=`$ $`\frac{1}{2}n_{,r}`$
where $`1^r`$ is $`\pm 1`$ depending on which sublattice site $`r`$ is on. We define the $`x`$ and $`y`$ component of quasispin by $`q_r^x=\frac{1}{2}(q_r^++q_r^{})`$ and $`q_r^y=\frac{1}{2i}(q_r^+q_r^{})`$. The factor appearing as $`(1)^r`$ can in fact be chosen as an arbitrary phase, but is fixed here to preserve symmetries specific to the Hubbard model. This ”nonlinear” transformation generalizes earlier work on transformations that were strictly canonical, and hence did not alter commutation relations of the fermi operators.mele ; gunn
The new operators obey the following algebra, which can be verified by straightforward manipulations: $`\{\widehat{c}_r^{},\widehat{c}_r^{}^{}\}=\delta _{r,r^{}}`$ , $`\{\widehat{c}_r^{},\widehat{c}_r^{}^{}\}=0`$ , $`[\widehat{c}_r^{},q_r^{}^i]=0`$ , $`[q_r^i,q_r^{}^j]=i\delta _{rr^{}}_kϵ_{ijk}q_r^k`$
We thus find that $`q_r^i`$, $`\widehat{c}_r^{}`$ and $`\widehat{c}_r^{}`$ are respectively independent spin half bosonic and “spinless” Fermionic operators. The usual operators are given by
$`c_{,r}^{}`$ $`=`$ $`\widehat{c}_r^{}(\frac{1}{2}+q_r^z)+(1)^r\widehat{c}_r^{}(\frac{1}{2}q_r^z)`$ (3)
$`c_{,r}^{}`$ $`=`$ $`q_r^{}(\widehat{c}_r^{}(1)^r\widehat{c}_r^{}).`$
We define the quasicharge operator as $`n_r^f=\widehat{c}_r^{}\widehat{c}_r^{}`$. We find the following relations with the usual operators expressed in terms of ordinary Fermions: $`n_r=12n_r^fq_r^z`$, $`s_r^i=(1n_r^f)q_r^i`$ and $`n_r^f=(n_r1)^2`$ where $`s_r^i=\frac{1}{2}_{\alpha ,\beta }c_{\alpha ,r}^{}\sigma _{\alpha \beta }^ic_{\beta ,r}^{}`$, with $`\sigma ^i`$ the Pauli matrices. We also define the local ”pseudospin” operators $`p_r^i=n_r^fq_r^i`$ which are the generators of the SU(2) algebra which correponds to ”rotations” between the empty and doubly occupied states.shiba72 The formal resemblance to spin symmetry has led to the name pseudospin. The total z-component of pseudospin can therefore be seen to be half the number of doubly occupied sites minus the number of empty sites, which is precisely the charge relative to half filling. Pseudospin is a symmetry of the particle-hole symmetric Hubbard model but is broken by a chemical potential.
We note that, aside from phase factors, the quasicharge operators always create even charged states from odd and vice versa, as they must since they are Fermionic. However, there is a precise phase relationship which is not so easily interpreted that is enforced by the commutation relations and symmetries.
Our Hubbard model can now be rewritten exactly in terms of the quasiparticle operators as
$$H=t(T_0+T_1+T_1)+Uh_U$$
(4)
with $`h_U=\frac{1}{2}_r\widehat{c}_r^{}\widehat{c}_r^{}`$ and
$`T_0`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{r,r^{}}{}}(1+4𝒒_r𝒒_r^{})(\widehat{c}_r^{}\widehat{c}_r^{}^{}+CC)`$ (5)
$`T_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{r,r^{}}{}}1^r(14𝒒_r𝒒_r^{})(\widehat{c}_r^{}\widehat{c}_r^{}^{})`$
and $`T_1=T_1^{}`$. In this representation, the symmetry under the entire $`SU(2)`$ quasispin rotation is manifest and we note that the Hubbard interaction becomes simply the quasicharge number operator.
We define total quasispin $`𝐐`$ , spin $`𝐒`$ and pseudospin $`𝐏`$ respectively by $`𝑸=_r𝒒_r`$ etc. $`𝑷,𝑸`$ and $`𝑺`$ all obey the SU(2) algebra $`[P_i,P_j]=i_kϵ_{ijk}P_k`$ etc. We see from the definitions above Eq. 4 the relationship, $`p_r^i=n_r^fq_r^i`$ and $`s_r^i=(1n_r^f)q_r^i`$. Thus $`𝑸=𝑷+𝑺`$, i.e. quasispin can be exactly split into pseudospin and spin, with $`n_r^f`$ and $`(1n_r^f)`$ providing the projection operator of $`𝑸`$ into either spin or pseudospin. Since $`𝑸`$ commutes with $`\widehat{c}_r^{}`$ and the Hubbard Hamiltonian $`H`$ depends on $`𝒒_r`$ only through rotationally invariant terms, we conclude that $`[H,𝑸]=0`$. We further know that the Hubbard model commutes with ordinary spin whereby $`[H,𝑺]=0`$. We can therefore conclude that $`[H,𝑷]=0`$, which confirms the well known invariance under pseudospin rotations.
The operators $`T_{\pm 1,0}`$ are identical to the operators defined in Ref. macdonald and used in a perturbation expansion in $`t/U`$. It is clear from their definitions that $`T_{\pm 1}`$ couples different Hubbard bands defined by the sectors given by different values of $`h_U`$ whereas $`T_0`$ does not. We observe $`[T_{0,\pm 1},𝑸]=0`$. It can be inferred that $`[T_{0\pm 1},𝑷]=0`$ since $`T_{0\pm 1}`$ is known to commute with ordinary spin.
We now define the usual canonical transformation $`𝒮`$ which splits the Hamiltonian into pieces that preserve total quasicharge. We define $`𝒮=i(T_1T_1)+i[T_0,(T_1+T_1)]`$ and it straightforward to show that $`[𝒮,h_U]=(T_1+T_1)`$, whereby it follows that $`H_{eff}=e^{it𝒮/U}H_{hub}e^{it𝒮/U}=`$ $`H_{hub}+`$ $`it[𝒮,H_{hub}]/U`$ $`t^2[𝒮,[𝒮,H_{hub}]]/(2U^2)+\mathrm{}`$ $`H_{q\widehat{c}}+H_q+O(t^3/U^2)`$ splits into two terms. The first term $`H_q`$ is the ordinary spin interaction in the “spin-only” sector
$$H_q=J\underset{rr^{}}{}(𝒒_r𝒒_r^{}\frac{1}{4})$$
(6)
where $`J=4t^2/U`$ and the second term $`H_{qc}`$ couples quasispin and quasicharge
$$\begin{array}{c}H_{q\widehat{c}}=\underset{rr^{}}{}(Uh_U+tT_0+\frac{t^2}{4U}(n_r^f+n_r^{}^f)(𝒒_r𝒒_r^{}\frac{1}{4}))+\hfill \\ \hfill +\frac{t^2}{4U}\underset{rr^{}r^{\prime \prime }}{}(\widehat{c}_r^{}\widehat{c}_{r^{\prime \prime }}^{}+\widehat{c}_{r^{\prime \prime }}^{}\widehat{c}_r^{})(1𝒒_r𝒒_r^{}+𝒒_r𝒒_{r^{\prime \prime }}+\\ \hfill 𝒒_r^{}𝒒_{r^{\prime \prime }}+i𝒒_r(𝒒_r^{}\times 𝒒_{r^{\prime \prime }}).)\end{array}$$
(7)
The last summation runs over all three-site neighbors connected by links. Terms of order $`t^2/U`$ which are already present to order $`t`$ or $`U`$ are ignored. According to the results of Ref. macdonald , to all orders in $`t/U`$ the operator $`𝒮`$ which diagonalizes total quasicharge is a polynomial of $`T_{0,\pm 1}`$, and hence to all orders in perturbation theory, the transformation by $`e^{it𝒮/U}`$ will preserve all the symmetries $`Q`$ and $`P`$ in this formulation.
## The spin-only sector
Under this canonical transformation the eigenstates of $`H_{eff}`$ will map to the eigenstates of the Hubbard model represented by the bare quasiparticle operators through the transformation $`e^{it𝒮/U}`$. We let $`|\mathrm{\Psi }_{eff}`$ be eigenstates of $`H_{eff}`$ and $`|\mathrm{\Psi }=e^{it𝒮/U}|\mathrm{\Psi }_{eff}`$. Adapting Ref. macdonald to our problem, there will be a “zeroth Hubbard band” or “spin-only sector” of solutions corresponding to $`n_r^f|\mathrm{\Psi }_{eff}`$ being identically zero so that $`H_{eff}`$ reduces exactly to $`H_q`$. We conclude $`H_{q\widehat{c}}_{\mathrm{\Psi }_{eff}}`$ is identically zero and the effective Hamiltonian of the zeroth subband to order $`t^2/U`$ is exactly the Heisenberg model.
What can we conclude from symmetries about the properties of the zeroth Hubbard band? It must be remembered that $`N_{tot}^f=_rn_r^f`$ is not a “good quantum number” i.e. is neither invariant under the symmetries nor under $`e^{i𝒮}`$. We have therefore no information about $`n_r^f_\mathrm{\Psi }`$ from symmetry arguments alone. However, we see that $`n_r^f|\mathrm{\Psi }_{eff}=0`$ throughout the entire lattice in the entire zeroth Hubbard subband. We can therefore conclude that $`𝑷_{\mathrm{\Psi }_{eff}}=0`$. Since $`𝑷`$ commutes with $`𝒮`$, this assertion survives to all orders in perturbation theory and we can conclude that $`𝑷_\mathrm{\Psi }=0.`$ In fact quasispin and ordinary spin coincide. Since $`n=\mathrm{\hspace{0.17em}1}2P_z`$ is equal to physical charge, we see that the spin-only sector has indeed charge density corresponding to half filling.
## Extended states in the charge-only sector
The pseudospin is given by $`𝑷^2=_{r,r^{}}𝒒_r𝒒_r^{}n_r^fn_r^{}^f`$ If there is exactly one quasicharge, we conclude from the fact that $`𝒒_r^2=3/4`$ that $`𝑷^2=\frac{3}{4}`$. We have thus identified the lowest nonzero quasicharge sector with the spin half representation of pseudospin, exactly what we need to associate it with one added positive or negative physical charge.
We can derive a set of extended charge states in the dilute limit by a mean field treatment of Eq. 7, replacing $`𝒒_r𝒒_r^{}=𝒒_r𝒒_r^{}_{H_q}`$. In two dimensions, these have the spectrum
$$\begin{array}{c}E_k=\frac{1}{2}U4t\alpha (\mathrm{cos}k_x+\mathrm{cos}k_y)+\hfill \\ \hfill \frac{t^2}{4U}((\mathrm{cos}2k_x+\mathrm{cos}2k_y)(1e_2)+\\ \hfill 2\mathrm{cos}k_x\mathrm{cos}k_y(1e_\sqrt{2}))\end{array}$$
(8)
where $`\alpha =|(\frac{1}{4}+e_1)|`$ and where $`e_1=𝒒_r𝒒_r^{}`$, $`e_\sqrt{2}`$ and $`e_2`$ denote nearest neighbor and longer range Heisenberg correlations. Using the numerical value $`e_1=0.33972(2)`$ heis we find $`\alpha =.08972`$. We find a quasicharge spectrum bounded by $`U/2\pm 8\alpha t+O(t^2/U)`$ and shown in the shaded region of Fig. 2.
## The Nagaoka instability and localized states
The present calculation makes the “Nagaoka theorem” natural.nagaoka Consider Eq. 4 in the case where the quasispins are perfectly aligned along the quasispin z axis. In this case $`𝒒_r𝒒_r^{}=\frac{1}{4}`$, the operator $`T_{\pm 1}`$ will be identically zero and the kinetic energy $`T_0`$ becomes precisely $`_{rr^{}}\widehat{c}_r^{}\widehat{c}_r^{}^{}+\widehat{c}_r^{}^{}\widehat{c}_r^{}`$. These are quasicharges moving independently of the ferromagnetic quasispins and are ordinary holes in a perfect 2-D ferromagnet which the kinetic energy $`2t(\mathrm{cos}k_xa+\mathrm{cos}k_ya)+\frac{1}{2}U`$ in the thermodynamic limit.
The Nagaoka theorem states that as $`U\mathrm{}`$ a single charge causes the ground state of $`H`$ to become ferromagnetic. This Nagaoka ferromagnet with charge one has Hubbard energy $`E_{ferro}=\frac{1}{2}U4t`$, whereas the antiferromagnetic solution with a single quasicharge has energy $`E_{antiferro}=2\mathrm{\Omega }J(e_1\frac{1}{4})+U/28\alpha t`$ where $`\mathrm{\Omega }`$ is the number of lattice sites in the system. We therefore find that $`E_{ferro}<E_{antiferro}`$ when the inequality $`U>\mathrm{\Omega }t(1+2\alpha )/(12\alpha )`$ holds. We see that whether or not there is Nagaoka ferromagnetism depends on which order we take $`U`$ and $`\mathrm{\Omega }`$ to infinity.
The argument can be used to estimate when the antiferromagnet becomes unstable to ferromagnetic ordering. (See also Ref. phasesep ). Let us consider a quasicharge density of $`\rho 1`$. We see that if $`\rho `$ obeys $`\rho <t/U(1+2\alpha )/(12\alpha )`$ the dilute ferromagnet will be energetically favored over a collection of extended quasicharge states, essentially reproducing the argument by Ioffe and Larkin.phasesep
Using the previous ideas, we shall study localized states in a 2-D antiferromagnet. As we have seen, a quasicharge hops freely on a ferromagnetic bond where $`𝒒_r𝒒_r^{}=\frac{1}{4}`$. In a mean field treatment of the quasispins, to linear order in $`t`$ a quasicharge cannot hop an a bond with antialigned quasispins since $`𝒒_r𝒒_r^{}=\frac{1}{4}`$.
Starting with a patch of locally Neel ordered quasispins oriented along $`\pm q^z`$ we consider a quasispin state with nearby spins flipped relative to the Neel state on diagonally adjacent sites of the spin-down sublattice. We note that in the presence of a quasispin ferromagnet oriented along $`+z`$, quasicharge corresponds to physical charge measured from half filling.
The core of the region will thus be a ferromagnet surrounded by a closed boundary of antialigned spins. Outside this patch, quasispins are allowed to relax, so that far away the system will be a Heisenberg antiferromagnet. The two smallest such configurations are shown in Fig. 1.
We identify the quasispin of the state by constructing the Neel state with the same rotational symmetry around the cluster. The smallest cluster (Fig.1,a) has one flipped spin thus has $`Q=\frac{1}{2}`$ whereas the cluster with two flipped spins (Fig. 1,b) has $`Q=2`$. We now add $`n_f`$ quasicharges to the cluster. According to the definition of pseudospin, in the case where quasicharge is confined to a quasiferromagnetic region with $`𝒒_r𝒒_r^{}=\frac{1}{4}+\frac{1}{2}\delta _{rr^{}}`$ , it can be shown that $`𝑷^2=\frac{n_f}{2}(\frac{n_f}{2}+1)`$. A cluster with $`n_f`$ units of quasicharge will thus have pseudospin $`P=\frac{1}{2}n_f`$ and, from the identity $`𝑸=𝑺+𝑷`$, physical spin $`S=|Q\frac{1}{2}n_f|`$. The clusters with no quasicharge have pseudospin zero and physical spin equal to quasispin, whereas the charge one cluster will have pseudospin $`P=\frac{1}{2}`$ and spin $`S=Q\frac{1}{2}`$.
The energies of these states are approximate but the spin and pseudospin quantum numbers are exact. As quasicharge is added in the cluster, quasispin, spin and pseudospin continue to be good quantum numbers. Since quasispin is conserved upon adding quasicharge each such configuration has a well defined value of quasispin independent of physical spin or charge and can be calculated in the effective theory. Quasicharge, however, is not a good quantum number, and will change upon transforming back the the physical state with $`e^{it𝒮/U}`$.
We estimate the energy by keeping only the part of the kinetic energy linear in $`t`$ and ignore the spin exchange energy beyond the antiferromagnetic boundary of the cluster. For the smallest $`Q=\frac{1}{2}`$ cluster we see sixteen bonds that are strongly affected by the localized charge. Four bonds are ferromagnetic with zero energy and twelve bonds are antiferromagnetic with energy $`J\times (\frac{1}{2})`$, i.e. they retain about 85% of the energy of the unaffected Neel state $`J\times (e_1\frac{1}{4}).59J`$.
The energy of these states is compactly written with the spectroscopic-like notation $`^QE_S^c`$ where $`Q`$ is the quasispin, $`c=2P^z`$ is the charge relative to half filling and $`S`$ is the spin. Thus $`^{{\scriptscriptstyle \frac{1}{2}}}E_{{\scriptscriptstyle \frac{1}{2}}}^0`$ is the energy of the the smallest uncharged cluster measured with respect to the Heisenberg ground state. We find $`^{{\scriptscriptstyle \frac{1}{2}}}E_{{\scriptscriptstyle \frac{1}{2}}}^0=J(12\times (\frac{1}{2})+16\times |e_1+\frac{1}{4}|)13.74t^2/U`$. The five kinetic energy eigenvalues of this cluster are $`t(2,0,0,0,2)`$. We thus find $`^{{\scriptscriptstyle \frac{1}{2}}}E_0^1=U/2+^{{\scriptscriptstyle \frac{1}{2}}}E_{{\scriptscriptstyle \frac{1}{2}}}^02tU/2+\mathrm{\hspace{0.33em}13.74}t^2/U2t`$. For the $`Q=2`$ cluster we find a kinetic energy for a single charge $`\sqrt{6}t`$; there are 16 antiferromagnetic bonds and 8 ferromagnetic bonds giving $`^2E_{\frac{3}{2}}^1=U/2+^2E_2^0\sqrt{6}tU/2+\mathrm{\hspace{0.33em}24.61}t^2/U2.45t`$. The energy of these states are plotted in Fig. 2 and compared to the energy of the band of extended states from Eq. 8.
We find that for large $`U`$ the clusters with localized charge are energetically favored over the extended states. As $`U`$ becomes larger, larger clusters lower the energy in accordance with the Nagaoka theorem.
## Conclusions
We have demonstrated an exact local mapping from a Fermionic many-body system to system of interacting Bosons and Fermions that does not introduce constraints. The mapping provides a paradigm for a non-Fermi liquid behavior for correlated Fermions valid in any dimension. Applied to the 2-D Hubbard model the transformation makes precise earlier attempts that relied on approximations and constraints to discuss spin and charge as composite operators. Our calculations supports the idea of low energy localized states where spin is bound to charge in the limit of large $`U`$. |
warning/0507/math0507217.html | ar5iv | text | # Invariant Metrics and Laplacians on Siegel-Jacobi Disk
## 1. Introduction
For a given fixed positive integer $`n`$, we let
$$_n=\{\mathrm{\Omega }^{(n,n)}|\mathrm{\Omega }=^t\mathrm{\Omega },\text{Im}\mathrm{\Omega }>0\}$$
be the Siegel upper half plane of degree $`n`$ and let
$$Sp(n,)=\{M^{(2n,2n)}|^tMJ_nM=J_n\}$$
be the symplectic group of degree $`n`$, where
$$J_n=\left(\begin{array}{cc}0& I_n\\ I_n& 0\end{array}\right).$$
We see that $`Sp(n,)`$ acts on $`_n`$ transitively by
(1.1)
$$M\mathrm{\Omega }=(A\mathrm{\Omega }+B)(C\mathrm{\Omega }+D)^1,$$
where $`M=\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)Sp(n,)`$ and $`\mathrm{\Omega }_n.`$
For two positive integers $`n`$ and $`m`$, we consider the Heisenberg group
$$H_{}^{(n,m)}=\{(\lambda ,\mu ;\kappa )|\lambda ,\mu ^{(m,n)},\kappa ^{(m,m)},\kappa +\mu ^t\lambda \text{symmetric}\}$$
endowed with the following multiplication law
$$(\lambda ,\mu ;\kappa )(\lambda ^{},\mu ^{};\kappa ^{})=(\lambda +\lambda ^{},\mu +\mu ^{};\kappa +\kappa ^{}+\lambda ^t\mu ^{}\mu ^t\lambda ^{}).$$
We define the semidirect product of $`Sp(n,)`$ and $`H_{}^{(n,m)}`$
$$G^J:=Sp(n,)H_{}^{(n,m)}$$
endowed with the following multiplication law
$$(M,(\lambda ,\mu ;\kappa ))(M^{},(\lambda ^{},\mu ^{};\kappa ^{}))=(MM^{},(\stackrel{~}{\lambda }+\lambda ^{},\stackrel{~}{\mu }+\mu ^{};\kappa +\kappa ^{}+\stackrel{~}{\lambda }^t\mu ^{}\stackrel{~}{\mu }^t\lambda ^{}))$$
with $`M,M^{}Sp(n,),(\lambda ,\mu ;\kappa ),(\lambda ^{},\mu ^{};\kappa ^{})H_{}^{(n,m)}`$ and $`(\stackrel{~}{\lambda },\stackrel{~}{\mu })=(\lambda ,\mu )M^{}`$. We call this group $`G^J`$ the Jacobi group of degree $`n`$ and index $`m`$. Then we get the natural action of $`G^J`$ on $`_n\times ^{(m,n)}`$ (cf. \[1-2\], \[7-9\], ) defined by
(1.2)
$$(M,(\lambda ,\mu ;\kappa ))(\mathrm{\Omega },Z)=(M\mathrm{\Omega },(Z+\lambda \mathrm{\Omega }+\mu )(C\mathrm{\Omega }+D)^1),$$
where $`M=\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)Sp(n,),(\lambda ,\mu ;\kappa )H_{}^{(n,m)}`$ and $`(\mathrm{\Omega },Z)_n\times ^{(m,n)}.`$ We note that the action (1.2) is transitive.
For brevity, we write $`_{n,m}:=_n\times ^{(m,n)}.`$ For a coordinate $`(\mathrm{\Omega },Z)_{n,m}`$ with $`\mathrm{\Omega }=(\omega _{\mu \nu })_n`$ and $`Z=(z_{kl})^{(m,n)},`$ we put
$`\mathrm{\Omega }`$ $`=`$ $`X+iY,X=(x_{\mu \nu }),Y=(y_{\mu \nu })\text{real},`$
$`Z`$ $`=`$ $`U+iV,U=(u_{kl}),V=(v_{kl})\text{real},`$
$`d\mathrm{\Omega }`$ $`=`$ $`(d\omega _{\mu \nu }),d\overline{\mathrm{\Omega }}=(d\overline{\omega }_{\mu \nu }),`$
$`dZ`$ $`=`$ $`(dz_{kl}),d\overline{Z}=(d\overline{z}_{kl}),`$
$`{\displaystyle \frac{}{\mathrm{\Omega }}}=\left({\displaystyle \frac{1+\delta _{\mu \nu }}{2}}{\displaystyle \frac{}{\omega _{\mu \nu }}}\right),{\displaystyle \frac{}{\overline{\mathrm{\Omega }}}}=\left({\displaystyle \frac{1+\delta _{\mu \nu }}{2}}{\displaystyle \frac{}{\overline{\omega }_{\mu \nu }}}\right),`$
$$\frac{}{Z}=\left(\begin{array}{ccc}\frac{}{z_{11}}\hfill & \mathrm{}& \frac{}{z_{m1}}\hfill \\ \mathrm{}& \mathrm{}& \mathrm{}\\ \frac{}{z_{1n}}\hfill & \mathrm{}& \frac{}{z_{mn}}\hfill \end{array}\right),\frac{}{\overline{Z}}=\left(\begin{array}{ccc}\frac{}{\overline{z}_{11}}\hfill & \mathrm{}& \frac{}{\overline{z}_{m1}}\\ \mathrm{}& \mathrm{}& \mathrm{}\\ \frac{}{\overline{z}_{1n}}& \mathrm{}& \frac{}{\overline{z}_{mn}}\hfill \end{array}\right),$$
where $`\delta _{ij}`$ denotes the Kronecker delta symbol.
C. L. Siegel introduced the symplectic metric $`ds_n^2`$ on $`_n`$ invariant under the action (1.1) of $`Sp(n,)`$ given by
(1.3) $`ds_n^2=\sigma \left(Y^1d\mathrm{\Omega }Y^1d\overline{\mathrm{\Omega }}\right)`$
and H. Maass proved that the differential operator
(1.4) $`\mathrm{\Delta }_n=\mathrm{\hspace{0.17em}4}\sigma \left(Y^t\left(Y{\displaystyle \frac{}{\overline{\mathrm{\Omega }}}}\right){\displaystyle \frac{}{\mathrm{\Omega }}}\right)`$
is the Laplacian of $`_n`$ for the symplectic metric $`ds_n^2.`$ Here $`\sigma (A)`$ denotes the trace of a square matrix $`A`$.
In , the author proved the following theorems.
Theorem A. For any two positive real numbers $`A`$ and $`B`$, the following metric
$`ds_{n,m;A,B}^2`$ $`=`$ $`A\sigma \left(Y^1d\mathrm{\Omega }Y^1d\overline{\mathrm{\Omega }}\right)`$
$`+B\{\sigma \left(Y^1{}_{}{}^{t}VVY^1d\mathrm{\Omega }Y^1d\overline{\mathrm{\Omega }}\right)+\sigma \left(Y^1{}_{}{}^{t}(dZ)d\overline{Z}\right)`$
$`\sigma \left(VY^1d\mathrm{\Omega }Y^1{}_{}{}^{t}(d\overline{Z})\right)\sigma \left(VY^1d\overline{\mathrm{\Omega }}Y^1{}_{}{}^{t}(dZ)\right)\}`$
is a Riemannian metric on $`_{n,m}`$ which is invariant under the action (1.2) of the Jacobi
group $`G^J`$.
Theorem B. For any two positive real numbers $`A`$ and $`B`$, the Laplacian $`\mathrm{\Delta }_{n,m;A,B}`$ of
$`(_{n,m},ds_{n,m;A,B}^2)`$ is given by
$`\mathrm{\Delta }_{n,m;A,B}`$ $`=`$ $`{\displaystyle \frac{4}{A}}\{\sigma \left(Y^t\left(Y{\displaystyle \frac{}{\overline{\mathrm{\Omega }}}}\right){\displaystyle \frac{}{\mathrm{\Omega }}}\right)+\sigma \left(VY^1{}_{}{}^{t}V_{}^{t}\left(Y{\displaystyle \frac{}{\overline{Z}}}\right){\displaystyle \frac{}{Z}}\right)`$
$`+\sigma \left(V^t\left(Y{\displaystyle \frac{}{\overline{\mathrm{\Omega }}}}\right){\displaystyle \frac{}{Z}}\right)+\sigma \left({}_{}{}^{t}V_{}^{t}\left(Y{\displaystyle \frac{}{\overline{Z}}}\right){\displaystyle \frac{}{\mathrm{\Omega }}}\right)\}`$
$`+{\displaystyle \frac{4}{B}}\sigma \left(Y{\displaystyle \frac{}{Z}}{}_{}{}^{t}\left({\displaystyle \frac{}{\overline{Z}}}\right)\right)`$
Let
$$G_{}=SU(n,n)Sp(n,)$$
be the symplectic group and let
$$𝔻_n=\left\{W^{(n,n)}\right|W={}_{}{}^{t}W,I_n\overline{W}W>0\}$$
be the generalized unit disk. Then $`G_{}`$ acts on $`𝔻_n`$ transitively by
$$\left(\begin{array}{cc}P& Q\\ \overline{Q}& \overline{P}\end{array}\right)W=(PW+Q)(\overline{Q}W+\overline{P})^1,$$
where $`\left(\begin{array}{cc}P& Q\\ \overline{Q}& \overline{P}\end{array}\right)G_{}`$ and $`W𝔻_n.`$ Using the Cayley transform of $`𝔻_n`$ onto $`_n`$, we can see (cf. ) that
(1.7)
$$ds_{}^2=4\sigma \left((I_nW\overline{W})^1dW(I_n\overline{W}W)^1d\overline{W}\right)$$
is a $`G_{}`$-invariant Kähler metric on $`𝔻_n`$ and H. Maass showed that its Laplacian is given by
(1.8)
$$\mathrm{\Delta }_{}=\sigma \left((I_nW\overline{W}){}_{}{}^{t}\left((I_nW\overline{W})\frac{}{\overline{W}}\right)\frac{}{W}\right).$$
Let
$$G_{}^J=\left\{(\left(\begin{array}{cc}P& Q\\ \overline{Q}& \overline{P}\end{array}\right),(\xi ,\overline{\xi };i\kappa ))\right|\left(\begin{array}{cc}P& Q\\ \overline{Q}& \overline{P}\end{array}\right)G_{},\xi ^{(m,n)},\kappa ^{(m,m)}\}$$
be the Jacobi group with the following multiplication
$`(\left(\begin{array}{cc}P& Q\\ \overline{Q}& \overline{P}\end{array}\right),(\xi ,\overline{\xi };i\kappa ))(\left(\begin{array}{cc}P^{}& Q^{}\\ \overline{Q^{}}& \overline{P^{}}\end{array}\right),(\xi ^{},\overline{\xi ^{}};i\kappa ^{}))`$
$`=`$ $`(\left(\begin{array}{cc}P& Q\\ \overline{Q}& \overline{P}\end{array}\right)\left(\begin{array}{cc}P^{}& Q^{}\\ \overline{Q^{}}& \overline{P^{}}\end{array}\right),(\stackrel{~}{\xi }+\xi ^{},\stackrel{~}{\overline{\xi ^{}}}+\overline{\xi ^{}};i\kappa +i\kappa ^{}+\stackrel{~}{\xi }{}_{}{}^{t}\overline{\xi ^{}}\stackrel{~}{\overline{\xi ^{}}}{}_{}{}^{t}\xi _{}^{})),`$
where $`\stackrel{~}{\xi }=\xi P^{}+\overline{\xi }\overline{Q^{}}`$ and $`\stackrel{~}{\overline{\xi }}=\xi Q^{}+\overline{\xi }\overline{P^{}}.`$ Then we have the canonical action of $`G_{}^J`$ on the Siegel-Jacobi disk $`𝔻_n\times ^{(m,n)}`$ (see (2.6)) given by
(1.9)
$$(\left(\begin{array}{cc}P& Q\\ \overline{Q}& \overline{P}\end{array}\right),(\xi ,\overline{\xi };i\kappa ))(W,\eta )=((PW+Q)(\overline{Q}W+\overline{P})^1,(\eta +\xi W+\overline{\xi })(\overline{Q}W+\overline{P})^1),$$
where $`W𝔻_n`$ and $`\eta ^{(m,n)}.`$
For brevity, we write $`𝔻_{n,m}:=𝔻_n\times ^{(m,n)}.`$ For a coordinate $`(W,\eta )𝔻_{n,m}`$ with $`W=(w_{\mu \nu })𝔻_n`$ and $`\eta =(\eta _{kl})^{(m,n)},`$ we put
$`dW`$ $`=`$ $`(dw_{\mu \nu }),d\overline{W}=(d\overline{w}_{\mu \nu }),`$
$`d\eta `$ $`=`$ $`(d\eta _{kl}),d\overline{\eta }=(d\overline{\eta }_{kl})`$
and
$`{\displaystyle \frac{}{W}}=\left({\displaystyle \frac{1+\delta _{\mu \nu }}{2}}{\displaystyle \frac{}{w_{\mu \nu }}}\right),{\displaystyle \frac{}{\overline{W}}}=\left({\displaystyle \frac{1+\delta _{\mu \nu }}{2}}{\displaystyle \frac{}{\overline{w}_{\mu \nu }}}\right),`$
$$\frac{}{\eta }=\left(\begin{array}{ccc}\frac{}{\eta _{11}}\hfill & \mathrm{}& \frac{}{\eta _{m1}}\hfill \\ \mathrm{}& \mathrm{}& \mathrm{}\\ \frac{}{\eta _{1n}}\hfill & \mathrm{}& \frac{}{\eta _{mn}}\hfill \end{array}\right),\frac{}{\overline{\eta }}=\left(\begin{array}{ccc}\frac{}{\overline{\eta }_{11}}\hfill & \mathrm{}& \frac{}{\overline{\eta }_{m1}}\\ \mathrm{}& \mathrm{}& \mathrm{}\\ \frac{}{\overline{\eta }_{1n}}& \mathrm{}& \frac{}{\overline{\eta }_{mn}}\hfill \end{array}\right).$$
In this paper, we find the $`G_{}^J`$-invariant Riemannian metrics on $`𝔻_{n,m}`$ and their Laplacians. In fact, we prove the following theorems.
###### Theorem 1.1.
For any two positive real numbers $`A`$ and $`B`$, the following metric $`d\stackrel{~}{s}_{n,m;A,B}^2`$ defined by
$`d\stackrel{~}{s}_{n,m;A,B}^2`$ $`=`$ $`4A\sigma \left((I_nW\overline{W})^1dW(I_n\overline{W}W)^1d\overline{W}\right)`$
$`+\mathrm{\hspace{0.17em}4}B\{\sigma \left((I_nW\overline{W})^1{}_{}{}^{t}(d\eta )d\overline{\eta }\right)`$
$`+\sigma \left((\eta \overline{W}\overline{\eta })(I_nW\overline{W})^1dW(I_n\overline{W}W)^1{}_{}{}^{t}(d\overline{\eta })\right)`$
$`+\sigma \left((\overline{\eta }W\eta )(I_n\overline{W}W)^1d\overline{W}(I_nW\overline{W})^1{}_{}{}^{t}(d\eta )\right)`$
$`\sigma \left((I_nW\overline{W})^1{}_{}{}^{t}\eta \eta (I_n\overline{W}W)^1\overline{W}dW(I_n\overline{W}W)^1d\overline{W}\right)`$
$`\sigma \left(W(I_n\overline{W}W)^1{}_{}{}^{t}\overline{\eta }\overline{\eta }(I_nW\overline{W})^1dW(I_n\overline{W}W)^1d\overline{W}\right)`$
$`+\sigma \left((I_nW\overline{W})^1{}_{}{}^{t}\eta \overline{\eta }(I_nW\overline{W})^1dW(I_n\overline{W}W)^1d\overline{W}\right)`$
$`+\sigma \left((I_n\overline{W})^1{}_{}{}^{t}\overline{\eta }\eta \overline{W}(I_nW\overline{W})^1dW(I_n\overline{W}W)^1d\overline{W}\right)`$
$`+\sigma ((I_n\overline{W})^1(I_nW)(I_n\overline{W}W)^1{}_{}{}^{t}\overline{\eta }\eta (I_n\overline{W}W)^1`$
$`\times (I_n\overline{W})(I_nW)^1dW(I_n\overline{W}W)^1d\overline{W})`$
$`\sigma ((I_nW\overline{W})^1(I_nW)(I_n\overline{W})^1{}_{}{}^{t}\overline{\eta }\eta (I_nW)^1`$
$`\times dW(I_n\overline{W}W)^1d\overline{W})\}`$
is a Riemannian metric on $`𝔻_{n,m}`$ which is invariant under the action (1.9) of the Jacobi group $`G_{}^J`$. We note that if $`n=m=1`$ and $`A=B=1,`$ we get
$`{\displaystyle \frac{1}{4}}d\stackrel{~}{s}_{1,1;1,1}^2`$ $`=`$ $`{\displaystyle \frac{dWd\overline{W}}{(1|W|^2)^2}}+{\displaystyle \frac{1}{(1|W|^2)}}d\eta d\overline{\eta }`$
$`+{\displaystyle \frac{(1+|W|^2)|\eta |^2\overline{W}\eta ^2W\overline{\eta }^2}{(1|W|^2)^3}}dWd\overline{W}`$
$`+{\displaystyle \frac{\eta \overline{W}\overline{\eta }}{(1|W|^2)^2}}dWd\overline{\eta }+{\displaystyle \frac{\overline{\eta }W\eta }{(1|W|^2)^2}}d\overline{W}d\eta .`$
###### Theorem 1.2.
For any two positive real numbers $`A`$ and $`B`$, the Laplacian $`\stackrel{~}{\mathrm{\Delta }}_{n,m;A,B}`$ of $`(𝔻_{n,m},d\stackrel{~}{s}_{n,m;A,B}^2)`$ is given by
$`\stackrel{~}{\mathrm{\Delta }}_{n,m;A,B}`$ $`=`$ $`{\displaystyle \frac{1}{A}}\{\sigma \left((I_nW\overline{W}){}_{}{}^{t}\left((I_nW\overline{W}){\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{W}}\right)`$
$`+\sigma \left({}_{}{}^{t}(\eta \overline{\eta }W){}_{}{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{W}}\right)`$
$`+\sigma \left((\overline{\eta }\eta \overline{W}){}_{}{}^{t}\left((I_nW\overline{W}){\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{\eta }}\right)`$
$`\sigma \left(\eta \overline{W}(I_nW\overline{W})^1{}_{}{}^{t}\eta {}_{}{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{\eta }}\right)`$
$`\sigma \left(\overline{\eta }W(I_n\overline{W}W)^1{}_{}{}^{t}\overline{\eta }{}_{}{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{\eta }}\right)`$
$`+\sigma \left(\overline{\eta }(I_nW\overline{W})^1{}_{}{}^{t}\eta {}_{}{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{\eta }}\right)`$
$`+\sigma \left(\eta \overline{W}W(I_n\overline{W}W)^1{}_{}{}^{t}\overline{\eta }{}_{}{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{\eta }}\right)\}`$
$`+{\displaystyle \frac{1}{B}}\sigma \left((I_n\overline{W}W){\displaystyle \frac{}{\eta }}{}_{}{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)\right).`$
We note that if $`n=m=1`$ and $`A=B=1,`$ we get
$`\stackrel{~}{\mathrm{\Delta }}_{1,1;1,1}`$ $`=`$ $`\left(1|W|^2\right)^2{\displaystyle \frac{^2}{W\overline{W}}}+\left(1|W|^2\right){\displaystyle \frac{^2}{\eta \overline{\eta }}}`$
$`+(1|W|^2)(\eta \overline{\eta }W){\displaystyle \frac{^2}{W\overline{\eta }}}+(1|W|^2)(\overline{\eta }\eta \overline{W}){\displaystyle \frac{^2}{\overline{W}\eta }}`$
$`\left(\overline{W}\eta ^2+W\overline{\eta }^2\right){\displaystyle \frac{^2}{\eta \overline{\eta }}}+\left(1+|W|^2\right)|\eta |^2{\displaystyle \frac{^2}{\eta \overline{\eta }}}.`$
The main ingredients for the proof of Theorem 1.1 and Theorem 1.2 are the partial Cayley transform, Theorem A and Theorem B. The paper is organized as follows. In Section 2, we review the partial Cayley transform that was dealt with in . In Section 3, we prove Theorem 1.1. In Section 4, we prove Theorem 1.2. In the final section we briefly remark the theory of harmonic analysis on the Siegel-Jacobi disk.
Notations : We denote by $``$ and $``$ the field of real numbers, and the field of complex numbers respectively. The symbol “:=” means that the expression on the right is the definition of that on the left. For two positive integers $`k`$ and $`l`$, $`F^{(k,l)}`$ denotes the set of all $`k\times l`$ matrices with entries in a commutative ring $`F`$. For a square matrix $`AF^{(k,k)}`$ of degree $`k`$, $`\sigma (A)`$ denotes the trace of $`A`$. For $`\mathrm{\Omega }_g,\text{Re}\mathrm{\Omega }`$ (resp. $`\text{Im}\mathrm{\Omega })`$ denotes the real (resp. imaginary) part of $`\mathrm{\Omega }.`$ For any $`MF^{(k,l)},^tM`$ denotes the transpose matrix of $`M`$. For a matrix $`AF^{(k,k)}`$ and $`BF^{(k,l)},`$ we write $`A[B]=^tBAB`$. $`I_n`$ denotes the identity matrix of degree $`n`$.
## 2. The partial Cayley transform
In this section, we review the partial Cayley transform of $`𝔻_{n,m}`$ onto $`_{n,m}`$ needed for the proof of Theorem 1.1 and Theorem 1.2.
We can identify an element $`g=(M,(\lambda ,\mu ;\kappa ))`$ of $`G^J,M=\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)Sp(n,)`$ with the element
$$\left(\begin{array}{cccc}A& 0& B& A^t\mu B^t\lambda \\ \lambda & I_m& \mu & \kappa \\ C& 0& D& C^t\mu D^t\lambda \\ 0& 0& 0& I_m\end{array}\right)$$
of $`Sp(m+n,).`$ We set
$$T_{}=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}I_{m+n}& I_{m+n}\\ iI_{m+n}& iI_{m+n}\end{array}\right).$$
We now consider the group $`G_{}^J`$ defined by
$$G_{}^J:=T_{}^1G^JT_{}.$$
If $`g=(M,(\lambda ,\mu ;\kappa ))G^J`$ with $`M=\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)Sp(n,)`$, then $`T_{}^1gT_{}`$ is given by
(2.1)
$$T_{}^1gT_{}=\left(\begin{array}{cc}P_{}& Q_{}\\ \overline{Q}_{}& \overline{P}_{}\end{array}\right),$$
where
$$P_{}=\left(\begin{array}{cc}P& \frac{1}{2}\{Q{}_{}{}^{t}(\lambda +i\mu )P{}_{}{}^{t}(\lambda i\mu )\}\\ \frac{1}{2}(\lambda +i\mu )& I_h+i\frac{\kappa }{2}\end{array}\right),$$
$$Q_{}=\left(\begin{array}{cc}Q& \frac{1}{2}\{P{}_{}{}^{t}(\lambda i\mu )Q{}_{}{}^{t}(\lambda +i\mu )\}\\ \frac{1}{2}(\lambda i\mu )& i\frac{\kappa }{2}\end{array}\right),$$
and $`P,Q`$ are given by the formulas
(2.2)
$$P=\frac{1}{2}\left\{(A+D)+i(BC)\right\}$$
and
(2.3)
$$Q=\frac{1}{2}\left\{(AD)i(B+C)\right\}.$$
From now on, we write
$$(\left(\begin{array}{cc}P& Q\\ \overline{Q}& \overline{P}\end{array}\right),(\frac{1}{2}(\lambda +i\mu ),\frac{1}{2}(\lambda i\mu );i\frac{\kappa }{2})):=\left(\begin{array}{cc}P_{}& Q_{}\\ \overline{Q}_{}& \overline{P}_{}\end{array}\right).$$
In other words, we have the relation
$$T_{}^1(\left(\begin{array}{cc}A& B\\ C& D\end{array}\right),(\lambda ,\mu ;\kappa ))T_{}=(\left(\begin{array}{cc}P& Q\\ \overline{Q}& \overline{P}\end{array}\right),(\frac{1}{2}(\lambda +i\mu ),\frac{1}{2}(\lambda i\mu );i\frac{\kappa }{2})).$$
Let
$$H_{}^{(n,m)}:=\left\{(\xi ,\eta ;\zeta )\right|\xi ,\eta ^{(m,n)},\zeta ^{(m,m)},\zeta +\eta {}_{}{}^{t}\xi \text{symmetric}\}$$
be the complex Heisenberg group endowed with the following multiplication
$$(\xi ,\eta ;\zeta )(\xi ^{},\eta ^{};\zeta ^{}):=(\xi +\xi ^{},\eta +\eta ^{};\zeta +\zeta ^{}+\xi {}_{}{}^{t}\eta _{}^{}\eta {}_{}{}^{t}\xi _{}^{})).$$
We define the semidirect product
$$SL(2n,)H_{}^{(n,m)}$$
endowed with the following multiplication
$`(\left(\begin{array}{cc}P& Q\\ R& S\end{array}\right),(\xi ,\eta ;\zeta ))(\left(\begin{array}{cc}P^{}& Q^{}\\ R^{}& S^{}\end{array}\right),(\xi ^{},\eta ^{};\zeta ^{}))`$
$`=`$ $`(\left(\begin{array}{cc}P& Q\\ R& S\end{array}\right)\left(\begin{array}{cc}P^{}& Q^{}\\ R^{}& S^{}\end{array}\right),(\stackrel{~}{\xi }+\xi ^{},\stackrel{~}{\eta }+\eta ^{};\zeta +\zeta ^{}+\stackrel{~}{\xi }{}_{}{}^{t}\eta _{}^{}\stackrel{~}{\eta }{}_{}{}^{t}\xi _{}^{})),`$
where $`\stackrel{~}{\xi }=\xi P^{}+\eta R^{}`$ and $`\stackrel{~}{\eta }=\xi Q^{}+\eta S^{}.`$
If we identify $`H_{}^{(n,m)}`$ with the subgroup
$$\left\{(\xi ,\overline{\xi };i\kappa )\right|\xi ^{(m,n)},\kappa ^{(m,m)}\}$$
of $`H_{}^{(n,m)},`$ we have the following inclusion
$$G_{}^JSU(n,n)H_{}^{(n,m)}SL(2n,)H_{}^{(n,m)}.$$
We define the mapping $`\mathrm{\Theta }:G^JG_{}^J`$ by
(2.4)
$$\mathrm{\Theta }(\left(\begin{array}{cc}A& B\\ C& D\end{array}\right),(\lambda ,\mu ;\kappa )):=(\left(\begin{array}{cc}P& Q\\ \overline{Q}& \overline{P}\end{array}\right),(\frac{1}{2}(\lambda +i\mu ),\frac{1}{2}(\lambda i\mu );i\frac{\kappa }{2})),$$
where $`P`$ and $`Q`$ are given by (2.2) and (2.3). We can see that if $`g_1,g_2G^J`$, then $`\mathrm{\Theta }(g_1g_2)=\mathrm{\Theta }(g_1)\mathrm{\Theta }(g_2).`$
According to , p. 250, $`G_{}^J`$ is of the Harish-Chandra type (cf. , p. 118). Let
$$g_{}=(\left(\begin{array}{cc}P& Q\\ \overline{Q}& \overline{P}\end{array}\right),(\lambda ,\mu ;\kappa ))$$
be an element of $`G_{}^J.`$ Since the Harish-Chandra decomposition of an element $`\left(\begin{array}{cc}P& Q\\ R& S\end{array}\right)`$ in $`SU(n,n)`$ is given by
$$\left(\begin{array}{cc}P& Q\\ R& S\end{array}\right)=\left(\begin{array}{cc}I_n& QS^1\\ 0& I_n\end{array}\right)\left(\begin{array}{cc}PQS^1R& 0\\ 0& S\end{array}\right)\left(\begin{array}{cc}I_n& 0\\ S^1R& I_n\end{array}\right),$$
the $`P_{}^+`$-component of the following element
$$g_{}(\left(\begin{array}{cc}I_n& W\\ 0& I_n\end{array}\right),(0,\eta ;0)),W𝔻_n$$
of $`SL(2n,)H_{}^{(n,m)}`$ is given by
(2.5)
$$(\left(\begin{array}{cc}I_n& (PW+Q)(\overline{Q}W+\overline{P})^1\\ 0& I_n\end{array}\right),(0,(\eta +\lambda W+\mu )(\overline{Q}W+\overline{P})^1;0)).$$
We can identify $`𝔻_{n,m}`$ with the subset
$$\left\{(\left(\begin{array}{cc}I_n& W\\ 0& I_n\end{array}\right),(0,\eta ;0))\right|W𝔻_n,\eta ^{(m,n)}\}$$
of the complexification of $`G_{}^J.`$ Indeed, $`𝔻_{n,m}`$ is embedded into $`P_{}^+`$ given by
$$P_{}^+=\left\{(\left(\begin{array}{cc}I_n& W\\ 0& I_n\end{array}\right),(0,\eta ;0))\right|W={}_{}{}^{t}W^{(n,n)},\eta ^{(m,n)}\}.$$
This is a generalization of the Harish-Chandra embedding (cf. , p. 119). Then we get the canonical transitive action of $`G_{}^J`$ on $`𝔻_{n,m}`$ defined by
(2.6)
$$(\left(\begin{array}{cc}P& Q\\ \overline{Q}& \overline{P}\end{array}\right),(\xi ,\overline{\xi };i\kappa ))(W,\eta )=((PW+Q)(\overline{Q}W+\overline{P})^1,(\eta +\xi W+\overline{\xi })(\overline{Q}W+\overline{P})^1),$$
where $`\left(\begin{array}{cc}P& Q\\ \overline{Q}& \overline{P}\end{array}\right)G_{},\xi ^{(m,n)},\kappa ^{(m,m)}`$ and $`(W,\eta )𝔻_{n,m}.`$
The author proved that the action (1.2) of $`G^J`$ on $`_{n,m}`$ is compatible with the action (2.6) of $`G_{}^J`$ on $`𝔻_{n,m}`$ through the partial Cayley transform $`\mathrm{\Phi }:𝔻_{n,m}_{n,m}`$ defined by
(2.7)
$$\mathrm{\Phi }(W,\eta ):=(i(I_n+W)(I_nW)^1,\mathrm{\hspace{0.17em}2}i\eta (I_nW)^1).$$
In other words, if $`g_0G^J`$ and $`(W,\eta )𝔻_{n,m}`$,
(2.8)
$$g_0\mathrm{\Phi }(W,\eta )=\mathrm{\Phi }(g_{}(W,\eta )),$$
where $`g_{}=T_{}^1g_0T_{}`$. $`\mathrm{\Phi }`$ is a biholomorphic mapping of $`𝔻_{n,m}`$ onto $`_{n,m}`$ which gives the partially bounded realization of $`_{n,m}`$ by $`𝔻_{n,m}`$. The inverse of $`\mathrm{\Phi }`$ is
$$\mathrm{\Phi }^1(\mathrm{\Omega },Z)=((\mathrm{\Omega }iI_n)(\mathrm{\Omega }+iI_n)^1,Z(\mathrm{\Omega }+iI_n)^1).$$
## 3. Proof of Theorem 1.1
For $`(W,\eta )𝔻_{n,m},`$ we write
$$(\mathrm{\Omega },Z):=\mathrm{\Phi }(W,\eta ).$$
Thus
(3.1)
$$\mathrm{\Omega }=i(I_n+W)(I_nW)^1,Z=2i\eta (I_nW)^1.$$
Since
$$d(I_nW)^1=(I_nW)^1dW(I_nW)^1$$
and
$$I_n+(I_n+W)(I_nW)^1=2(I_nW)^1,$$
we get the following formulas from (3.1)
(3.2) $`Y`$ $`=`$ $`{\displaystyle \frac{1}{2i}}(\mathrm{\Omega }\overline{\mathrm{\Omega }})=(I_nW)^1(I_nW\overline{W})(I_n\overline{W})^1,`$
(3.3) $`V`$ $`=`$ $`{\displaystyle \frac{1}{2i}}(Z\overline{Z})=\eta (I_nW)^1+\overline{\eta }(I_n\overline{W})^1,`$
(3.4) $`d\mathrm{\Omega }`$ $`=`$ $`2i(I_nW)^1dW(I_nW)^1,`$
(3.5) $`dZ`$ $`=`$ $`2i\left\{d\eta +\eta (I_nW)^1dW\right\}(I_nW)^1.`$
According to the formulas (3.2) and (3.4), we obtain
(3.6)
$$\sigma \left(Y^1d\mathrm{\Omega }Y^1d\overline{\mathrm{\Omega }}\right)=4\sigma \left((I_nW\overline{W})^1dW(I_n\overline{W}W)^1d\overline{W}\right).$$
From the formulas (3.2)-(3.4), we get
(3.7)
$$\sigma \left(Y^1{}_{}{}^{t}VVY^1d\mathrm{\Omega }Y^1d\overline{\mathrm{\Omega }}\right)=(a)+(b)+(c)+(d),$$
where
$`(a):`$ $`=`$ $`\mathrm{\hspace{0.17em}4}\sigma ((I_nW\overline{W})^1{}_{}{}^{t}\eta \eta (I_n\overline{W}W)^1(I_n\overline{W})(I_nW)^1`$
$`\times dW(I_n\overline{W}W)^1d\overline{W}),`$
$`(b):`$ $`=`$ $`\mathrm{\hspace{0.17em}4}\sigma \left((I_nW\overline{W})^1{}_{}{}^{t}\eta \overline{\eta }(I_nW\overline{W})^1dW(I_n\overline{W}W)^1d\overline{W}\right),`$
$`(c):`$ $`=`$ $`\mathrm{\hspace{0.17em}4}\sigma ((I_n\overline{W})^1(I_nW)(I_n\overline{W}W)^1{}_{}{}^{t}\overline{\eta }\eta (I_n\overline{W}W)^1`$
$`\times (I_n\overline{W})(I_nW)^1dW(I_n\overline{W}W)^1d\overline{W}),`$
$`(d):`$ $`=`$ $`\mathrm{\hspace{0.17em}4}\sigma ((I_n\overline{W})^1(I_nW)(I_n\overline{W}W)^1{}_{}{}^{t}\overline{\eta }\overline{\eta }(I_nW\overline{W})^1`$
$`\times dW(I_n\overline{W}W)^1d\overline{W}).`$
According to the formulas (3.2) and (3.5), we get
(3.8)
$$\sigma \left(Y^1{}_{}{}^{t}(dZ)d\overline{Z}\right)=(e)+(f)+(g)=(h),$$
where
$`(e):`$ $`=`$ $`\mathrm{\hspace{0.17em}4}\sigma \left((I_nW\overline{W})^1{}_{}{}^{t}(d\eta )d\overline{\eta }\right),`$
$`(f):`$ $`=`$ $`\mathrm{\hspace{0.17em}4}\sigma \left((I_nW\overline{W})^1dW(I_nW)^1{}_{}{}^{t}\eta d\overline{\eta }\right),`$
$`(g):`$ $`=`$ $`\mathrm{\hspace{0.17em}4}\sigma \left((I_nW\overline{W})^1{}_{}{}^{t}(d\eta )\overline{\eta }(I_n\overline{W})^1d\overline{W}\right),`$
$`(h):`$ $`=`$ $`\mathrm{\hspace{0.17em}4}\sigma \left((I_nW\overline{W})^1dW(I_nW)^1{}_{}{}^{t}\eta \overline{\eta }(I_n\overline{W})^1d\overline{W}\right).`$
From the formulas (3.2)-(3.5), we get
(3.9)
$$\sigma \left(VY^1d\mathrm{\Omega }Y^1{}_{}{}^{t}(d\overline{Z})\right)=(i)+(j)+(k)+(l),$$
where
$`(i):`$ $`=`$ $`4\sigma \left(\eta (I_nW)^1(I_n\overline{W})(I_nW\overline{W})^1dW(I_n\overline{W}W)^1{}_{}{}^{t}(d\overline{\eta })\right),`$
$`(j):`$ $`=`$ $`4\sigma ((I_n\overline{W})^1{}_{}{}^{t}\overline{\eta }\eta (I_nW)^1(I_n\overline{W})(I_nW\overline{W})^1`$
$`\times dW(I_n\overline{W}W)^1d\overline{W}),`$
$`(k):`$ $`=`$ $`4\sigma \left(\overline{\eta }(I_nW\overline{W})^1dW(I_n\overline{W}W)^1{}_{}{}^{t}(d\overline{\eta })\right),`$
$`(l):`$ $`=`$ $`4\sigma \left((I_n\overline{W})^1{}_{}{}^{t}\overline{\eta }\overline{\eta }(I_nW\overline{W})^1dW(I_n\overline{W}W)^1d\overline{W}\right).`$
Conjugating the formula (3.9), we get
(3.10)
$$\sigma \left(VY^1d\overline{\mathrm{\Omega }}Y^1{}_{}{}^{t}(dZ)\right)=(m)+(n)+(o)+(p),$$
where
$`(m):`$ $`=`$ $`4\sigma \left(\overline{\eta }(I_n\overline{W})^1(I_nW)(I_n\overline{W}W)^1d\overline{W}(I_nW\overline{W})^1{}_{}{}^{t}(d\eta )\right),`$
$`(n):`$ $`=`$ $`4\sigma ((I_nW)^1{}_{}{}^{t}\eta \overline{\eta }(I_n\overline{W})^1(I_nW)(I_n\overline{W}W)^1`$
$`\times d\overline{W}(I_nW\overline{W})^1dW),`$
$`(o):`$ $`=`$ $`4\sigma \left(\eta (I_n\overline{W}W)^1d\overline{W}(I_nW\overline{W})^1{}_{}{}^{t}(d\eta )\right),`$
$`(p):`$ $`=`$ $`4\sigma \left((I_nW)^1{}_{}{}^{t}\eta \eta (I_n\overline{W}W)^1d\overline{W}(I_nW\overline{W})^1dW\right).`$
If we add the formulas $`(f),(i)`$ and $`(k)`$, we get
(3.11)
$$(f)+(i)+(k)=\mathrm{\hspace{0.17em}4}\sigma \left((\eta \overline{W}\overline{\eta })(I_nW\overline{W})^1dW(I_n\overline{W}W)^1{}_{}{}^{t}(d\overline{\eta })\right).$$
Indeed, transposing the matrix inside the formula $`(f)`$, we can express the formula $`(f)`$ as
$$(f)=\mathrm{\hspace{0.17em}4}\sigma \left(\eta (I_nW)^1dW(I_n\overline{W}W)^1{}_{}{}^{t}(d\overline{\eta })\right)$$
and adding the formulas $`(f)`$ and $`(i)`$ together with the formula $`(k)`$, we get the formula (3.11) because
$`(I_nW)^1(I_nW)^1(I_n\overline{W})(I_nW\overline{W})^1`$
$`=`$ $`(I_nW)^1\left\{(I_nW\overline{W})(I_n\overline{W})\right\}(I_nW\overline{W})^1`$
$`=`$ $`(I_nW)^1(I_nW)\overline{W}(I_nW\overline{W})^1`$
$`=`$ $`\overline{W}(I_nW\overline{W})^1.`$
If we add the formulas $`(g),(m)`$ and $`(o)`$, we get
(3.12)
$$(g)+(m)+(o)=\mathrm{\hspace{0.17em}4}\sigma \left((\overline{\eta }W\eta )(I_n\overline{W}W)^1d\overline{W}(I_nW\overline{W})^1{}_{}{}^{t}(d\eta )\right).$$
Indeed, we can express the formula $`(g)`$ as
$$(g)=\mathrm{\hspace{0.17em}4}\sigma \left(\overline{\eta }(I_n\overline{W})^1d\overline{W}(I_nW\overline{W})^1{}_{}{}^{t}(d\eta )\right)$$
and adding the formulas $`(g)`$ and $`(m)`$ together with the formula $`(o)`$, we get the formula (3.12) because
$`(I_n\overline{W})^1(I_n\overline{W})^1(I_nW)(I_n\overline{W}W)^1`$
$`=`$ $`(I_n\overline{W})^1\left\{(I_n\overline{W}W)(I_nW)\right\}(I_n\overline{W}W)^1`$
$`=`$ $`(I_n\overline{W})^1(I_n\overline{W})W(I_n\overline{W}W)^1`$
$`=`$ $`W(I_n\overline{W}W)^1.`$
If we add the formulas $`(a)`$ and $`(p)`$, we get
$`(a)+(p)`$ $`=`$ $`4\sigma ((I_nW\overline{W})^1{}_{}{}^{t}\eta \eta (I_n\overline{W}W)^1\overline{W}`$
$`\times dW(I_n\overline{W}W)^1d\overline{W}).`$
Indeed, transposing the matrix inside the formula $`(p)`$, we can express the formula $`(p)`$ as
$$(p)=4\sigma \left((I_nW\overline{W})^1{}_{}{}^{t}\eta \eta (I_nW)^1dW(I_n\overline{W}W)^1d\overline{W}\right)$$
and adding the formulas $`(a)`$ and $`(p)`$, we get the formula (3.13) because
$`(I_n\overline{W}W)^1(I_n\overline{W})(I_nW)^1(I_nW)^1`$
$`=`$ $`(I_n\overline{W}W)^1\left\{(I_n\overline{W})(I_n\overline{W}W)\right\}(I_nW)^1`$
$`=`$ $`(I_n\overline{W}W)^1(\overline{W})(I_nW)(I_nW)^1`$
$`=`$ $`(I_n\overline{W}W)^1\overline{W}.`$
Adding the formulas $`(d)`$ and $`(l)`$, we get
$`(d)+(l)`$ $`=`$ $`4\sigma (W(I_n\overline{W}W)^1{}_{}{}^{t}\overline{\eta }\overline{\eta }(I_nW\overline{W})^1`$
$`\times dW(I_n\overline{W}W)^1d\overline{W})`$
because
$`(I_n\overline{W})^1(I_nW)(I_n\overline{W}W)^1(I_n\overline{W})^1`$
$`=`$ $`(I_n\overline{W})^1\left\{(I_nW)(I_n\overline{W}W)\right\}(I_n\overline{W}W)^1`$
$`=`$ $`(I_n\overline{W})^1(I_n\overline{W})(W)(I_n\overline{W}W)^1`$
$`=`$ $`W(I_n\overline{W}W)^1.`$
Adding the formulas $`(h)`$ and $`(j)`$, we get
$`(h)+(j)`$ $`=`$ $`\mathrm{\hspace{0.17em}4}\sigma ((I_n\overline{W})^1{}_{}{}^{t}\overline{\eta }\eta \overline{W}(I_nW\overline{W})^1`$
$`\times dW(I_n\overline{W}W)^1d\overline{W}).`$
Indeed, transposing the matrix inside the formula $`(h)`$, we can express the formula $`(h)`$ as
$$(h)=\mathrm{\hspace{0.17em}4}\sigma \left((I_n\overline{W})^1{}_{}{}^{t}\overline{\eta }\eta (I_nW)^1dW(I_n\overline{W}W)^1d\overline{W}\right)$$
and adding the formulas $`(h)`$ and $`(j)`$, we get the formula (3.15) because
$`(I_nW)^1(I_nW)^1(I_n\overline{W})(I_nW\overline{W})^1`$
$`=`$ $`(I_nW)^1\left\{(I_nW\overline{W})(I_n\overline{W})\right\}(I_nW\overline{W})^1`$
$`=`$ $`(I_nW)^1(I_nW)\overline{W}(I_nW\overline{W})^1`$
$`=`$ $`\overline{W}(I_nW\overline{W})^1.`$
Transposing the matrix inside the formula $`(n)`$, we get
$`(n)`$ $`=`$ $`4\sigma ((I_nW\overline{W})^1(I_nW)(I_n\overline{W})^1{}_{}{}^{t}\overline{\eta }\eta (I_nW)^1`$
$`\times dW(I_n\overline{W}W)^1d\overline{W}).`$
From the formulas (3.7)-(3.16), we obtain
$`\sigma \left(Y^1{}_{}{}^{t}VVY^1d\mathrm{\Omega }Y^1d\overline{\mathrm{\Omega }}\right)+\sigma \left(Y^1{}_{}{}^{t}(dZ)d\overline{Z}\right)`$
$`\sigma \left(VY^1d\mathrm{\Omega }Y^1{}_{}{}^{t}(d\overline{Z})\right)\sigma \left(VY^1d\overline{\mathrm{\Omega }}Y^1{}_{}{}^{t}(dZ)\right)`$
$`=(a)+(b)+(c)+(d)+\mathrm{}+(m)+(n)+(o)+(p)`$
$`=4\sigma \left((I_nW\overline{W})^1{}_{}{}^{t}(d\eta )d\overline{\eta }\right)`$
$`+\mathrm{\hspace{0.17em}4}\sigma \left((\eta \overline{W}\overline{\eta })(I_nW\overline{W})^1dW(I_n\overline{W}W)^1{}_{}{}^{t}(d\overline{\eta })\right)`$
$`+\mathrm{\hspace{0.17em}4}\sigma \left((\overline{\eta }W\eta )(I_n\overline{W}W)^1d\overline{W}(I_nW\overline{W})^1{}_{}{}^{t}(d\eta )\right)`$
$`\mathrm{\hspace{0.17em}4}\sigma \left((I_nW\overline{W})^1{}_{}{}^{t}\eta \eta (I_n\overline{W}W)^1\overline{W}dW(I_n\overline{W}W)^1d\overline{W}\right)`$
$`\mathrm{\hspace{0.17em}4}\sigma \left(W(I_n\overline{W}W)^1{}_{}{}^{t}\overline{\eta }\overline{\eta }(I_nW\overline{W})^1dW(I_n\overline{W}W)^1d\overline{W}\right)`$
$`+\mathrm{\hspace{0.17em}4}\sigma \left((I_nW\overline{W})^1{}_{}{}^{t}\eta \overline{\eta }(I_nW\overline{W})^1dW(I_n\overline{W}W)^1d\overline{W}\right)`$
$`+\mathrm{\hspace{0.17em}4}\sigma \left((I_n\overline{W})^1{}_{}{}^{t}\overline{\eta }\eta \overline{W}(I_nW\overline{W})^1dW(I_n\overline{W}W)^1d\overline{W}\right)`$
$`+\mathrm{\hspace{0.17em}4}\sigma ((I_n\overline{W})^1(I_nW)(I_n\overline{W}W)^1{}_{}{}^{t}\overline{\eta }\eta (I_n\overline{W}W)^1`$
$`\times (I_n\overline{W})(I_nW)^1dW(I_n\overline{W}W)^1d\overline{W})`$
$`\mathrm{\hspace{0.17em}4}\sigma ((I_nW\overline{W})^1(I_nW)(I_n\overline{W})^1{}_{}{}^{t}\overline{\eta }\eta (I_nW)^1`$
$`\times dW(I_n\overline{W}W)^1d\overline{W}).`$
Consequently the complete proof follows from the above formula, the formula (3.6), Theorem A and the fact that the action (1.2) of $`G^J`$ on $`_{n,m}`$ is compatible with the action (2.6) of $`G_{}^J`$ on $`𝔻_{n,m}`$ through the partial Cayley transform. $`\mathrm{}`$
## 4. Proof of Theorem 1.2
From the formulas (3.1), (3.4) and (3.5), we get
(4.1)
$$\frac{}{\mathrm{\Omega }}=\frac{1}{2i}(I_nW)[{}_{}{}^{t}\left\{(I_nW)\frac{}{W}\right\}{}_{}{}^{t}\left\{{}_{}{}^{t}\eta {}_{}{}^{t}\left(\frac{}{\eta }\right)\right\}]$$
and
(4.2)
$$\frac{}{Z}=\frac{1}{2i}(I_nW)\frac{}{\eta }.$$
We need the following lemma for the proof of Theorem 1.2. H. Maass observed the following useful fact. Lemma 4.1. (a) Let $`A`$ be an $`m\times n`$ matrix and $`B`$ an $`n\times l`$ matrix. Assume that the entries of $`A`$ commute with the entries of $`B`$. Then $`{}_{}{}^{t}(AB)={}_{}{}^{t}B{}_{}{}^{t}A.`$ (b) Let $`A,B`$ and $`C`$ be a $`k\times l`$, an $`n\times m`$ and an $`m\times l`$ matrix respectively. Assume that the entries of $`A`$ commute with the entries of $`B`$. Then
$`{}_{}{}^{t}(A{}_{}{}^{t}(BC))=B{}_{}{}^{t}(A^tC).`$
Proof. The proof follows immediately from the direct computation. $`\mathrm{}`$
From the formulas (3.2), (4.1) and Lemma 4.1, we get the following formula
(4.3)
$$4\sigma \left(Y^t\left(Y\frac{}{\overline{\mathrm{\Omega }}}\right)\frac{}{\mathrm{\Omega }}\right)=(\alpha )+(\beta )+(\gamma )+(\delta ),$$
where
$`(\alpha ):`$ $`=`$ $`\sigma \left((I_nW\overline{W}){}_{}{}^{t}\left((I_nW\overline{W}){\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{W}}\right),`$
$`(\beta ):`$ $`=`$ $`\sigma \left(\eta (I_nW)^1(I_nW\overline{W})^t\left((I_nW\overline{W}){\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{\eta }}\right),`$
$`(\gamma ):`$ $`=`$ $`\sigma \left((I_nW\overline{W})(I_n\overline{W})^1{}_{}{}^{t}\overline{\eta }{}_{}{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{W}}\right),`$
$`(\delta ):`$ $`=`$ $`\sigma \left(\eta (I_nW)^1(I_nW\overline{W})(I_n\overline{W})^1{}_{}{}^{t}\overline{\eta }_{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{\eta }}\right).`$
According to the formulas (3.2) and (4.2), we get
(4.4)
$$4\sigma \left(Y\frac{}{Z}^t\left(\frac{}{\overline{Z}}\right)\right)=\sigma \left((I_n\overline{W}W)\frac{}{\eta }^t\left(\frac{}{\overline{\eta }}\right)\right).$$
From the formulas (3.2), (3.3) and (4.2),we get
(4.5)
$$4\sigma \left(VY^1{}_{}{}^{t}V_{}^{t}\left(Y\frac{}{\overline{Z}}\right)\frac{}{Z}\right)=(ϵ)+(\zeta )+(\eta )+(\theta ),$$
where
$`(ϵ):`$ $`=`$ $`\sigma \left(\eta (I_n\overline{W}W)^1{}_{}{}^{t}\overline{\eta }_{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{\eta }}\right),`$
$`(\zeta ):`$ $`=`$ $`\sigma \left(\overline{\eta }(I_nW\overline{W})^1{}_{}{}^{t}\eta _{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{\eta }}\right),`$
$`(\eta ):`$ $`=`$ $`\sigma \left(\eta (I_nW)^1(I_n\overline{W})(I_nW\overline{W})^1{}_{}{}^{t}\eta _{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{\eta }}\right),`$
$`(\theta ):`$ $`=`$ $`\sigma \left(\overline{\eta }(I_n\overline{W})^1(I_nW)(I_n\overline{W}W)^1{}_{}{}^{t}\overline{\eta }_{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{\eta }}\right).`$
Using the formulas (3.2), (3.3), (4.1), (4.2) and Lemma 4.1, we get
(4.6)
$$4\sigma \left(V^t\left(Y\frac{}{\overline{\mathrm{\Omega }}}\right)\frac{}{Z}\right)=(\iota )+(\kappa )+(\lambda )+(\mu ),$$
where
$`(\iota ):`$ $`=`$ $`\sigma \left(\overline{\eta }^t\left((I_nW\overline{W}){\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{\eta }}\right),`$
$`(\kappa ):`$ $`=`$ $`\sigma \left(\eta (I_nW)^1(I_n\overline{W})^t\left((I_nW\overline{W}){\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{\eta }}\right),`$
$`(\lambda ):`$ $`=`$ $`\sigma \left(\eta (I_nW)^1{}_{}{}^{t}\overline{\eta }_{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{\eta }}\right),`$
$`(\mu ):`$ $`=`$ $`\sigma \left(\overline{\eta }(I_n\overline{W})^1{}_{}{}^{t}\overline{\eta }_{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{\eta }}\right).`$
Using the formulas (3.2), (3.3), (4.1), (4.2) and Lemma 4.1, we get
(4.7)
$$4\sigma \left({}_{}{}^{t}V_{}^{t}\left(Y\frac{}{\overline{Z}}\right)\frac{}{\mathrm{\Omega }}\right)=(\nu )+(\xi )+(o)+(\pi ),$$
where
$`(\nu ):`$ $`=`$ $`\sigma \left({}_{}{}^{t}\eta _{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{W}}\right),`$
$`(\xi ):`$ $`=`$ $`\sigma \left((I_nW)(I_n\overline{W})^1{}_{}{}^{t}\overline{\eta }_{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{W}}\right),`$
$`(o):`$ $`=`$ $`\sigma \left(\eta (I_nW)^1{}_{}{}^{t}\eta _{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{\eta }}\right),`$
$`(\pi ):`$ $`=`$ $`\sigma \left(\eta (I_n\overline{W})^1{}_{}{}^{t}\overline{\eta }_{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{\eta }}\right).`$
Adding the formulas $`(\gamma ),(\nu )`$ and $`(\xi )`$, we get
(4.8)
$$(\gamma )+(\nu )+(\xi )=\sigma \left({}_{}{}^{t}(\eta \overline{\eta }W)_{}^{t}\left(\frac{}{\overline{\eta }}\right)(I_n\overline{W}W)\frac{}{W}\right)$$
because
$`(I_nW\overline{W})(I_n\overline{W})^1+(I_nW)(I_n\overline{W})^1`$
$`=`$ $`W(I_n\overline{W})(I_n\overline{W})^1=W.`$
Adding the formulas $`(\beta ),(\iota )`$ and $`(\kappa )`$, we get
(4.9)
$$(\beta )+(\iota )+(\kappa )=\sigma \left((\overline{\eta }\eta \overline{W})^t\left((I_nW\overline{W})\frac{}{\overline{W}}\right)\frac{}{\eta }\right)$$
because
$`(I_nW)^1(I_nW\overline{W})+(I_nW)^1(I_n\overline{W})`$
$`=`$ $`(I_nW)^1(I_nW)\overline{W}=\overline{W}.`$
If we add the formulas $`(\eta )`$ and $`(o)`$, we get
(4.10)
$$(\eta )+(o)=\sigma \left(\eta \overline{W}(I_nW\overline{W})^1{}_{}{}^{t}\eta _{}^{t}\left(\frac{}{\overline{\eta }}\right)(I_n\overline{W}W)\frac{}{\eta }\right)$$
because
$`(I_nW)^1(I_n\overline{W})(I_nW\overline{W})^1(I_nW)^1`$
$`=`$ $`(I_nW)^1\left\{I_n\overline{W}(I_nW\overline{W})\right\}(I_nW\overline{W})^1`$
$`=`$ $`(I_nW)^1(I_nW)(\overline{W})(I_nW\overline{W})^1`$
$`=`$ $`\overline{W}(I_nW\overline{W})^1.`$
If we add the formulas $`(\theta )`$ and $`(\mu )`$, we get
(4.11)
$$(\theta )+(\mu )=\sigma \left(\overline{\eta }W(I_n\overline{W}W)^1{}_{}{}^{t}\overline{\eta }_{}^{t}\left(\frac{}{\overline{\eta }}\right)(I_n\overline{W}W)\frac{}{\eta }\right)$$
because
$`(I_n\overline{W})^1(I_nW)(I_n\overline{W}W)^1(I_n\overline{W})^1`$
$`=`$ $`(I_n\overline{W})^1\left\{I_nW(I_n\overline{W}W)\right\}(I_n\overline{W}W)^1`$
$`=`$ $`(I_n\overline{W})^1(I_n\overline{W})(W)(I_n\overline{W}W)^1`$
$`=`$ $`W(I_n\overline{W}W)^1.`$
If we add the formulas $`(\delta ),(ϵ),(\lambda )`$ and $`(\pi )`$, we get
(4.12)
$$(\delta )+(ϵ)+(\lambda )+(\pi )=\sigma \left(\eta \overline{W}W(I_n\overline{W}W)^1{}_{}{}^{t}\overline{\eta }_{}^{t}\left(\frac{}{\overline{\eta }}\right)(I_n\overline{W}W)\frac{}{\eta }\right)$$
because
$`(I_nW)^1(I_nW\overline{W})(I_n\overline{W})^1+(I_n\overline{W}W)^1`$
$`(I_nW)^1(I_n\overline{W})^1`$
$`=`$ $`(I_nW)^1\left\{(I_nW\overline{W})(I_n\overline{W})\right\}(I_n\overline{W})^1`$
$`+(I_n\overline{W}W)^1(I_n\overline{W})^1`$
$`=`$ $`\overline{W}(I_n\overline{W})^1+(I_n\overline{W}W)^1(I_n\overline{W})^1`$
$`=`$ $`(I_n\overline{W})(I_n\overline{W})^1+(I_n\overline{W}W)^1`$
$`=`$ $`I_n+(I_n\overline{W}W)^1`$
$`=`$ $`\left\{(I_n\overline{W}W)+I_n\right\}(I_n\overline{W}W)^1`$
$`=`$ $`\overline{W}W(I_n\overline{W}W)^1.`$
From the formulas (4.3) and (4.5)-(4.12), we obtain
$`\sigma \left(Y^t\left(Y{\displaystyle \frac{}{\overline{\mathrm{\Omega }}}}\right){\displaystyle \frac{}{\mathrm{\Omega }}}\right)+\sigma \left(VY^1{}_{}{}^{t}V_{}^{t}\left(Y{\displaystyle \frac{}{\overline{Z}}}\right){\displaystyle \frac{}{Z}}\right)`$
$`+\sigma \left(V^t\left(Y{\displaystyle \frac{}{\overline{\mathrm{\Omega }}}}\right){\displaystyle \frac{}{Z}}\right)+\sigma \left({}_{}{}^{t}V_{}^{t}\left(Y{\displaystyle \frac{}{\overline{Z}}}\right){\displaystyle \frac{}{\mathrm{\Omega }}}\right)`$
$`=(\alpha )+(\beta )+(\gamma )+(\delta )+\mathrm{}+(\nu )+(\xi )+(o)+(\pi )`$
$`=\sigma \left((I_nW\overline{W}){}_{}{}^{t}\left((I_nW\overline{W}){\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{W}}\right)`$
$`+\sigma \left({}_{}{}^{t}(\eta \overline{\eta }W){}_{}{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{W}}\right)`$
$`+\sigma \left((\overline{\eta }\eta \overline{W}){}_{}{}^{t}\left((I_nW\overline{W}){\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{\eta }}\right)`$
$`\sigma \left(\eta \overline{W}(I_nW\overline{W})^1{}_{}{}^{t}\eta {}_{}{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{\eta }}\right)`$
$`\sigma \left(\overline{\eta }W(I_n\overline{W}W)^1{}_{}{}^{t}\overline{\eta }{}_{}{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{\eta }}\right)`$
$`+\sigma \left(\overline{\eta }(I_nW\overline{W})^1{}_{}{}^{t}\eta {}_{}{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{\eta }}\right)`$
$`+\sigma \left(\eta \overline{W}W(I_n\overline{W}W)^1{}_{}{}^{t}\overline{\eta }{}_{}{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{\eta }}\right).`$
Consequently the complete proof follows from the formula (4.4), the above formula, Theorem B and the fact that the action (1.2) of $`G^J`$ on $`_{n,m}`$ is compatible with the action (2.6) of $`G_{}^J`$ on $`𝔻_{n,m}`$ through the partial Cayley transform. $`\mathrm{}`$
Remark 4.1. We proved in that the following two differential operators $`D`$ and $`L:=\frac{1}{4}\mathrm{\Delta }_{n,m;1,1}D`$ on $`_{n,m}`$ defined by
$$D=\sigma \left(Y\frac{}{Z}{}_{}{}^{t}\left(\frac{}{\overline{Z}}\right)\right)$$
and
$`L`$ $`=`$ $`\sigma \left(Y^t\left(Y{\displaystyle \frac{}{\overline{\mathrm{\Omega }}}}\right){\displaystyle \frac{}{\mathrm{\Omega }}}\right)+\sigma \left(VY^1{}_{}{}^{t}V_{}^{t}\left(Y{\displaystyle \frac{}{\overline{Z}}}\right){\displaystyle \frac{}{Z}}\right)`$
$`+\sigma \left(V^t\left(Y{\displaystyle \frac{}{\overline{\mathrm{\Omega }}}}\right){\displaystyle \frac{}{Z}}\right)+\sigma \left({}_{}{}^{t}V_{}^{t}\left(Y{\displaystyle \frac{}{\overline{Z}}}\right){\displaystyle \frac{}{\mathrm{\Omega }}}\right)`$
are invariant under the action (1.2) of $`G^J`$. By the formula (4.4) and the proof of Theorem 1.2, we see that the following differential operators $`\stackrel{~}{D}`$ and $`\stackrel{~}{L}:=\stackrel{~}{\mathrm{\Delta }}_{n,m;1,1}`$ on $`𝔻_{n,m}`$ defined by
$$\stackrel{~}{D}=\sigma \left((I_n\overline{W}W)\frac{}{\eta }^t\left(\frac{}{\overline{\eta }}\right)\right)$$
and
$`\stackrel{~}{L}`$ $`=`$ $`\sigma \left((I_nW\overline{W}){}_{}{}^{t}\left((I_nW\overline{W}){\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{W}}\right)`$
$`+\sigma \left({}_{}{}^{t}(\eta \overline{\eta }W){}_{}{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{W}}\right)`$
$`+\sigma \left((\overline{\eta }\eta \overline{W}){}_{}{}^{t}\left((I_nW\overline{W}){\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{\eta }}\right)`$
$`\sigma \left(\eta \overline{W}(I_nW\overline{W})^1{}_{}{}^{t}\eta {}_{}{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{\eta }}\right)`$
$`\sigma \left(\overline{\eta }W(I_n\overline{W}W)^1{}_{}{}^{t}\overline{\eta }{}_{}{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{\eta }}\right)`$
$`+\sigma \left(\overline{\eta }(I_nW\overline{W})^1{}_{}{}^{t}\eta {}_{}{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{\eta }}\right)`$
$`+\sigma \left(\eta \overline{W}W(I_n\overline{W}W)^1{}_{}{}^{t}\overline{\eta }{}_{}{}^{t}\left({\displaystyle \frac{}{\overline{\eta }}}\right)(I_n\overline{W}W){\displaystyle \frac{}{\eta }}\right)`$
are invariant under the action (2.6) of $`G_{}^J.`$ Indeed it is very complicated and difficult at this moment to express the generators of the algebra of all $`G_{}^J`$-invariant differential operators on $`𝔻_{n,m}`$ explicitly. We propose an open problem to find other explicit $`G_{}^J`$-invariant differential operators on $`𝔻_{n,m}`$.
## 5. Remark on Harmonic Analysis on Siegel-Jacobi Disk
It might be interesting to develop the theory of harmonic analysis on the Siegel-Jacobi disk $`𝔻_{n,m}`$. The theory of harmonic analysis on the generalized unit disk $`𝔻_n`$ can be done explicitly by the work of Harish-Chandra because $`𝔻_n`$ is a symmetric space. However the Siegel-Jacobi disk $`𝔻_{n,m}`$ is not a symmetric space. The work for developing the theory of harmonic analysis on $`𝔻_{n,m}`$ explicitly is complicated and difficult at this moment. We observe that this work on $`𝔻_{n,m}`$ generalizes the work on the generalized unit disk $`𝔻_n`$.
More precisely, if we put $`G_{}=SU(n,n)Sp(n,)`$, then the Jacobi group
$$G_{}^J=\left\{(\left(\begin{array}{cc}P& Q\\ \overline{Q}& \overline{P}\end{array}\right),(\xi ,\overline{\xi };i\kappa ))\right|\left(\begin{array}{cc}P& Q\\ \overline{Q}& \overline{P}\end{array}\right)G_{},\xi ^{(m,n)},\kappa ^{(m,m)}\}$$
acts on the Siegel-Jacobi disk $`𝔻_{n,m}`$ transitively via the transformation behavior (1.9). It is easily seen that the stabilizer $`K_{}^J`$ of the action (1.9) at the base point $`(0,0)`$ is given by
$$K_{}^J=\left\{(\left(\begin{array}{cc}P& 0\\ 0& \overline{P}\end{array}\right),(0,0;i\kappa ))\right|PU(n),\kappa ^{(m,m)}\}.$$
Therefore $`G_{}^J/K_{}^J`$ is biholomorphic to $`𝔻_{n,m}`$ via the correspondence
$$gK_{}^Jg(0,0),gG_{}^J.$$
We observe that the Siegel-Jacobi disk $`𝔻_{n,m}`$ is not a reductive symmetric space.
We let
$$\mathrm{\Gamma }_{n,m}:=Sp(n,)H_{}^{(n,m)},$$
where $`Sp(n,)`$ is the Siegel modular group of degree $`n`$ and
$$H_{}^{(n,m)}=\left\{(\lambda ,\mu ;\kappa )H_{}^{(n,m)}\right|\lambda ,\mu ,\kappa \text{are integral}\}.$$
We set
$$\mathrm{\Gamma }_{n,m}^{}:=T_{}^1\mathrm{\Gamma }_{n,m}T_{},$$
where $`T_{}`$ was already defined in Section 2. Clearly the arithmetic subgroup $`\mathrm{\Gamma }_{n,m}^{}`$ acts on $`𝔻_{n,m}`$ properly continuously. We can describe a fundamental domain $`_{n,m}^{}`$ for $`\mathrm{\Gamma }_{n,m}^{}\backslash 𝔻_{n,m}`$ explicitly using the partial Cayley transform and a fundamental domain $`_{n,m}`$ for $`\mathrm{\Gamma }_{n,m}\backslash _{n,m}`$ which is described explicitly in . The $`G_{}^J`$-invariant metric $`d\stackrel{~}{s}_{n,m;A,B}`$ on $`𝔻_{n,m}`$ induces a metric on $`_{n,m}^{}`$ naturally. It may be intersting to investigate the spectral theory of the Laplacian $`\stackrel{~}{\mathrm{\Delta }}_{n,m;A,B}`$ on a fundamental domain $`_{n,m}^{}`$. But this work is very complicated and difficult at this moment.
For instance, we consider the case $`n=m=1`$ and $`A=B=1`$. In this case
$$G_{}^J=\left\{(\left(\begin{array}{cc}p& q\\ \overline{q}& \overline{p}\end{array}\right),(\xi ,\overline{\xi };i\kappa ))\right|p,q,\xi ,|p|^2|q|^2=1,\kappa \}$$
and
$$K_{}^J=\left\{(\left(\begin{array}{cc}p& 0\\ 0& \overline{p}\end{array}\right),(0,0;i\kappa ))\right|p,|p|=1,\kappa \}.$$
$`d\stackrel{~}{s}_{1,1;1,1}`$ is a $`G_{}^J`$-invariant Riemannian metric on $`𝔻_{1,1}=𝔻_1\times `$ (cf. Theorem 1.1) and $`\stackrel{~}{\mathrm{\Delta }}_{1,1;1,1}`$ is its Laplacian. It is well known that the theory of harmonic analysis on the unit disk $`𝔻_1`$ has been well developed explicitly (cf. , pp. 29-72). I think that so far nobody has not investigated the theory of harmonic analysis on $`𝔻_{1,1}`$ explicitly. For example, inversion formula, Plancherel formula, Paley-Wiener theorem on $`𝔻_{1,1}`$ have not been described explicitly until now. It seems that it is interesting to develop the theory of harmonic analysis on the Siegel-Jacobi disk $`𝔻_{1,1}`$ explicitly. |
warning/0507/hep-ph0507163.html | ar5iv | text | # Decays of the 𝑋(3872) into 𝐽/𝜓 and Light Hadrons
## I Introduction
The $`X(3872)`$ is a narrow resonance near $`3872`$ MeV discovered by the Belle collaboration in electron-positron collisions through the $`B`$-meson decay $`B^\pm XK^\pm `$ followed by the decay $`XJ/\psi \pi ^+\pi ^{}`$ Choi:2003ue . Its existence has been confirmed by the CDF and DØ collaborations through its inclusive production in proton-antiproton collisions Acosta:2003zx ; Abazov:2004kp and by the Babar collaboration through the discovery mode $`B^\pm XK^\pm `$ Aubert:2004ns . The combined measurement of the mass of the $`X`$ is Olsen:2004fp
$`m_X=3871.9\pm 0.5\mathrm{MeV},`$ (1)
which is within 1 MeV of the threshold for the charm mesons $`D^0`$ and $`\overline{D}^0`$. The presence of the $`J/\psi `$ among the decay products of the $`X(3872)`$ motivates its interpretation as a charmonium state with constituents $`c\overline{c}`$ Barnes:2003vb ; Eichten:2004uh ; Quigg:2004nv ; Olsen:2004fp ; Quigg:2004vf . Two possibilities motivated by the proximity of the mass in Eq. (1) to the $`D^0\overline{D}^0`$ threshold are a hadronic molecule with constituents $`DD^{}`$ Tornqvist ; Voloshin:2003nt ; Wong:2003xk ; Braaten:2003he ; Swanson:2003tb ; Braaten:2004fk ; Swanson:2004pp ; Voloshin:2004mh ; Braaten:2004ai ; Braaten:2005jj ; AlFiky:2005jd and a “cusp” at the $`D^0\overline{D}^0`$ threshold associated with strong coupling to $`D^0\overline{D}^0`$ or $`D^0\overline{D}^0`$ Bugg:2004rk ; Bugg:2004sh . Other proposed interpretations include a tetraquark with constituents $`c\overline{c}q\overline{q}`$ Vijande:2004vt , a “hybrid charmonium” state with constituents $`c\overline{c}g`$ Close:2003mb ; Li:2004st , a glueball with constituents $`ggg`$ Seth:2004zb , and a diquark-antidiquark bound state with constituents $`cu+\overline{c}\overline{u}`$ Maiani:2004vq . The interpretation as a $`DD^{}`$ molecule is particularly predictive because the small binding energy implies that the molecule has universal properties that are completely determined by the binding energy Voloshin:2003nt ; Braaten:2003he ; Braaten:2004fk ; Braaten:2004ai . The small binding energy can be further exploited through factorization formulas for production and decay rates of the $`X`$ Braaten:2005jj .
Measurements of the decays of the $`X`$ can be used to determine its quantum numbers and narrow down the possibilities Olsen:2004fp ; Close:2003sg ; Pakvasa:2003ea ; Rosner:2004ac ; Kim:2004cz ; Abe:2005iy . The upper bound on the decay width of the $`X`$ is Choi:2003ue
$`\mathrm{\Gamma }_X<2.3\mathrm{MeV}(90\%\mathrm{C}.\mathrm{L}.),`$ (2)
which is much narrower than other charmonium states above the $`D\overline{D}`$ threshold. The product of the branching fractions associated with the discovery channel is Choi:2003ue ; Aubert:2004ns ; Olsen:2004fp
$`\mathrm{Br}[B^+XK^+]\mathrm{Br}[XJ/\psi \pi ^+\pi ^{}]=(1.3\pm 0.3)\times 10^5.`$ (3)
The invariant mass distribution of the two pions from the decay $`XJ/\psi \pi ^+\pi ^{}`$ seems to peak near the upper endpoint, which suggests that the pions come from a virtual $`\rho `$ resonance. Recently the Belle collaboration has observed the $`X(3872)`$ in the decay mode $`XJ/\psi \pi ^+\pi ^{}\pi ^0`$ Abe:2005iy with the branching ratio
$`{\displaystyle \frac{\mathrm{Br}[XJ/\psi \pi ^+\pi ^{}\pi ^0]}{\mathrm{Br}[XJ/\psi \pi ^+\pi ^{}]}}`$ $`=`$ $`1.0\pm 0.4_{\mathrm{stat}}\pm 0.3_{\mathrm{syst}}.`$ (4)
The invariant mass distribution of the three pions indicates that they come predominantly from a virtual $`\omega `$ resonance Abe:2005ix . If the decays into $`J/\psi \pi ^+\pi ^{}`$ and $`J/\psi \pi ^+\pi ^{}\pi ^0`$ are interpreted as $`J/\psi \rho ^{}`$ and $`J/\psi \omega ^{}`$, the approximate equality of the branching fractions in Eq. (4) implies a large violation of isospin symmetry. The Belle collaboration also reported evidence for the decay $`XJ/\psi \gamma `$ Abe:2005iy with the branching ratio
$`{\displaystyle \frac{\mathrm{Br}[XJ/\psi \gamma ]}{\mathrm{Br}[XJ/\psi \pi ^+\pi ^{}]}}`$ $`=`$ $`0.14\pm 0.05.`$ (5)
The observation of the decay into $`J/\psi \gamma `$ establishes the charge conjugation of the $`X`$ to be $`+`$. By analyzing angular distributions in the decay of $`X`$ into $`J/\psi \pi ^+\pi ^{}`$, the Belle collaboration has ruled out all $`J^{P+}`$ assignments for $`X`$ with $`J2`$ other than $`1^{++}`$ and $`2^{++}`$ Abe:2005iy . Upper limits have been placed on the branching fractions for other decay modes of the $`X`$, including $`D^0\overline{D}^0`$, $`D^+D^{}`$, $`D^0\overline{D}^0\pi ^0`$ Abe:2003zv , $`\chi _{c1}\gamma `$, $`\chi _{c2}\gamma `$, $`J/\psi \pi ^0\pi ^0`$ Abe:2004sd , and $`J/\psi \eta `$ Aubert:2004fc . Upper limits have also been placed on the partial widths for the decay of $`X`$ into $`e^+e^{}`$ Yuan:2003yz ; Metreveli:2004px and into $`\gamma \gamma `$ Metreveli:2004px .
The possibility that charm mesons might form molecular states was considered shortly after the discovery of charm Bander:1975fb ; Voloshin:ap ; DeRujula:1976qd ; Nussinov:1976fg . In 1993, Tornqvist made a quantitative study of the possibility of molecular states of charm mesons using a one-pion-exchange potential model Tornqvist:1993ng . He found that the isospin-0 combinations of $`D\overline{D}^{}`$ and $`D^{}\overline{D}`$ could form weakly-bound states in the S-wave $`1^{++}`$ channel and in the P-wave $`0^+`$ channel. After the discovery of the $`X(3872)`$, Tornqvist pointed out that because the binding energy is small compared to the splitting between the $`D^+D^{}`$ and $`D^0\overline{D}^0`$ thresholds, there will be large violations of isospin symmetry Tornqvist . Swanson considered a potential model that includes both one-pion-exchange and quark exchange and found that the $`C=+`$ superposition of $`D^0\overline{D}^0`$ and $`D^0\overline{D}^0`$ could form a weakly-bound state in the S-wave $`1^{++}`$ channel Swanson:2003tb . Swanson’s model included not only the charmed mesons $`D^0\overline{D}^0`$, $`D^0\overline{D}^0`$, $`D^+D^{}`$, and $`D^+D^{}`$, but also two other pairs of hadrons with nearby thresholds: $`J/\psi \rho ^0`$ and $`J/\psi \omega `$. His prediction that the branching fraction for the decay of $`X`$ into $`J/\psi \pi ^+\pi ^{}\pi ^0`$ should be comparable to that for decay into $`J/\psi \pi ^+\pi ^{}`$ was verified by the Belle collaboration Abe:2005ix ; Abe:2005iy .
In this paper, we analyze decays of the $`X(3872)`$ into $`J/\psi `$ and light hadrons under the assumption that $`X`$ is a loosely-bound $`DD^{}`$ molecule and that these decays proceed through transitions of $`X`$ to $`J/\psi \rho `$ and $`J/\psi \omega `$. In Section II, we summarize some of the universal results for a system with large scattering length and we give the current constraints on the real and imaginary parts of the large scattering length for the $`DD^{}`$ system. In Section III, we discuss other hadronic states with thresholds near the $`D^0\overline{D}^0`$ threshold and we summarize results from Swanson’s model for the $`X(3872)`$. In Section IV, we calculate the differential distributions for decays of $`X`$ into $`J/\psi \pi ^+\pi ^{}`$, $`J/\psi \pi ^+\pi ^{}\pi ^0`$, and $`J/\psi \pi ^0\gamma `$ using an effective lagrangian that reproduces decays of the vector mesons. We also use vector meson dominance to calculate the partial width for the decay into $`J/\psi \gamma `$. The normalizations of the decay rates are determined by unknown coupling constants for the transitions of $`X`$ to $`J/\psi \rho `$ and $`J/\psi \omega `$. The dependence of the coupling constants on the binding energy and the total width of the $`X`$ is deduced using a factorization formula. In Section V, we use a two-channel scattering model to illustrate how the coupling constants can be related to probabilities for the $`J/\psi \rho `$ and $`J/\psi \omega `$ components of the $`X`$. We use the probability for $`J/\psi \omega `$ in Swanson’s model to give a quantitative prediction for the partial decay rate of the $`X`$ into $`J/\psi \pi ^+\pi ^{}\pi ^0`$ as a function of the binding energy and the total width of the $`X`$. A summary of our results is given in Section VI. An updated determination of the parameters in the effective lagrangian for light pseudoscalar and vector mesons is given in an Appendix.
## II Universality and the $`DD^{}`$ System
The mass of the $`X`$ is extremely close to the $`D^0\overline{D}^0`$ threshold: $`m_{D^0}+m_{D^0}=3871.3\pm 1.0`$ MeV. From the mass measurement in Eq. (1), the difference is
$$m_X(m_{D^0}+m_{D^0})=+0.6\pm 1.1\mathrm{MeV}.$$
(6)
Most of the uncertainty comes from the experimental uncertainty in $`2m_{D^0}`$, because the mass difference $`m_{D^0}m_{D^0}`$ has a much smaller uncertainty. If the $`X`$ were a $`D^0\overline{D}^0`$/$`D^0\overline{D}^0`$ molecule, the energy difference in Eq. (6) would have to be negative, corresponding to a positive binding energy defined by
$$E_X=(m_{D^0}+m_{D^0})m_X.$$
(7)
The measurement of the binding energy in Eq. (6) is compatible with a small negative value corresponding to a $`DD^{}`$ molecule. However, the central value of the energy difference in Eq. (6) is positive, corresponding to a resonance in $`D^0\overline{D}^0`$ and $`D^0\overline{D}^0`$ scattering rather than a bound state. Bugg has referred to this possibility as a “cusp state” Bugg:2004rk ; Bugg:2004sh , because the line shape of the $`X`$ in some of its decay modes has a cusp at the $`D^0\overline{D}^0`$ threshold.
The energy difference in Eq. (6) is tiny compared to the natural energy scale for binding by the pion exchange interaction: $`m_\pi ^2/2\mu 10`$ MeV, where $`\mu `$ is the reduced mass of $`D^0`$ and $`\overline{D}^0`$:
$`\mu ={\displaystyle \frac{m_{D^0}m_{D^0}}{m_{D^0}+m_{D^0}}}=966.5\pm 0.3\mathrm{MeV}.`$ (8)
Whether the energy difference is positive or negative, its unnaturally small value implies that if the $`X`$ couples to $`D^0\overline{D}^0`$ and $`D^0\overline{D}^0`$, the S-wave scattering lengths for those channels must be large compared to the natural length scale $`1/m_\pi `$ associated with the pion exchange interaction. Since the experimental evidence favors the charge conjugation quantum number $`C=+`$, we assume that there is a large scattering length $`a`$ in the $`C=+`$ channel and that the scattering length in the $`C=`$ channel is negligible in comparison. In this case, the scattering lengths for elastic $`D^0\overline{D}^0`$ scattering and elastic $`D^0\overline{D}^0`$ scattering are both $`a/2`$.
Nonrelativistic few-body systems with short-range interactions and a large scattering length have universal properties that depend on the scattering length but are otherwise insensitive to details at distances small compared to $`|a|`$ Braaten:2004rn . We consider the scattering length to be large if it is much larger than the natural momentum scale associated with low-energy scattering. The universal results are encoded in the truncated connected transition amplitude, which is a function of the total energy $`E`$ of the two particles in the rest frame of the pair:
$$𝒜(E)=\frac{2\pi /\mu }{1/a+\sqrt{2\mu E}},$$
(9)
where $`\mu `$ is the reduced mass of the two particles. If $`a`$ is real and positive, this amplitude has a pole on the physical sheet at $`E=1/(2\mu a^2)`$, indicating the existence of a weakly-bound state with the universal binding energy
$$E_X=\frac{1}{2\mu a^2}.$$
(10)
The universal momentum-space wavefunction of this bound state is
$`\psi (p)={\displaystyle \frac{(8\pi /a)^{1/2}}{p^2+1/a^2}}.`$ (11)
The universal amplitude for transitions from the bound state to a scattering state consisting of two particles with small relative momentum is determined by the residue of the pole in $`𝒜(E)`$:
$$𝒜_X=\frac{\sqrt{2\pi }}{\mu }a^{1/2}.$$
(12)
If there is an inelastic scattering channel, the large scattering length $`a`$ has a negative imaginary part. It is convenient to express the complex scattering length in the form
$`{\displaystyle \frac{1}{a}}=\gamma _{\mathrm{re}}+i\gamma _{\mathrm{im}},`$ (13)
where $`\gamma _{\mathrm{re}}`$ and $`\gamma _{\mathrm{im}}`$ are real and $`\gamma _{\mathrm{im}}0`$. The universal expression for the binding energy in Eq. (10) has an imaginary part $`i\mathrm{\Gamma }_X/2`$, where
$`\mathrm{\Gamma }_X=2\gamma _{\mathrm{re}}\gamma _{\mathrm{im}}/\mu .`$ (14)
If $`\gamma _{\mathrm{im}}<\gamma _{\mathrm{re}}`$, $`\mathrm{\Gamma }_X`$ is the full width at half maximum of a resonance in the inelastic channel Braaten:2005jj . It therefore can be interpreted as the rate for the decay of the bound state into the inelastic channel. The peak of the resonance is below the threshold by the amount
$`E_X=\gamma _{\mathrm{re}}^2/(2\mu ).`$ (15)
We can therefore interpret this expression as the binding energy of the resonance Braaten:2005jj .
The observed decays of the $`X`$ imply that there are inelastic scattering channels, so $`a`$ has a negative imaginary part. It can be parameterized in terms of the real and imaginary parts of $`1/a`$ as in Eq. (13). Our interpretation of $`X`$ as a bound state requires $`\gamma _{\mathrm{re}}>0`$. The energy difference in Eq. (6) puts an upper bound on $`\gamma _{\mathrm{re}}`$:
$`\gamma _{\mathrm{re}}<40\mathrm{MeV}(90\%\mathrm{C}.\mathrm{L}.).`$ (16)
The upper bound on the width in Eq. (2) puts an upper bound on the product of $`\gamma _{\mathrm{re}}`$ and $`\gamma _{\mathrm{im}}`$:
$`\gamma _{\mathrm{re}}\gamma _{\mathrm{im}}<(33\mathrm{MeV})^2(90\%\mathrm{C}.\mathrm{L}.).`$ (17)
There is also a lower bound on the width of the $`X`$ from its decays into $`D^0\overline{D}^0\pi ^0`$ and $`D^0\overline{D}^0\gamma `$, which both proceed through the decay of a constituent $`D^{}`$. These decays involve interesting interference effects, but the decay rates have smooth limits as the binding energy is tuned to 0 Voloshin:2003nt . In this limit, the constructive interference increases the decay rate by a factor of 2, so the partial widths of $`X`$ into $`D^0\overline{D}^0\pi ^0`$ and $`D^0\overline{D}^0\gamma `$ add up to $`2\mathrm{\Gamma }[D^0]`$. The width of $`D^0`$ has not been measured, but it can be deduced from other information about the decays of $`D^0`$ and $`D^+`$. Using the total width of the $`D^+`$, its branching fraction into $`D^+\pi ^0`$, and isospin symmetry, we can deduce the partial width of $`D^0`$ into $`D^0\pi ^0`$ to be $`42\pm 10`$ keV. The total width of the $`D^0`$ can then be obtained by dividing by its branching fraction into $`D^0\pi ^0`$: $`\mathrm{\Gamma }[D^0]=68\pm 16`$ keV. The sum of the partial widths of $`X`$ into $`D^0\overline{D}^0\pi ^0`$ and $`D^0\overline{D}^0\gamma `$ is therefore $`136\pm 32`$ keV. The resulting lower bound on the product of $`\gamma _{\mathrm{re}}`$ and $`\gamma _{\mathrm{im}}`$ is
$`\gamma _{\mathrm{re}}\gamma _{\mathrm{im}}>(7\mathrm{MeV})^2(90\%\mathrm{C}.\mathrm{L}.).`$ (18)
By combining this with the upper bound on $`\gamma _{\mathrm{re}}`$ in Eq. (16), we can infer that $`\gamma _{\mathrm{im}}>1`$ MeV.
## III Hadronic states with nearby thresholds
The state of the $`X`$ can be written schematically as
$`|X={\displaystyle \frac{Z^{1/2}}{\sqrt{2}}}\left(|D^0\overline{D}^0+|D^0\overline{D}^0\right)+{\displaystyle \underset{H}{}}Z_H^{1/2}|H,`$ (19)
where $`Z`$ is the probability for the $`X`$ to be in the $`D^0\overline{D}^0/D^0\overline{D}^0`$ state and $`Z_H`$ is the probability for the $`X`$ to be in another hadronic state $`H`$. The hadronic states $`H`$ in (19) could include charmonium states, other charm meson pairs such as $`D^\pm D^{}`$, states consisting of a charmonium and a light hadron such as $`J/\psi \rho `$ and $`J/\psi \omega `$, etc. An expansion of $`|X`$ in terms of hadronic states can be valid only if there is an ultraviolet cutoff on the energy difference with respect to the $`D^0\overline{D}^0`$ threshold. The probabilities $`Z_H`$ depend on that cutoff. Universality implies that as $`a\mathrm{}`$, $`Z`$ approaches to $`1`$ and $`Z_H`$ scales as $`1/a`$ Braaten:2003he .
The small binding energy of $`X`$ compared to the natural energy scale $`m_\pi ^2/(2\mu )=10`$ MeV associated with pion exchange implies resonant S-wave interactions in the $`D^0\overline{D}^0`$ and $`D^0\overline{D}^0`$ systems. If there are other hadronic channels whose thresholds differ from the $`D^0\overline{D}^0`$ threshold by less than 10 MeV, there would be resonant S-wave interactions in those channels as well. In this case, it would be necessary to treat all the resonating channels as a coupled-channel system, with a large elastic scattering length for each channel and a large transition scattering length for each pair of channels. The hadronic states $`D^\pm D^{}`$, $`J/\psi \rho `$, and $`J/\psi \omega `$ have thresholds that are relatively close to the $`D^0\overline{D}^0`$ threshold. The energy gaps between these other thresholds and the $`D^0\overline{D}^0`$ threshold are
$`m_{D^\pm }+m_D^{}(m_{D^0}+m_{D^0})`$ $`=`$ $`+8.1\pm 0.1\mathrm{MeV},`$ (20a)
$`m_{J/\psi }+m_\rho (m_{D^0}+m_{D^0})`$ $`=`$ $`+1.4\pm 1.1\mathrm{MeV},`$ (20b)
$`m_{J/\psi }+m_\omega (m_{D^0}+m_{D^0})`$ $`=`$ $`+8.2\pm 1.0\mathrm{MeV}.`$ (20c)
The small uncertainty in Eq. (20a) comes from using mass differences between charm mesons to calculate the energy gap. The uncertainties in Eqs. (20b) and (20c) are dominated by the uncertainty in $`2m_{D^0}`$. The energy gaps in the $`D^\pm D^{}`$ and $`J/\psi \omega `$ channels are comparable to the natural energy scale of about $`10`$ MeV associated with pion exchange. The energy gap in Eq. (20b) for the $`J/\psi \rho `$ channel is much smaller. However whether any of these channels can have resonant interactions with $`D^0\overline{D}^0`$ or $`D^0\overline{D}^0`$ is determined not only by the real parts of the energy gaps, which are given in Eqs. (20), but also by the imaginary parts, which can be obtained by replacing each mass $`m`$ by $`mi\mathrm{\Gamma }/2`$, where $`\mathrm{\Gamma }`$ is the width of the particle. If there are large differences between the widths of the various particles, it is necessary only to take into account the largest width among the particles in each channel. The largest width in each of the three channels is
$`\mathrm{\Gamma }[D^\pm ]`$ $`=`$ $`0.096\pm 0.022\mathrm{MeV},`$ (21a)
$`\mathrm{\Gamma }[\rho ]`$ $`=`$ $`150.3\pm 1.6\mathrm{MeV},`$ (21b)
$`\mathrm{\Gamma }[\omega ]`$ $`=`$ $`8.49\pm 0.08\mathrm{MeV}.`$ (21c)
For the $`D^\pm D^{}`$ and $`J/\psi \omega `$ channels, the magnitude $`|\mathrm{\Delta }|`$ of the complex energy gap is comparable to the natural energy scale 10 MeV associated with pion exchange between $`D`$ and $`D^{}`$. The large width of the $`\rho `$ makes $`|\mathrm{\Delta }|`$ for the $`J/\psi \rho `$ channel much larger than the natural energy scale.
Because the complex energy gap $`\mathrm{\Delta }`$ for the other hadronic channels with nearby thresholds are comparable to or larger than the natural energy scale, these channels need not be taken into account explicitly in calculations of quantities that have nontrivial universal limits as $`a\pm \mathrm{}`$. Their dominant effects enter through the complex-valued scattering length $`a`$. A coupled-channel model that includes other hadronic states with nearby thresholds could still be useful for estimating nonuniversal quantities or for calculating nonuniversal corrections to the universal predictions.
Swanson has constructed a model of the $`X`$ and the hadronic states with nearby thresholds and used it to predict some of the properties of the $`X`$ Swanson:2003tb ; Swanson:2004pp . In particular, he predicted correctly that the branching fraction for $`XJ/\psi \pi ^+\pi ^{}\pi ^0`$ is comparable to that for $`XJ/\psi \pi ^+\pi ^{}`$. In addition to the channel $`D^0\overline{D}^0+D^0\overline{D}^0`$, Swanson’s model includes $`D^+D^{}+D^+D^{}`$, $`J/\psi \rho `$, and $`J/\psi \omega `$. It includes the S-wave and D-wave channels for $`DD^{}`$, but only the S-wave channel for $`J/\psi V`$, where $`V`$ is the vector meson $`\rho `$ or $`\omega `$. Thus the model has 6 coupled channels. The interactions between the hadrons are modeled by potentials: one-pion-exchange potentials for the S-wave and D-wave $`DD^{}`$ channels and for transitions between those channels and Gaussian potentials for the transitions between the S-wave $`DD^{}`$ channels and the S-wave $`J/\psi V`$ channels to simulate the effects of quark exchange. The one-pion-exchange potential is singular at short distances and it was regularized by an ultraviolet momentum cutoff $`\mathrm{\Lambda }`$. The nonrelativistic Schrödinger equation for the 6 coupled channels was solved numerically. A bound state with the quantum numbers $`J^{PC}=1^{++}`$ of $`X`$ appeared when the ultraviolet cutoff exceeded the critical value $`\mathrm{\Lambda }_c=1.45`$ GeV. The binding energy of the $`X`$ could be adjusted by varying the ultraviolet cutoff.
Swanson solved the coupled channel problem under the assumption that the $`\rho `$ and $`\omega `$ are stable hadrons with equal masses $`m_\rho =m_\omega =782.6`$ MeV. The reason for using an unphysical value for $`m_\rho `$ is that the central PDG value from 2002 and earlier, $`m_\rho =771.1`$ MeV, is below the $`D^0\overline{D}^0`$ threshold. If such a value had been used, it would have been necessary to treat $`J/\psi \rho `$ states as scattering states. This complication was avoided by using a value of $`m_\rho `$ above the $`D^0\overline{D}^0`$ threshold. Note that in Eq. (20b), we have taken the updated 2004 PDG value $`m_\rho =775.8\pm 0.5`$ MeV Eidelman:2004wy , which gives a $`J/\psi \rho `$ threshold that is a few MeV above the mass of the $`X`$.
Swanson calculated the probabilities for each component of the wavefunction of $`X`$ for values of the ultraviolet cutoff that correspond to varying the binding energy $`E_X`$ from 0.7 MeV to 23.2 MeV Swanson:2003tb . His results for the probabilities $`Z_{\psi \omega }`$ and $`Z_{\psi \rho }`$ are shown as dots in Fig. 1. Since the binding energy of the $`X`$ is known to be less than 1 MeV, only the lowest two values of $`E_X`$ could be physically relevant. For the lowest value $`E_X=0.7`$ MeV, the probabilities were $`Z_{\psi \omega }=9.6`$%, $`Z_{D^\pm D^{}}=7.9`$%, and $`Z_{\psi \rho }=0.86`$%. The total probability for the $`D^0\overline{D}^0`$ and $`D^0\overline{D}^0`$ components of the wavefunction is 81.6%.
In Fig. 1, the dotted lines have the scaling behavior $`E_X^{1/2}`$ predicted by universality and are normalized so that they pass through the dot at $`E_X=0.7`$ MeV. The probability $`Z_{\psi \omega }`$ clearly exhibits the universal behavior. The probability $`Z_{\psi \rho }`$ does not. This may be related to the fact that $`Z_{\psi \rho }`$ is more than an order of magnitude smaller than $`Z_{\psi \omega }`$. Because of the weaker coupling of the $`X`$ to isospin-1 states, the scaling region for isospin-1 states may not set in until a much smaller value of $`E_X`$.
Swanson estimated the partial widths for the decays of $`X`$ into $`J/\psi h`$, where $`h`$ is the light hadronic state $`\pi ^+\pi ^{}`$, $`\pi ^+\pi ^{}\pi ^0`$, $`\pi ^0\gamma `$, or $`\pi ^+\pi ^{}\gamma `$, using a simple ad hoc recipe. The partial width into $`J/\psi h`$ was taken to be the sum over the vector mesons $`V=\rho ,\omega `$ of the product of the probability $`Z_{\psi V}`$ for the $`J/\psi V`$ component of the wavefunction and the partial width $`\mathrm{\Gamma }[Vh]`$ for the decay of the vector meson:
$$\mathrm{\Gamma }[XJ/\psi h]\underset{V}{}Z_{\psi V}\mathrm{\Gamma }[Vh].$$
(22)
For the smallest value of the binding energy that was considered, $`E_X=0.7`$ MeV, the resulting estimates of the partial widths for decay into $`J/\psi h`$ were 1290 keV, 720 keV, 70 keV, and 13 keV for $`h=\pi ^+\pi ^{}`$, $`\pi ^+\pi ^{}\pi ^0`$, $`\pi ^0\gamma `$, and $`\pi ^+\pi ^{}\gamma `$, respectively. The ratio of the partial widths into $`J/\psi \pi ^+\pi ^{}\pi ^0`$ and $`J/\psi \pi ^+\pi ^{}`$ was predicted to be 0.56 for $`E_X=0.7`$ MeV. Remarkably, this prediction agrees with the subsequent measurement by the Belle collaboration given in Eq. (4) to within the experimental errors Abe:2004sd . The approximately equal branching fractions are a fortuitous result of an amplitude for $`XJ/\psi \omega `$ that is much larger than the amplitude for $`XJ/\psi \rho `$ and an amplitude for $`\omega \pi ^+\pi ^{}\pi ^0`$ that is much smaller than that for $`\rho \pi ^+\pi ^{}`$. The suppression of the amplitude for $`XJ/\psi \rho `$ is related to the fact that in the isospin symmetric limit in which the mass difference between neutral and charged $`D`$’s is neglected, there is binding in the isospin-0 channel but not in the isospin-1 channel Tornqvist:1993ng .
Swanson has also used his model to calculate the rates for several other decay modes of the $`X`$ Swanson:2004pp . The decay rate into $`J/\psi \gamma `$ has contributions from transitions to $`J/\psi \rho `$ and $`J/\psi \omega `$ that can be calculated using vector meson dominance. It also has contributions from the annihilation of the $`u`$ and $`\overline{u}`$ from the charm mesons that are the constituents of the $`X`$. Swanson’s prediction for the partial width into $`J/\psi \gamma `$ for an $`X`$ with a binding energy of 1 MeV is 8 keV. Decay modes that receive contributions only from $`u\overline{u}`$ annihilation, such as $`\psi (2S)\gamma `$, $`KK^{}`$, and $`\pi \rho `$, have much smaller partial widths.
## IV Decays of $`𝑿`$ into $`𝑱\mathbf{/}𝝍𝒉`$
In this section, we calculate the differential decay rates of the $`X`$ into $`J/\psi h`$, where the hadronic system $`h`$ is $`\pi ^+\pi ^{}\pi ^0`$, $`\pi ^+\pi ^{}`$, $`\pi ^0\gamma `$, or $`\gamma `$. We assume that these decays proceed through transitions of $`X`$ to $`J/\psi \rho `$ and $`J/\psi \omega `$. We calculate the differential decay rates in terms of two unknown complex coupling constants using an effective lagrangian that reproduces the decays of the light vector mesons.
### IV.1 Vector Meson Decay Amplitudes
We assume that the decay of $`X`$ into $`J/\psi h`$, where $`h`$ is a system of light hadrons, proceeds through transitions to $`J/\psi V`$, where $`V`$ is one of the vector mesons $`\rho `$ or $`\omega `$, followed by the decay of the vector meson into $`h`$. Because the mass of the $`X`$ is so close to the threshold for $`J/\psi V`$, the vector meson is almost on its mass shell. Any model that reproduces the decays of the vector mesons should also accurately describe the decay of the virtual vector meson in the $`J/\psi V`$ component of $`X`$. In Ref. Braaten:1989zn , the semileptonic branching fractions for the $`\tau `$ lepton were calculated using an effective lagrangian for light pseudoscalar and vector mesons with $`U(3)\times U(3)`$ chiral symmetry. All the parameters in the effective lagrangian, aside from the pion decay constant, were determined directly from decays of the vector mesons $`\rho `$ and $`\omega `$. That same effective lagrangian can be used to calculate the partial widths of $`X`$ into $`J/\psi h`$. An updated determination of the parameters in that effective lagrangian is given in the Appendix.
The T-matrix element for the decay of a vector meson $`V`$ into the light hadronic state $`h`$ can be expressed in the form
$`𝒯[Vh]`$ $`=`$ $`ϵ_V^\mu 𝒜_\mu [Vh],`$ (23)
where $`ϵ_V`$ is the polarization vector of the vector meson. The amplitude $`𝒜_\mu `$ for the decay $`\rho \pi ^+\pi ^{}`$ is
$`𝒜_\mu [\rho \pi ^+\pi ^{}]=\frac{1}{2}G_{v\pi \pi }\left(p_+p_{}\right)_\mu .`$ (24)
The value of the coupling constant $`G_{v\pi \pi }`$ is given in Eq. (67a). The amplitude $`𝒜_\mu `$ for the decay $`\omega \pi ^+\pi ^{}\pi ^0`$ is
$`𝒜_\mu [\omega \pi ^+\pi ^{}\pi ^0]`$ $`=`$ $`{\displaystyle \frac{4\sqrt{3}(\mathrm{cos}\theta _v+\sqrt{2}\mathrm{sin}\theta _v)}{F_\pi ^3}}\epsilon _{\mu \nu \alpha \beta }p_+^\nu p_{}^\alpha p_0^\beta `$ (25)
$`\times \left(C_{v3\pi }+{\displaystyle \frac{G_{v\pi \pi }C_{vv\pi }F_\pi ^2}{m_v^2}}\left(1\frac{1}{3}\left[f_\rho (s_{12})+f_\rho (s_{23})+f_\rho (s_{31})\right]\right)\right),`$
where $`s_{12}`$, $`s_{23}`$, and $`s_{31}`$ are the invariant masses of the three different pion pairs and
$$f_V(s)\frac{s}{sM_V^2+im_V\mathrm{\Gamma }_V}$$
(26)
is a vector meson resonance factor that vanishes at $`s=0`$. We have denoted the 4-momenta of $`\pi ^+`$, $`\pi ^{}`$, and $`\pi ^0`$ by $`p_+`$, $`p_{}`$, and $`p_0`$, respectively. The pion decay constant is $`F_\pi =93`$ MeV, the values of the parameters $`C_{v3\pi }`$ and $`G_{v\pi \pi }C_{vv\pi }F_\pi ^2/m_v^2`$ are given by Eqs. (67b) and (67c), and the value of the light vector meson mixing angle $`\theta _v`$ is given by Eq. (69). The amplitudes $`𝒜_\mu `$ for the radiative decays of the vector mesons are
$`𝒜_\mu [\rho ^+\pi ^+\gamma ]`$ $`=`$ $`{\displaystyle \frac{4e}{3F_\pi }}\left(C_{v\pi \gamma }+{\displaystyle \frac{G_{v\gamma }C_{vv\pi }F_\pi ^2}{m_v^2}}\right)\epsilon _{\mu \nu \alpha \beta }Q^\nu p^\alpha ϵ_\gamma ^\beta ,`$ (27a)
$`𝒜_\mu [\omega \pi ^0\gamma ]`$ $`=`$ $`{\displaystyle \frac{4(\mathrm{cos}\theta _v+\sqrt{2}\mathrm{sin}\theta _v)e}{\sqrt{3}F_\pi }}\left(C_{v\pi \gamma }+{\displaystyle \frac{G_{v\gamma }C_{vv\pi }F_\pi ^2}{m_v^2}}\right)\epsilon _{\mu \nu \alpha \beta }Q^\nu p^\alpha ϵ_\gamma ^\beta ,`$ (27b)
where $`Q`$ and $`p`$ are the 4-momenta of the vector meson and the pion and $`ϵ_\gamma `$ is the polarization vector of the photon. The values of the parameters $`C_{v\pi \gamma }`$ and $`G_{v\gamma }C_{vv\pi }F_\pi ^2/m_v^2`$ are given by Eqs. (68b) and (68c).
The amplitudes $`𝒜_\nu `$ in Eqs. (24), (25), and (27) all satisfy $`Q^\nu 𝒜_\nu =0`$, where $`Q`$ is the 4-momentum of the vector meson. This condition is satisfied in any model consistent with vector meson dominance. The assumption of vector meson dominance is that the amplitude for the production of a real photon in a hadronic process can be expressed as the sum of over vector mesons $`V`$ of the amplitude for producing $`V`$ multiplied by a coupling constant for the transition $`V\gamma `$. The condition $`Q^\nu 𝒜_\nu =0`$ is required for the gauge invariance of the resulting amplitude for real photon production.
The T-matrix element for $`X`$ to decay into $`J/\psi `$ and a light hadronic system $`h`$ through a virtual vector meson resonance $`V`$ can be expressed as
$`𝒯[XJ/\psi h]=𝒜_\mu [XJ/\psi V]{\displaystyle \frac{g^{\mu \nu }}{Q^2m_V^2+im_V\mathrm{\Gamma }_V}}𝒜_\nu [Vh],`$ (28)
where $`Q`$ is the total 4-momentum of the hadronic system $`h`$ or, equivalently, of the virtual vector meson. We have used the condition $`Q^\nu 𝒜_\nu =0`$ to simplify the numerator of the vector meson propagator. The quantum numbers of the particles, together with Lorentz invariance, constrains the amplitude for $`XJ/\psi V`$ to be the sum of two terms. One of them is
$$𝒜_\mu [XJ/\psi V]=G_{X\psi V}\epsilon _{\mu \nu \alpha \beta }Q^\nu ϵ_X^\alpha ϵ_\psi ^\beta ,$$
(29)
where $`ϵ_X`$ and $`ϵ_\psi `$ are the polarization 4-vectors of the $`X`$ and the $`J/\psi `$ and $`G_{X\psi V}`$ is a dimensionless constant. The contraction of this amplitude with the polarization vector $`ϵ_V^{}`$ of the vector meson reduces in the rest frame of the vector meson to $`G_{X\psi V}m_V\mathit{ϵ}_𝑿(\mathit{ϵ}_\psi \times \mathit{ϵ}_V)^{}`$. The other independent amplitude $`𝒜_\mu `$ has the Lorentz structure $`\epsilon _{\mu \nu \alpha \beta }P^\nu ϵ_X^\alpha ϵ_\psi ^\beta `$. In the rest frame of the $`X`$, its contraction with $`ϵ_V^{}`$ is $`m_X\mathit{ϵ}_𝑿(\mathit{ϵ}_\psi \times \mathit{ϵ}_V)^{}`$. Since the mass of the $`X`$ is so close to the threshold for $`J/\psi V`$, the rest frames of the $`X`$ and $`V`$ are essentially identical. Thus the two independent Lorentz structures are essentially equivalent for decays that are dominated by the vector meson resonance. They give similar predictions for the partial widths for $`X`$ into $`J/\psi h`$ for $`h=\pi ^+\pi ^{}\pi ^0`$, $`\pi ^+\pi ^{}`$, or $`\pi ^0\gamma `$. The amplitude in Eq. (29) has the advantage that it is also consistent with the constraint $`Q^\mu 𝒜_\mu =0`$ required by vector meson dominance. Thus this amplitude can be used to calculate the decay of $`X`$ into $`J/\psi \gamma `$. We therefore take the transition amplitude for $`X`$ into $`J/\psi V`$ to be the expression in Eq. (29).
### IV.2 Decay into $`𝑱\mathbf{/}𝝍𝝅^\mathbf{+}𝝅^{\mathbf{}}`$
We assume that the decay of $`X`$ into $`J/\psi \pi ^+\pi ^{}`$ proceeds through a transition of $`X`$ to $`J/\psi \rho `$. The T-matrix element is then given in terms of the unknown coupling constant $`G_{X\psi \rho }`$ by Eqs. (28) and (29) with $`V=\rho `$. The expression for the amplitude $`𝒜_\nu `$ for $`\rho \pi ^+\pi ^{}`$ is given in Eq. (24). We obtain the decay rate by squaring the amplitude, summing over spins, and integrating over phase space. The differential decay rate into $`J/\psi \pi ^+\pi ^{}`$ as a function of the invariant mass $`Q`$ of the two pions is
$`{\displaystyle \frac{d\mathrm{\Gamma }}{dQ}}[XJ/\psi \pi ^+\pi ^{}]`$ $`=`$ $`{\displaystyle \frac{|G_{X\psi \rho }|^2G_{v\pi \pi }^2}{9216\pi ^3m_X^5m_\psi ^2}}{\displaystyle \frac{(Q^24m_\pi ^2)^{3/2}\lambda ^{1/2}(m_X,m_\psi ,Q)}{(Q^2m_\rho ^2)^2+m_\rho ^2\mathrm{\Gamma }_\rho ^2}}`$ (30)
$`\times \left[(m_X^2+m_\psi ^2)(m_X^2m_\psi ^2)^22(m_X^44m_X^2m_\psi ^2+m_\psi ^4)Q^2+(m_X^2+m_\psi ^2)Q^4\right],`$
where $`\lambda (x,y,z)`$ is the triangle function:
$`\lambda (x,y,z)=x^4+y^4+z^42(x^2y^2+y^2z^2+z^2x^2).`$ (31)
After integrating over the pion invariant mass, the decay rate is
$`\mathrm{\Gamma }[XJ/\psi \pi ^+\pi ^{}]`$ $`=`$ $`|G_{X\psi \rho }|^2(223\mathrm{keV}).`$ (32)
The shape of the pion invariant mass distribution for the decay of $`X`$ into $`J/\psi \pi ^+\pi ^{}`$ is shown in Fig. 2. Its qualitative features are dominated by the phase space factor $`\lambda ^{1/2}(m_X,m_\psi ,Q)`$, which cuts the distribution off at the endpoint $`Q=m_Xm_\psi `$, and the vector meson resonance factor, which has its maximum at $`Q=m_\rho `$ just outside the kinematic region. Most of the support for $`d\mathrm{\Gamma }/dQ`$ comes from within $`\mathrm{\Gamma }_\rho `$ of the upper endpoint.
### IV.3 Decay into $`𝑱\mathbf{/}𝝍𝝅^\mathbf{+}𝝅^{\mathbf{}}𝝅^\mathrm{𝟎}`$
We assume that the decay of $`X`$ into $`J/\psi \pi ^+\pi ^{}\pi ^0`$ proceeds through a transition of $`X`$ to $`J/\psi \omega `$. The T-matrix element is then given in terms of the unknown coupling constant $`G_{X\psi \omega }`$ by Eqs. (28) and (29) with $`V=\omega `$. The expression for the amplitude for $`\omega 3\pi `$ is given in Eq. (25). We obtain the decay rate by squaring the amplitude, summing over spins, and integrating over phase space. The differential decay rate into $`J/\psi \pi ^+\pi ^{}\pi ^0`$ as a function of the invariant mass $`Q`$ of the 3 pions can be reduced to a 2-dimensional integral:
$`{\displaystyle \frac{d\mathrm{\Gamma }}{dQ}}[XJ/\psi \pi ^+\pi ^{}\pi ^0]`$ $`=`$ $`{\displaystyle \frac{|G_{X\psi \omega }|^2(\mathrm{cos}\theta _v+\sqrt{2}\mathrm{sin}\theta _v)^2}{3072\pi ^5m_X^5m_\psi ^2F_\pi ^6}}{\displaystyle \frac{\lambda ^{1/2}(m_X,m_\psi ,Q)}{Q[(Q^2m_\omega ^2)^2+m_\omega ^2\mathrm{\Gamma }_\omega ^2]}}`$ (33)
$`\times \left[(m_X^2+m_\psi ^2)(m_X^2m_\psi ^2)^22(m_X^44m_X^2m_\psi ^2+m_\psi ^4)Q^2+(m_X^2+m_\psi ^2)Q^4\right]`$
$`\times {\displaystyle }ds_{12}{\displaystyle }ds_{23}[s_{12}s_{23}s_{31}m_\pi ^2(Q^2m_\pi ^2)^2]`$
$`\times \left|C_{v3\pi }+{\displaystyle \frac{G_{v\pi \pi }C_{vv\pi }F_\pi ^2}{m_v^2}}\left(1\frac{1}{3}\left[f_\rho (s_{12})+f_\rho (s_{23})+f_\rho (s_{31})\right]\right)\right|^2,`$
where $`s_{12}`$, $`s_{23}`$, and $`s_{31}`$ are the squares of the invariant masses of the three pairs of pions. We have suppressed the limits of integration in the integrals over $`s_{12}`$ and $`s_{23}`$. After integrating over the pion invariant masses, the decay rate is
$`\mathrm{\Gamma }[XJ/\psi \pi ^+\pi ^{}\pi ^0]`$ $`=`$ $`|G_{X\psi \omega }|^2(19.4\mathrm{keV}).`$ (34)
The shape of the pion invariant mass distributions for the decay of $`X`$ into $`J/\psi \pi ^+\pi ^{}\pi ^0`$ is shown in Fig. 3. Its qualitative features are dominated by the phase space factor $`\lambda ^{1/2}(m_X,m_\psi ,Q)`$, which cuts the distribution off at the endpoint $`Q=m_Xm_\psi `$, and the vector meson resonance factor, which has its maximum at $`Q=m_\omega `$ just outside the kinematic region. Most of the support for $`d\mathrm{\Gamma }/dQ`$ comes from within a few widths $`\mathrm{\Gamma }_\omega `$ of the upper endpoint.
The ratio of the decay rates in Eqs. (32) and (34) is
$`{\displaystyle \frac{\mathrm{\Gamma }[XJ/\psi \pi ^+\pi ^{}\pi ^0]}{\mathrm{\Gamma }[XJ/\psi \pi ^+\pi ^{}]}}=0.0870{\displaystyle \frac{|G_{X\psi \omega }|^2}{|G_{X\psi \rho }|^2}}.`$ (35)
By comparing this to Belle’s result in Eq. (4) for the ratio of the branching fractions, we can obtain an estimate of the ratio of the coupling constants:
$`{\displaystyle \frac{|G_{X\psi \omega }|^2}{|G_{X\psi \rho }|^2}}11.5\pm 5.7.`$ (36)
### IV.4 Decay into $`𝑱\mathbf{/}𝝍𝝅^\mathrm{𝟎}𝜸`$
We assume that the decay of $`X`$ into $`J/\psi \pi ^0\gamma `$ proceeds through transitions of $`X`$ to $`J/\psi \rho `$ and $`J/\psi \omega `$. The T-matrix element is then given by Eq. (28) summed over $`V=\rho ,\omega `$. The amplitudes for $`V\pi ^0\gamma `$ are given in Eqs. (27). The differential decay rate with respect to the invariant mass $`Q`$ of the $`\pi ^0\gamma `$ is
$`{\displaystyle \frac{d\mathrm{\Gamma }}{dQ}}[XJ/\psi \pi ^0\gamma ]`$ $`=`$ $`{\displaystyle \frac{\alpha _{\mathrm{em}}\left(C_{v\pi \gamma }+G_{v\gamma }C_{vv\pi }F_\pi ^2/m_v^2\right)^2}{648\pi ^2m_X^5m_\psi ^2F_\pi ^2}}{\displaystyle \frac{(Q^2m_\pi ^2)^3\lambda ^{1/2}(m_X,m_\psi ,Q)}{Q}}`$ (37)
$`\times \left[(m_X^2+m_\psi ^2)(m_X^2m_\psi ^2)^22(m_X^44m_X^2m_\psi ^2+m_\psi ^4)Q^2+(m_X^2+m_\psi ^2)Q^4\right]`$
$`\times \left|{\displaystyle \frac{G_{X\psi \rho }}{Q^2m_\rho ^2+im_\rho \mathrm{\Gamma }_\rho }}+{\displaystyle \frac{G_{X\psi \omega }\sqrt{3}(\mathrm{cos}\theta _v+\sqrt{2}\mathrm{sin}\theta _v)}{Q^2m_\omega ^2+im_\omega \mathrm{\Gamma }_\omega }}\right|^2.`$
After integrating over the $`\pi ^0\gamma `$ invariant mass, the decay rate is
$`\mathrm{\Gamma }[XJ/\psi \pi ^0\gamma ]`$ $`=`$ $`[|G_{X\psi \omega }|^2+0.026|G_{X\psi \rho }|^2`$ (38)
$`+(0.163\mathrm{cos}\varphi +0.215\mathrm{sin}\varphi )|G_{X\psi \omega }||G_{X\psi \rho }|](3.24\mathrm{keV}),`$
where $`\mathrm{exp}(i\varphi )`$ is the relative phase between $`G_{X\psi \omega }`$ and $`G_{X\psi \rho }`$. The estimate of the ratio $`|G_{X\psi \omega }|^2/|G_{X\psi \rho }|^2`$ in Eq. (36) suggests that the $`|G_{X\psi \omega }|^2`$ term in Eq. (38) dominates. If this is the case, the branching fraction for the decay of $`X`$ into $`J/\psi \pi ^0\gamma `$ should be smaller than that for $`J/\psi \pi ^+\pi ^{}\pi ^0`$ by a factor of about 0.17.
### IV.5 Decay into $`𝑱\mathbf{/}𝝍𝜸`$
Having chosen the transition amplitude in Eq. (29) so that it satisfies $`Q^\mu 𝒜_\mu =0`$, we can use vector meson dominance to calculate the partial width for the decay of $`X`$ into $`J/\psi \gamma `$. The T-matrix element is
$`𝒯[XJ/\psi \gamma ]`$ $`=`$ $`G_{v\gamma }F_\pi ^2e\left({\displaystyle \frac{G_{X\psi \rho }}{m_\rho ^2im_\rho \mathrm{\Gamma }_\rho }}+{\displaystyle \frac{G_{X\psi \omega }\mathrm{cos}\theta _v/\sqrt{3}}{m_\omega ^2im_\omega \mathrm{\Gamma }_\omega }}\right)\epsilon _{\mu \nu \alpha \beta }Q^\mu ϵ_X^\nu ϵ_{\psi }^{\alpha }{}_{}{}^{}ϵ_{\gamma }^{\beta }{}_{}{}^{},`$ (39)
where $`Q`$ is the 4-momentum of the photon. The value of the coupling constant $`G_{v\gamma }`$ is given in Eq. (68a). The result for the decay rate is
$`\mathrm{\Gamma }[XJ/\psi \gamma ]`$ $`=`$ $`{\displaystyle \frac{\alpha _{\mathrm{em}}G_{v\gamma }^2F_\pi ^4(m_X^2+m_\psi ^2)(m_X^2m_\psi ^2)^3}{24m_X^5m_\psi ^2}}`$ (40)
$`\times \left|{\displaystyle \frac{G_{X\psi \rho }}{m_\rho ^2im_\rho \mathrm{\Gamma }_\rho }}+{\displaystyle \frac{G_{X\psi \omega }\mathrm{cos}\theta _v/\sqrt{3}}{m_\omega ^2im_\omega \mathrm{\Gamma }_\omega }}\right|^2.`$
If the widths in the vector meson propagators are neglected and if we use $`m_\rho m_\omega `$, the decay rate in Eq. (40) reduces to
$`\mathrm{\Gamma }[XJ/\psi \gamma ]=|G_{X\psi \rho }+0.30G_{X\psi \omega }|^2(5.51\mathrm{keV}).`$ (41)
Our estimate in Eq. (36) implies that $`|G_{X\psi \omega }|`$ is much larger than $`|G_{X\psi \rho }|`$. However, the larger magnitude of $`G_{X\psi \omega }`$ is compensated by the vector meson mixing factor $`\mathrm{cos}\theta _v/\sqrt{3}=0.30`$, so the $`G_{X\psi \rho }`$ and $`G_{X\psi \omega }`$ terms may be equally important. Using the partial widths in Eqs. (32) and (34), we can relate the branching fractions for $`J/\psi \gamma `$ to those for $`J/\psi \pi ^+\pi ^{}`$ and $`J/\psi \pi ^+\pi ^{}\pi ^0`$:
$`\mathrm{Br}[XJ/\psi \gamma ]`$ $`=`$ $`0.025\mathrm{Br}[XJ/\psi \pi ^+\pi ^{}]+0.026\mathrm{Br}[XJ/\psi \pi ^+\pi ^{}\pi ^0]`$ (42)
$`+0.050\mathrm{cos}\varphi \left(\mathrm{Br}[XJ/\psi \pi ^+\pi ^{}]\mathrm{Br}[XJ/\psi \pi ^+\pi ^{}\pi ^0]\right)^{1/2},`$
where $`\mathrm{exp}(i\varphi )`$ is the relative phase between $`G_{X\psi \omega }`$ and $`G_{X\psi \rho }`$. This prediction is compatible with the measurements of the branching ratios in Eqs. (4) and (5) if the angle $`\varphi `$ is small.
### IV.6 Factorization of short-distance decay rates
The decay modes of the $`X(3872)`$ can be classified into long-distance decays and short-distance decays. The long-distance decay modes are $`D^0\overline{D}^0\pi ^0`$ and $`D^0\overline{D}^0\gamma `$, which proceed through the decay of a constituent $`D^0`$ or $`\overline{D}^0`$. These decays are dominated by a component of the wavefunction of the $`X`$ in which the separation of the $`D`$ and $`D^{}`$ is of order $`|a|`$. These long-distance decays involve interesting interference effects between the $`D^0\overline{D}^0`$ and $`D^0\overline{D}^0`$ components of the wavefunction Voloshin:2003nt . The short-distance decays are dominated by a component of the wavefunction in which the separation of the $`D`$ and $`D^{}`$ is of order $`m_\pi `$ or smaller. Examples are the observed decay modes $`J/\psi \pi ^+\pi ^{}`$, $`J/\psi \pi ^+\pi ^{}\pi ^0`$, and $`J/\psi \gamma `$.
Short-distance decays of the $`X`$ into a hadronic final state $`H`$ involve well-separated momentum scales. The $`DD^{}`$ wavefunction of the $`X`$ involves the momentum scale $`1/|a|`$ set by the large scattering length. The transition of the $`DD^{}`$ to $`H`$ involves momentum scales $`m_\pi `$ and larger. The separation of scales $`|a|1/m_\pi `$ can be exploited by using a factorization formula for the decay rate Braaten:2005jj . In limit $`|a|1/m_\pi `$, the leading term in the T-matrix element for the decay $`XH`$ can be separated into a short-distance factor and a long-distance factor:
$`𝒯[XH]=𝒜_{\mathrm{short}}[XH]\times 𝒜_X.`$ (43)
The short-distance factor $`𝒜_{\mathrm{short}}`$ in Eq. (43) has a well-behaved limit as $`|a|\mathrm{}`$. The long-distance factor $`𝒜_X`$ is the universal amplitude given in Eq. (12). If the complex scattering length is parameterized as in Eq. (13), this factor is
$`𝒜_X=\left(\sqrt{2\pi }/\mu \right)\left(\gamma _{\mathrm{re}}+i\gamma _{\mathrm{im}}\right)^{1/2}.`$ (44)
When applied to decays of $`X`$ into $`J/\psi `$ and light hadrons, the factorization formula in Eq. (43) implies that the coupling constants $`G_{X\psi \rho }`$ and $`G_{X\psi \omega }`$ have a long-distance factor $`𝒜_X`$.
The factorization formula for the T-matrix element in Eq. (43) implies a factorization formula for the decay rate:
$`\mathrm{\Gamma }[XH]=\mathrm{\Gamma }_{\mathrm{short}}[XH]\times |𝒜_X|^2.`$ (45)
The short-distance factor $`\mathrm{\Gamma }_{\mathrm{short}}`$ in Eq. (45) has a well-behaved limit as $`|a|\mathrm{}`$. Using the expressions in Eqs. (15) and (14) for the binding energy and the width of the molecule, the long-distance factor in Eq. (45) can be expressed as
$`|𝒜_X|^2=\sqrt{8\pi ^2/\mu ^3}\left[E_X+\mathrm{\Gamma }_X^2/(16E_X)\right]^{1/2}.`$ (46)
Predictions for the rates for short-distance decays of the $`X`$ can be obtained from models for low-energy hadrons in which the parameters have been tuned to obtain a small binding energy $`E_X`$, such as Swanson’s model Swanson:2003tb . In such models, calculations using the most straightforward numerical methods tend to become increasingly unstable as the binding energy is tuned toward 0, because the small binding energy results from a delicate cancellation. The factorization formula in Eqs. (45) and (46) can be useful for extrapolating the predictions of a model to other values of the binding energy $`E_X`$. In many models, it is difficult to take into account effects of the width $`\mathrm{\Gamma }_X`$ of the molecule. Given the prediction of a model in which the width has been neglected, the factorization formula in Eqs. (45) and (46) can be used to take into account the nonzero width $`\mathrm{\Gamma }_X`$ consistently.
In order to use the factorization formula in Eqs. (45) and (46) to extrapolate a partial width calculated using a model to other values of the binding energy and the width, the calculation must be carried out for small enough binding energy that the model is in the universal scaling regime where observables scale as powers of the binding energy. For example, the probabilities for components of the wavefunction other than $`D^0\overline{D}^0`$ and $`D^0\overline{D}^0`$ should scale as $`E_X^{1/2}`$. In Swanson’s model with $`E_X=0.7`$ MeV, this universal scaling behavior is satisfied by the probability $`Z_{\psi \omega }`$ but not by $`Z_{\psi \rho }`$, as is evident in Fig. 1. The delayed onset of the universal behavior for the probability $`Z_{\psi \rho }`$ can perhaps be attributed to the weaker coupling of $`X`$ to isospin-1 states. In the next section, we will use Swanson’s result for $`Z_{\psi \omega }`$ to estimate the coupling constant $`G_{X\psi \omega }`$.
## V Partial width for $`𝑿\mathbf{}𝑱\mathbf{/}𝝍𝝅^\mathbf{+}𝝅^{\mathbf{}}𝝅^\mathrm{𝟎}`$
The partial widths of the $`X`$ calculated in Section IV are expressed in terms of unknown coupling constants $`G_{X\psi \rho }`$ and $`G_{X\psi \omega }`$. In this section, we use a simple 2-channel scattering model to show that $`|G_{X\psi \omega }|`$ can be deduced from the probability $`Z_{\psi \omega }`$ for the $`J/\psi \omega `$ component of the $`X`$. We then use the probability $`Z_{\psi \omega }`$ in Swanson’s model to give a quantitative prediction for the partial width for the $`X`$ to decay into $`J/\psi \pi ^+\pi ^{}\pi ^0`$.
### V.1 Two-channel scattering model
Cohen, Gelman, and van Kolck have constructed a renormalizable effective field theory that describes two scattering channels with S-wave contact interactions Cohen:2004kf . We will refer to this model as the two-channel scattering model. An essentially equivalent model has been used to describe the effects of $`\mathrm{\Delta }\mathrm{\Delta }`$ states on the two-nucleon system Savage:1996tb . The parameters of this model can be tuned to produce a large scattering length in the lower energy channel. It can be used as a simple model for the effects on the $`D^0\overline{D}^0`$/$`D^0\overline{D}^0`$ system of other hadronic channels with nearby thresholds, such as $`J/\psi \rho `$ and $`J/\psi \omega `$.
The two-channel model of Ref. Cohen:2004kf describes two scattering channels with S-wave contact interactions only. We label the particles in the first channel $`1a`$ and $`1b`$ and those in the second channel $`2a`$ and $`2b`$. We denote the reduced masses in the two channels by $`\mu `$ and $`\mu _2`$. Renormalized observables in the 2-body sector are expressed in terms of 4 parameters: three interaction parameters $`a_{11}`$, $`a_{22}`$, and $`a_{12}=a_{21}`$ with dimensions of length and the energy gap $`\mathrm{\Delta }`$ between the two scattering channels, which is determined by the masses of the particles:
$`\mathrm{\Delta }`$ $`=`$ $`m_{2a}+m_{2b}(m_{1a}+m_{1b}).`$ (47)
The scattering parameters in Ref. Cohen:2004kf were defined in such a way that $`a_{11}`$ and $`a_{22}`$ reduce in the limit $`a_{12}\pm \mathrm{}`$ to the scattering lengths for the two channels. The truncated connected transition amplitude $`𝒜(E)`$ for this coupled-channel system is a $`2\times 2`$ matrix that depends on the energy $`E`$ in the center-of-mass frame. If that energy is measured relative to the threshold $`m_{1a}+m_{1b}`$ for the first scattering channel, the inverse of the matrix $`𝒜(E)`$ is
$`𝒜(E)^1={\displaystyle \frac{1}{2\pi }}\left(\begin{array}{cc}\mu \left[1/a_{11}+\sqrt{2\mu E}\right]& \sqrt{\mu \mu _2}/a_{12}\\ \sqrt{\mu \mu _2}/a_{12}& \mu _2\left[1/a_{22}+\sqrt{2\mu _2(\mathrm{\Delta }E)}\right]\end{array}\right).`$ (48)
The square roots are defined for negative real arguments by the prescription $`EE+iϵ`$ with $`ϵ0^+`$. The explicit expressions for the $`11`$ and $`12`$ entries of this matrix are
$`𝒜_{11}(E)`$ $`=`$ $`{\displaystyle \frac{2\pi }{\mu }}\left({\displaystyle \frac{1}{a_{11}}}+\sqrt{2\mu E}{\displaystyle \frac{1}{a_{12}^2}}\left[1/a_{22}+\sqrt{2\mu _2(\mathrm{\Delta }E)}\right]^1\right)^1,`$ (49a)
$`𝒜_{12}(E)`$ $`=`$ $`{\displaystyle \frac{2\pi }{\sqrt{\mu \mu _2}}}\left({\displaystyle \frac{1}{a_{12}}}a_{12}\left[{\displaystyle \frac{1}{a_{11}}}+\sqrt{2\mu E}\right]\left[{\displaystyle \frac{1}{a_{22}}}+\sqrt{2\mu _2(\mathrm{\Delta }E)}\right]\right)^1.`$ (49b)
The amplitudes defined by (48) are for transitions between states with the standard nonrelativistic normalizations. The transitions between states with the standard relativistic normalizations are obtained by multiplying by a factor $`\sqrt{2m_i}`$ for every particle in the initial and final state. The T-matrix element $`T_{11}(p)`$ for the elastic scattering of particles in the first channel with relative momentum $`p`$ is obtained by evaluating $`𝒜_{11}(E)`$ at the energy $`E=p^2/(2\mu )`$. The scattering length is determined by the T-matrix element at $`p=0`$: $`T_{11}(0)=2\pi a/\mu `$. The inverse scattering length $`1/a`$ is therefore
$`{\displaystyle \frac{1}{a}}`$ $`=`$ $`{\displaystyle \frac{1}{a_{11}}}+{\displaystyle \frac{1}{a_{12}^2}}\left[\sqrt{2\mu _2\mathrm{\Delta }}1/a_{22}\right]^1.`$ (50)
If the matrix $`𝒜(E)`$ given by Eq. (48) has a pole on the physical sheet at $`E=\kappa ^2/(2\mu )`$, there is a bound state below the scattering threshold for the first channel with binding energy $`E_X=\kappa ^2/(2\mu )`$. The binding momentum $`\kappa `$ satisfies
$`\kappa ={\displaystyle \frac{1}{a_{11}}}+{\displaystyle \frac{1}{a_{12}^2}}\left[1/a_{22}+\sqrt{2\mu _2\mathrm{\Delta }+(\mu _2/\mu )\kappa ^2}\right]^1.`$ (51)
The momentum-space wavefunction $`\psi (p)`$ for the bound state is a column vector whose two components are the amplitudes for the bound state to consist of particles with relative momentum $`p`$ in the first and second channel, respectively. The wavefunction can be deduced from the behavior of $`𝒜(E)`$ near the bound-state pole:
$`𝒜(E){\displaystyle \frac{1}{E+\kappa ^2/(2\mu )}}\left(\begin{array}{c}𝒜_{X1}\\ 𝒜_{X2}\end{array}\right)\left(\begin{array}{cc}𝒜_{X1}& 𝒜_{X2}\end{array}\right).`$ (52)
The components $`𝒜_{X1}`$ and $`𝒜_{X2}`$ of the column vector are the amplitudes for transitions from the bound state to particles in the first and second channels, respectively. They satisfy
$`\mu [1/a_{11}+\kappa ]𝒜_{X1}+[\sqrt{\mu \mu _2}/a_{12}]𝒜_{X2}=0.`$ (53)
Because the only interactions in the two-channel model are contact interactions, the dependence of the wavefunction on the relative momentum of the constituents comes only from propagators. The wavefunction can be expressed in the form
$`\psi (p)=N\left(\begin{array}{c}2\mu 𝒜_{X1}[p^2+\kappa ^2]^1\\ 2\mu _2𝒜_{X2}[p^2+2\mu _2\mathrm{\Delta }+(\mu _2/\mu )\kappa ^2]^1\end{array}\right),`$ (54)
where $`N`$ is a normalization constant. The normalization condition
$`{\displaystyle \frac{d^3p}{(2\pi )^3}\left(|\psi _1(p)|^2+|\psi _2(p)|^2\right)}=1`$ (55)
can be expressed as $`Z_1+Z_2=1`$, where $`Z_1`$ and $`Z_2`$ are the probabilities for the bound state to consist of the particles in the first and second channels, respectively. The probability $`Z_1`$ for the first channel is given by
$`Z_1^1=1+{\displaystyle \frac{(\mu _2/\mu )a_{12}^2(1/a_{11}+\kappa )^2\kappa }{\sqrt{2\mu _2\mathrm{\Delta }+(\mu _2/\mu )\kappa ^2}}}.`$ (56)
### V.2 Two-channel model with large scattering length
In the two-channel model of Ref. Cohen:2004kf , a large scattering length $`a`$ in the first channel can be obtained by fine-tuning the parameters $`a_{11}`$, $`a_{22}`$, $`a_{12}`$, and $`\mathrm{\Delta }`$. The natural momentum scale $`\mathrm{\Lambda }`$ associated with low-energy elastic scattering in the first channel is set by the magnitudes of $`a_{11}^1`$, $`a_{22}^1`$, $`a_{12}^1`$, and $`(2\mu _2\mathrm{\Delta })^{1/2}`$. There are various ways to tune the parameters so that $`|a|`$ is large compared to $`\mathrm{\Lambda }^1`$. For example, $`a`$ can be tuned to $`\pm \mathrm{}`$ by tuning the scattering parameter $`a_{11}`$ to the critical value $`a_{12}^2[\sqrt{2\mu _2\mathrm{\Delta }}1/a_{22}]`$.
As $`a`$ is tuned to be much larger than the natural momentum scale, the amplitude $`𝒜_{11}(E)`$ for $`|E|\mathrm{\Lambda }^2/(2\mu )`$ approaches the universal expression given in Eq. (9). The solution to Eq. (51) for the binding momentum $`\kappa `$ approaches $`1/a`$, so if $`a>0`$, there is a bound state with the universal binding energy in Eq. (10). The first component of the wavefunction in Eq. (54) approaches the universal expression in Eq. (11), while the probability of the second component approaches 0 as $`1/a`$. The amplitude $`𝒜_{X1}`$ for the transition from the bound state to particles in the first channel also approaches the universal amplitude $`𝒜_X`$ in Eq. (12).
There are also universal features associated with transitions from the bound state to particles in the second channel. If $`|a|\mathrm{\Lambda }^1`$, the leading term in the amplitude for the transition of the weakly-bound state $`X`$ to particles in the second channel is
$`𝒜_{X2}={\displaystyle \frac{\sqrt{\mu /\mu _2}}{a_{12}}}\left[\sqrt{2\mu _2\mathrm{\Delta }}1/a_{22}\right]^1𝒜_X,`$ (57)
where $`𝒜_X`$ is the universal amplitude given in Eq. (12). This equation is a factorization formula that expresses the transition amplitude as the product of a short-distance factor and the universal long-distance factor $`𝒜_X`$. Using Eq. (56), the probability $`Z_2=1Z_1`$ for the bound state to consist of particles in the second channel reduces to
$`Z_2={\displaystyle \frac{(\mu _2/\mu )}{a_{12}^2\sqrt{2\mu _2\mathrm{\Delta }}}}\left[\sqrt{2\mu _2\mathrm{\Delta }}1/a_{22}\right]^2{\displaystyle \frac{1}{a}}.`$ (58)
Note that the probability $`Z_2`$ differs from $`|𝒜_{X2}|^2`$ only by kinematic factors:
$`|𝒜_{X2}|^2=\sqrt{8\pi ^2\mathrm{\Delta }/\mu _2^3}Z_2.`$ (59)
This relation also follows directly from the wavefunction in Eq. (54) if we use the fact that the normalization factor $`N`$ approaches 1 as $`a\mathrm{}`$. Thus the relation between the probability and the transition amplitude in Eq. (59) is not specific to the 2-channel model. It applies more generally to any 2-particle component of the bound state whose wavefunction can be approximated by $`(p^2+2\mu _2\mathrm{\Delta })^1`$, where $`\mathrm{\Delta }`$ is the energy gap. It requires only that $`\mathrm{\Delta }`$ is small enough that the interaction in that channel can be approximated by an S-wave contact interaction at momenta comparable to $`\sqrt{2\mu _2\mathrm{\Delta }}`$.
### V.3 Partial width into $`𝑱\mathbf{/}𝝍𝝅^\mathbf{+}𝝅^{\mathbf{}}𝝅^\mathrm{𝟎}`$
We can use results from Swanson’s model to estimate $`|G_{X\psi \omega }|`$, thereby determining the unknown constant in the expression in Eq. (34) for the partial width for $`XJ/\psi \pi ^+\pi ^{}\pi ^0`$. The relativistic amplitude for the transition from $`X`$ to $`J/\psi \omega `$ is given by the contraction of the amplitude $`𝒜_\mu `$ in Eq. (29) with a polarization vector for the $`\omega `$. The corresponding nonrelativistic amplitude $`𝒜_{X\psi \omega }`$ is the analog of the transition amplitude $`𝒜_{X2}`$ in Eq. (57) for the 2-channel model. In the rest frame of the $`X`$, the relativistic amplitude differs from the nonrelativistic amplitude by a factor of $`\sqrt{2m_i}`$ for every external particle:
$`ϵ_{\omega }^{\mu }{}_{}{}^{}𝒜_\mu [XJ/\psi \omega ]=(8m_Xm_\psi m_\omega )^{1/2}𝒜_{X\psi \omega }.`$ (60)
Using the expression for the amplitude $`𝒜_\mu `$ in Eq. (29) and the fact that the rest frame of the $`X`$ is almost identical to that of the $`\omega `$, the left side of Eq. (60) is
$`ϵ_{\omega }^{\mu }{}_{}{}^{}𝒜_\mu [XJ/\psi \omega ]=G_{X\psi \omega }m_\omega \mathit{ϵ}_X(\mathit{ϵ}_\psi \times \mathit{ϵ}_\omega )^{}.`$ (61)
The transition amplitude $`𝒜_{X\psi \omega }`$ on the right side of Eq. (60) must have the same dependence on the polarization vectors of the $`J/\psi `$ and $`\omega `$. There are two independent pairs of spin states for $`J/\psi `$ and $`\omega `$ that couple to any given spin state of $`X`$. If the analog of the factorization formula in Eq. (59) is summed over the spin states of the $`J/\psi `$ and $`\omega `$, it gives
$`{\displaystyle \underset{\mathrm{spins}}{}}|𝒜_{X\psi \omega }|^2=\sqrt{8\pi ^2\mathrm{\Delta }_{\psi \omega }/\mu _{\psi \omega }^3}Z_{\psi \omega },`$ (62)
where $`\mathrm{\Delta }_{\psi \omega }`$ is the energy gap in Eq. (20c), $`\mu _{\psi \omega }`$ is the reduced mass of the $`J/\psi `$ and $`\omega `$, and $`Z_{\psi \omega }`$ is the probability for the $`J/\psi \omega `$ component of $`X`$. Squaring both sides of Eq. (60) and summing over the spin states of $`J/\psi `$ and $`\omega `$, we get
$`2m_\omega ^2|G_{X\psi \omega }|^2=16\pi m_X(m_\psi +m_\omega )\sqrt{2\mathrm{\Delta }_{\psi \omega }/\mu _{\psi \omega }}Z_{\psi \omega }.`$ (63)
Inserting Swanson’s result $`Z_{\psi \omega }=9.6\%`$ for $`E_X=0.7`$ MeV and using the factorization formula in Eqs. (45) and (46), this reduces to
$`|G_{X\psi \omega }|^2=9.59\left({\displaystyle \frac{E_X+\mathrm{\Gamma }_X^2/(16E_X)}{0.7\mathrm{MeV}}}\right)^{1/2}.`$ (64)
Inserting the result into the expression in Eq. (34), we get a quantitative result for the partial width:
$`\mathrm{\Gamma }[XJ/\psi \pi ^+\pi ^{}\pi ^0]`$ $`=`$ $`(222\mathrm{keV})\left({\displaystyle \frac{E_X+\mathrm{\Gamma }_X^2/(16E_X)}{1\mathrm{MeV}}}\right)^{1/2}.`$ (65)
We can use the result in Eq. (65) to set a lower bound on the partial width into $`J/\psi \pi ^+\pi ^{}\pi ^0`$. As a function of the binding energy $`E_X`$, the right side of Eq. (65) is minimized at $`E_X=\mathrm{\Gamma }_X/4`$. The lower bound on the width is $`\mathrm{\Gamma }_X>2\mathrm{\Gamma }[D^0]=136\pm 32`$ keV. Thus the lower bound on the partial width into $`J/\psi \pi ^+\pi ^{}\pi ^0`$ in Swanson’s model is about 58 keV.
As is evident in Fig. 1, Swanson did not calculate the probability $`Z_{\psi \rho }`$ for the $`J/\psi \rho `$ component of $`X`$ for a binding energy small enough to be in the scaling region where $`Z_{\psi \rho }`$ scales like $`E_X^{1/2}`$. If he had, we could use an equation analogous to Eq. (63) to determine $`|G_{X\psi \rho }|`$. If we assume that the smallest binding energy considered by Swanson is close to the scaling region, we can use his value $`Z_{\psi \rho }=0.86\%`$ for $`E_X=0.7`$ MeV to estimate $`|G_{X\psi \rho }|`$. In the analog of Eq. (63), we should set $`\mathrm{\Delta }_{\psi \rho }=\mathrm{\Delta }_{\psi \omega }`$ rather than using the value $`\mathrm{\Delta }_{\psi \rho }`$ in Eq. (20b), because Swanson set $`m_\rho =m_\omega `$ in his calculation. The resulting estimate is
$`|G_{X\psi \rho }|^20.86\left({\displaystyle \frac{E_X+\mathrm{\Gamma }_X^2/(16E_X)}{0.7\mathrm{MeV}}}\right)^{1/2}.`$ (66)
We can insert this estimate into Eq. (32) to get an estimate of the partial width for decay into $`J/\psi \pi ^+\pi ^{}`$. We can also insert this estimate of $`|G_{X\psi \rho }|^2`$ and the value of $`|G_{X\psi \omega }|^2`$ from Eq. (64) into Eqs. (38) and (41) to get ranges of estimates of the partial widths for the decays into $`J/\psi \pi ^0\gamma `$ and $`J/\psi \gamma `$. The ranges arise from the unknown relative phase between $`G_{X\psi \omega }`$ and $`G_{X\psi \rho }`$.
## VI Summary
Evidence is accumulating that the $`X(3872)`$ is a loosely-bound S-wave molecule corresponding to a $`C=+`$ superposition of $`D^0\overline{D}^0`$ and $`D^0\overline{D}^0`$. Because its binding energy is small compared to the natural energy scale associated with pion exchange, this molecule has universal properties that are completely determined by the large scattering length $`a`$ in the $`C=+`$ channel of $`D^0\overline{D}^0`$ and $`D^0\overline{D}^0`$.
We have analyzed the decays of $`X`$ into $`J/\psi `$ plus light hadrons under the assumption that $`X`$ is a $`DD^{}`$ molecule and that these decays proceed through transitions to $`J/\psi \rho `$ and $`J/\psi \omega `$. The differential decay rates were calculated in terms of unknown coupling constants $`G_{X\psi \rho }`$ and $`G_{X\psi \omega }`$ by using an effective lagrangian that reproduces the decays of the light vector mesons. The dependence on the unknown coupling constants enters only through multiplicative factors, so the angular distributions are completely determined.
Quantitative predictions of the partial widths for the decays of $`X`$ into $`J/\psi `$ plus light hadrons require numerical values for the coupling constants $`G_{X\psi \rho }`$ and $`G_{X\psi \omega }`$. We pointed out that the dependence of these coupling constants on the binding energy $`E_X`$ and the total width $`\mathrm{\Gamma }_X`$ are determined by factorization formulas. We showed how $`|G_{X\psi \omega }|^2`$ could be determined from the probability $`Z_{\psi \omega }`$ for the $`J/\psi \omega `$ component of $`X`$ in Swanson’s model. We used this result to give a quantitative prediction for the partial width for $`XJ/\psi \pi ^+\pi ^{}\pi ^0`$ as a function of $`E_X`$ and $`\mathrm{\Gamma }_X`$.
## Appendix A Decay amplitudes for vector mesons
In this appendix, we present an updated determination of the coupling constants in the effective lagrangian for the light pseudoscalar and vector mesons that was used in Ref. Braaten:1989zn to calculate the semileptonic branching fractions for the $`\tau `$ lepton. The same effective lagrangian is used in Section IV to calculate the decay rates of the $`X`$ into $`J/\psi `$ and light hadrons.
The pion decay constant $`F_\pi =93`$ MeV and the hadron masses have all been determined accurately Eidelman:2004wy . The other parameters in the effective lagrangian can be determined from the partial widths for decays of $`\rho ^0`$, $`\rho ^\pm `$, and $`\omega `$ given in Table 1. The most useful combinations of the parameters in the amplitudes for the decays of the vector mesons into pions in Eqs. (24) and (25) are
$`G_{v\pi \pi }`$ $`=`$ $`11.99\pm 0.06,`$ (67a)
$`C_{v3\pi }+G_{v\pi \pi }C_{vv\pi }F_\pi ^2/m_v^2`$ $`=`$ $`(8.03\pm 0.48)/(16\pi ^2),`$ (67b)
$`G_{v\pi \pi }C_{vv\pi }F_\pi ^2/m_v^2`$ $`=`$ $`(10.2\pm 1.3)/(16\pi ^2).`$ (67c)
The coupling constant $`G_{v\gamma }`$ associated with vector meson dominance and the most useful combinations of parameters in the amplitudes for the radiative decays of the vector mesons in Eqs. (27) are
$`G_{v\gamma }`$ $`=`$ $`14.01\pm 0.11,`$ (68a)
$`C_{v\pi \gamma }+G_{v\gamma }C_{vv\pi }F_\pi ^2/m_v^2`$ $`=`$ $`(7.99\pm 0.45)/(16\pi ^2),`$ (68b)
$`G_{v\gamma }C_{vv\pi }F_\pi ^2/m_v^2`$ $`=`$ $`(11.9\pm 1.5)/(16\pi ^2).`$ (68c)
The vector meson mixing angle is given by<sup>1</sup><sup>1</sup>1The cosine of the angle $`\theta _v`$ here is the sine of the vector meson mixing angle used in Ref. Braaten:1989zn .
$$\mathrm{cos}\theta _v=0.51\pm 0.01.$$
(69)
Another function of $`\theta _v`$ that is often encountered is $`\mathrm{cos}\theta _v+\sqrt{2}\mathrm{sin}\theta _v1.73\pm 0.01`$. The errors in the parameters in Eqs. (67), (68), and (69) are determined using the uncertainties in the measurements of the vector meson decay widths only. The uncertainties in the hadron masses and the pion decay constant are negligible in comparison. Variations in the parameters associated with $`U(3)\times U(3)`$ symmetry breaking are neglected in this analysis.
The inputs that were used to determine the parameters in Eqs. (67), (68), and (69) are listed in Table 1. Following Ref. Braaten:1989zn , we determine the parameters by the following steps:
1. The coupling constant $`G_{v\pi \pi }`$ in Eq. (67a) is determined from the partial width for $`\rho \pi ^+\pi ^{}`$:
$$\mathrm{\Gamma }[\rho \pi ^+\pi ^{}]=\frac{G_{v\pi \pi }^2m_\rho }{192\pi }\left(14m_\pi ^2/m_\rho ^2\right)^{3/2}.$$
(70)
2. The coupling constant $`G_{v\gamma }`$ in Eq. (68a) is determined from the partial width for $`\rho e^+e^{}`$:
$$\mathrm{\Gamma }[\rho e^+e^{}]=\frac{4\pi \alpha _{\mathrm{em}}^2G_{v\gamma }^2F_\pi ^4}{3m_\rho ^3}.$$
(71)
3. The combination of parameters in Eq. (68b) is determined from the partial width for $`\rho ^{}\pi ^{}\gamma `$:
$$\mathrm{\Gamma }[\rho ^{}\pi ^{}\gamma ]=\frac{2\alpha _{\mathrm{em}}m_\rho ^3}{27F_\pi ^2}\left(C_{v\pi \gamma }+\frac{G_{v\gamma }C_{vv\pi }F_\pi ^2}{m_v^2}\right)^2(1m_\pi ^2/m_\rho ^2)^3.$$
(72)
4. The combination of parameters $`G_{v\gamma }C_{vv\pi }F_\pi ^2/m_v^2`$ in Eq. (68c) is determined from the ratio of the partial widths for $`\omega \pi ^0\mu ^+\mu ^{}`$ and $`\omega \pi ^0\gamma `$. The possibility of a relative phase between $`C_{v\pi \gamma }`$ and $`G_{v\gamma }C_{vv\pi }F_\pi ^2/m_v^2`$ is ignored. The partial width for $`\omega \pi ^0\mu ^+\mu ^{}`$ is
$$\mathrm{\Gamma }[\omega \pi ^0\mu ^+\mu ^{}]=\frac{1}{256\pi ^3m_\omega ^3}𝑑s_{12}𝑑s_{23}\overline{}|𝒜[\omega \pi ^0\mu ^+\mu ^{}]|^2.$$
(73)
The squared amplitude, averaged over initial spin states and summed over final spin states, is
$`\overline{}|𝒜[\omega \pi ^0\mu ^+\mu ^{}]|^2`$ $`={\displaystyle \frac{128\pi ^2\alpha _{\mathrm{em}}^2}{9F_\pi ^2}}(\mathrm{cos}\theta _v+\sqrt{2}\mathrm{sin}\theta _v)^2`$ (74)
$`\times \left[(s_{23}^2+4m_\mu ^2)\left((m_\omega ^2s_{23}m_\pi ^2)^24m_\pi ^2s_{23}\right)+s_{23}(s_{12}s_{31})^2\right]`$
$`\times {\displaystyle \frac{1}{s_{23}^2}}\left|C_{v\pi \gamma }+{\displaystyle \frac{G_{v\gamma }C_{vv\pi }F_\pi ^2}{m_v^2}}(1f_\rho (s_{23}))\right|^2,`$
where $`s_{12}`$, $`s_{23}`$, and $`s_{31}`$ are the squares of the invariant masses of the $`\pi ^0\mu ^+`$, $`\mu ^+\mu ^{}`$, and $`\mu ^{}\pi ^0`$, respectively. The partial width for $`\omega \pi ^0\gamma `$ is
$$\mathrm{\Gamma }[\omega \pi ^0\gamma ]=3(\mathrm{cos}\theta _v+\sqrt{2}\mathrm{sin}\theta _v)^2\frac{m_\omega ^3(1m_\pi ^2/m_\omega ^2)^3}{m_\rho ^3(1m_\pi ^2/m_\rho ^2)^3}\mathrm{\Gamma }[\rho ^{}\pi ^{}\gamma ],$$
(75)
where $`\mathrm{\Gamma }[\rho ^{}\pi ^{}\gamma ]`$ is given in Eq. (72). Note that the factor $`(\mathrm{cos}\theta _v+\sqrt{2}\mathrm{sin}\theta _v)^2`$ cancels in the ratio of Eqs. (73) and (75).
5. The combination of parameters $`G_{v\pi \pi }C_{vv\pi }F_\pi ^2/m_v^2`$ appearing in Eq. (67c) is determined by multiplying the combination of parameters in Eq. (68c) by the ratio $`G_{v\pi \pi }/G_{v\gamma }`$ obtained from Eqs. (67a) and (68a).
6. The combination of parameters in Eq. (67b) is determined from the ratio of the partial widths for $`\omega `$ to decay into $`\pi ^+\pi ^{}\pi ^0`$ and $`\pi ^0\gamma `$ and from the value of the combination of parameters in Eq. (67c). The possibility of a relative phase between $`C_{v3\pi }`$ and $`G_{v\pi \pi }C_{vv\pi }F_\pi ^2/m_v^2`$ is ignored. The partial width for $`\omega \pi ^0\gamma `$ is given in Eq. (75). The partial width for $`\omega \pi ^+\pi ^{}\pi ^0`$ is
$$\mathrm{\Gamma }[\omega \pi ^0\pi ^+\pi ^{}]=\frac{1}{256\pi ^3m_\omega ^3}𝑑s_{12}𝑑s_{23}\overline{}|𝒜[\omega \pi ^+\pi ^{}\pi ^0]|^2.$$
(76)
The squared amplitude, averaged over the spin states of $`\omega `$, is
$`\overline{}|𝒜[\omega \pi ^+\pi ^{}\pi ^0]|^2`$ $`={\displaystyle \frac{4(\mathrm{cos}\theta _v+\sqrt{2}\mathrm{sin}\theta _v)^2}{F_\pi ^6}}\left(s_{12}s_{23}s_{31}m_\pi ^2(m_\omega ^2m_\pi ^2)^2\right)`$ (77)
$`\times \left|C_{v3\pi }+{\displaystyle \frac{G_{v\pi \pi }C_{vv\pi }F_\pi ^2}{m_v^2}}\left(1\frac{1}{3}\left[f_\rho (s_{12})+f_\rho (s_{23})+f_\rho (s_{31})\right]\right)\right|^2.`$
Note that the factor $`(\mathrm{cos}\theta _v+\sqrt{2}\mathrm{sin}\theta _v)^2`$ cancels in the ratio of Eqs. (76) and (75).
7. Finally, the cosine of the vector meson mixing angle in Eq. (69) is determined from the ratio of the partial widths for $`\omega e^+e^{}`$ and $`\rho ^0e^+e^{}`$:
$$\mathrm{\Gamma }[\omega e^+e^{}]=\frac{\mathrm{cos}^2\theta _vm_\rho ^3}{3m_\omega ^3}\mathrm{\Gamma }[\rho ^0e^+e^{}].$$
(78)
###### Acknowledgements.
We thank Mark Wise for useful discussions. We thank Eric Swanson for providing us with some of his numerical results. This research was supported in part by the Department of Energy under grant DE-FG02-91-ER4069. |
warning/0507/math0507355.html | ar5iv | text | # Problems on Bieberbach groups and flat manifolds
## 1 The classification problems
It is well known (cf. ) that any finite group $`G`$ is a holonomy group of some flat manifold. Hence we have.
###### Problem 1
Find the minimal dimension of a flat manifold with the holonomy group $`G.`$
###### Remark 1
The answer is known for cyclic groups (, ), elementary abelian $`p`$ \- groups, dihedral groups, semidihedral groups, generalized quaternion groups (see ) and simple groups $`PSL(2,p)`$ ($`p`$ is a prime number) . There are also some other classes of finite groups for which the minimal dimension is known. We do not pretend to give a complete list here. Unfortunately, for the most finite groups it seems to be very difficult question. For example, almost nothing is known for the symmetric groups. The main obstruction is the calculation of the second cohomology of the finite group with special coefficients.
###### Problem 2
Do the calculation for those $`p`$-groups for which the cohomology is well-understood.
An inductive classification of the Bieberbach groups follows from E. Calabi (, \[35, Section 3.6 on page 124\]). This is a consequence of the fact that any subgroup of the torsion free crystallographic group is again crystallographic. Hence any $`n`$-dimensional Bieberbach group $`\mathrm{\Gamma }`$ with non trivial center is an extension of some Bieberbach group $`\mathrm{\Gamma }^{}`$ of dimension $`n1`$ by integers $``$,
$$0\mathrm{\Gamma }^{}\mathrm{\Gamma }0.$$
Hence we have the following problem.
###### Problem 3
For a finite group $`G`$ find a Bieberbach group (with trivial center) of minimal dimension.
###### Remark 2
The problem was considered in (see also , ), where there are also some calculations. These include cyclic groups, elementary abelian $`p`$-groups (where $`p`$ is a prime number), dihedral $`2`$-groups, semidihedral $`2`$-groups, the generalized quaternion $`2`$-groups and simple groups of Lie type.
In 1972 A. T. Vasquez introduced, for any finite group $`G,`$ an invariant $`n(G).`$ This is related to the class of flat manifolds with holonomy group $`G.`$ In 1989 in the beautiful paper G. Cliff and A. Weiss , calculated $`n(G)`$ for any finite $`p`$-group. A purely algebraic definition of $`n(G)`$ is given in . There is also given a characterization of the groups with $`n(G)=1.`$
###### Problem 4
Calculate Vasquez invariant $`n(G)`$ for a finite group $`G.`$
###### Remark 3
There is a method, complementary to that of Calabi for the classification of Bieberbach groups with given holonomy group $`G.`$ It says that any Bieberbach group $`\mathrm{\Gamma }`$ with holonomy group $`G`$ and dimension $`nn(G)`$ can be defined by a short exact sequence of groups
$$0^{nn(G)}\mathrm{\Gamma }\mathrm{\Gamma }_G0,$$
where $`\mathrm{\Gamma }_G`$ is a Bieberbach group of dimension $`n(G)`$, (cf. ).
Let $`f:M^nM^n`$ be a continuous map, and let $`\stackrel{~}{f}:^n^n`$ be its cover in the Euclidean space. From the third Bieberbach theorem (see ), we have the following short exact sequence (cf. )
$$0Aff_0(M^n)Aff(M^n)Out(\mathrm{\Gamma })0,$$
where
$$Aff(M^n)=\{f:M^nM^n\stackrel{~}{f}:^n^nGL(n,)^n\}.$$
$`Aff_0(M^n)`$ denotes the connected component of the identity of the group of affine diffeomorphisms of the flat manifold $`M^n.`$ It is isomorphic (see ) to $`(S^1)^{\beta _1(M^n)},`$ where $`\beta _1`$ denotes the first Betti number of the manifold, see . Moreover we have the following a short exact sequence (see and )
$$0H^1(H,^n)Out(\mathrm{\Gamma })N_\alpha /H0,$$
where $`N=N_{GL(n,)}(H)`$ is the normalizer of the holonomy group $`H`$ in $`GL(n,)`$ and
$$N_\alpha =\{nNn\alpha =\alpha \}.$$
Here $``$ is the standard action of the normalizer in cohomology, see \[4, page 168\]. Hence we have.
###### Theorem 1
() The following conditions are equivalent: (i) $`Out(\mathrm{\Gamma })`$ is a finite group, (ii) the normalizer $`N`$ is a finite group, (iii) the holonomy representation $`\varphi _\mathrm{\Gamma }`$ is $``$-multiplicity free and any $``$-irreducible component is also $``$-irreducible.
Proof: The equivalence of the conditions $`(i)`$ and $`(ii)`$ is obvious. The proof of the last equivalence is much more difficult and we refer the reader to .
We can now fomulate the following,
###### Problem 5
Classify finite groups which are holonomy groups of flat manifolds $`M^n=^n/\mathrm{\Gamma }`$ with finite outer automorphism group Out$`\mathrm{\Gamma }.`$
###### Remark 4
Some progress was made in , e.g. for finite abelian groups.
From a finite group representation theory point of view, it might be interesting to consider a similar classification for Kähler flat manifolds (see ). Note that, here $`Out(\mathrm{\Gamma })`$ is finite if and only if the holonomy representation is $``$-multiplicity free and any $``$-irreducible component is $``$-reducible. The next conjecture is related to the last problem.
###### Conjecture 1
() For any finite group $`H,`$ there exist a Bieberbach group $`\mathrm{\Gamma }`$ with $``$ multiplicity free holonomy representation $`\varphi _\mathrm{\Gamma }.`$
###### Remark 5
If the first Betti number of the manifold $`M^n`$ is zero then, from above, we have
$$Aff(M^n)Out(\mathrm{\Gamma }).$$
This means that all information about the symmetries of the manifild is built into the outer automorphism group.
It is interesting to compare the last problem with the work of M. Belolipetsky and A. Lubotzky where they proved that for any finite group $`G`$ there exist a fundamental group $`\mathrm{\Gamma }`$ of some compact hyperbolic manifold such that $`Out(\mathrm{\Gamma })=G.`$ Hence we have.
###### Problem 6
Which finite groups $`G`$ occur as outer automorphism groups of Bieberbach groups with trivial center.
In this connection we can ask if every finite subgroup of $`GL(n,)`$ has a realisation as the centraliser of some finite subgroup of $`GL(n,),`$ for some $`n.`$
###### Remark 6
In there is an example of a flat manifold whose fundamental group has trivial center and trivial outer automorphism group.
## 2 The Anosov relation for flat manifolds and Entropy Conjecture
Let $`f:MM`$ be a continuous map of a smooth manifold $`M.`$ In fixed point theory, there exist two numbers associated to $`f`$ which provide some information on the fixed point set of $`f:`$ the Lefschetz number $`L(f)`$ and the Nielsen number $`N(f).`$ It is known that the Nielsen number provides more information about the fixed point set than $`L(f)`$ does, but $`L(f)`$ is easier to calculate. For nilmanifolds however, D. Anosov showed that $`N(f)=L(f)`$ for any self map $`f`$ of the nilmanifold. On the other hand, he also observed that this result could not be extended to the class of all infranilmanifolds, since he was already able to construct a counter-example on the Klein Bottle. It was recently shown that for large families of flat manifolds (e.g. all flat manifolds with an odd order holonomy group) the Anosov relation $`N(f)=L(f)`$ does hold for any self map (see and ). It is therefore natural to ask.
###### Problem 7
(K. Dekimpe) Describe the class of flat manifolds (or more generally the class of infra-nilmanifolds) on which the Anosov relation
$$N(f)=L(f)$$
does hold for any map $`f:MM.`$
One can also consider this problem from a different point of view and fix the flat manifold $`M`$ and try to find all self maps $`f`$ for which the Anosov relation is valid.
###### Problem 8
(K. Dekimpe) Given a flat manifold $`M.`$ Is it possible to determine classes of self maps $`f:MM`$ for which the Anosov relation
$$N(f)=L(f)$$
does hold ?
A first result in this direction can be found in .
Let $`n^+`$ and $`ϵ^{>0}`$. Moreover let $`(X,d)`$ be a compact metric space and $`f:XX`$ a continuous map. Put
$`r(f,ϵ)`$ := lim<sub>n→∞</sub>sup$`\frac{1}{n}`$ log max$`\{\mathrm{\#}Q:QX\},`$
where $`Q`$ is such that for any distinct points $`x,y`$
max$`{}_{0jn1}{}^{}d(f^j(x),f^j(y))ϵ.`$
The topological entropy $`h(f)`$ is a nonnegative real number or $`\mathrm{}`$ defined
$`h(f)`$ = lim$`{}_{ϵ0}{}^{}r(f,ϵ)`$ = sup$`{}_{ϵ0}{}^{}r(f,ϵ).`$
Now assume that $`X`$ is a compact smooth manifold $`M`$ of dimension $`m.`$ Let
$$H^{}(f):H^{}(M,)H^{}(M,)$$
be the linear map induced by $`f`$ on the cohomology space $`H^{}(M,):=_{i=0}^mH^i(M,)`$ of $`M`$ with real coefficients. By $`sp(f)`$ we denote the spectral radius of $`H^{}(f),`$ which is a homotopy invariant. In 1974 M.Shub asked in \[26, page 36\] the extent to which the following inequality holds.
(EC) log sp($`f)h(f).`$
Then A. Katok in put the following Entropy Conjecture.
###### Conjecture 2
(EC) holds for all continuous mappings $`f:MM,`$ if $`M`$ is a manifold with the universal cover homeomorphic to $`^n.`$
It was proved for $`M`$ being tori and nilmanifolds, see .
###### Problem 9
() Give a proof of the (EC) for flat manifolds.
For any flat manifold $`M`$ with the first Betti number equal to zero and finite outer automorphism group of the fundamental group $`\pi _1(M)`$ the above conjecture is true. In fact, (see remark 5 and theorem 1) for any continuous map $`f:MM`$ the induced map $`H^{}(f)`$ has a finite order. We send the reader to for more informations about the Entropy Conjecture. $``$ We would like to thank W. Marzantowicz for putting our attention on this problem.
## 3 Generalized Hantzsche-Wendt flat manifolds
Let us agree to call any fundamental group of the $`n`$-dimensional flat manifold with holonomy group $`(_2)^{n1}`$ a generalized Hantzsche-Wendt $`(GHW)`$ Bieberbach group (see ). If the manifold is orientable, we call the $`GHW`$ Bieberbach group orientable. For dimension $`3`$ there is only one oriented flat manifold $`M^3`$ with holonomy group $`_2_2.`$ It was first considered by Hantzsche and Wendt. From the other side the fundamental group $`\pi _1(M^3)`$ is group $`F(2,6),`$ where $`F(r,n)`$ is the group defined by the presentation
$$<a_1,a_2,\mathrm{},a_na_1a_2\mathrm{}a_r=a_{r+1},a_2a_3\mathrm{}a_{r+1}=a_{r+2},\mathrm{},$$
$$a_{n1}a_na_1\mathrm{}a_{r2}=a_{r1},a_na_1a_2\mathrm{}a_{r1}=a_r>,$$
where $`r>0,n>0,`$ and all subscripts are assumed to be reduced modulo $`n`$. This group is named a Fibonacci group.
###### Problem 10
Explain the relation between orientable GHW Bieberbach groups and the Fibonacci groups.
###### Remark 7
For any even natural number $`2n`$ there exists an epimorphism
$$\mathrm{\Phi }_{2n}:F(2n,2(2n+1))\mathrm{\Gamma }_{2n},$$
where $`\mathrm{\Gamma }_{2n}E(2n+1)`$ denote an orientable $`GHW`$ Bieberbach group of rank $`(2n+1),`$ see . Nothing is known about the kernel of the epimorphism $`\mathrm{\Phi }_{2n},`$ for $`n>1.`$
###### Problem 11
Give definition of any GHW Bieberbach group in terms of generators and relations.
Let $`\mathrm{\Gamma }`$ be a GHW Bieberbach group of dimension $`n`$ and with trivial center. It is well known (\[10, Theorem 1.1, page 183\], ) that there is an epimorphism of $`\mathrm{\Gamma }`$ onto the infinite dihedral group $`_2_2.`$ Hence we have a decomposition
$$()\mathrm{\Gamma }\mathrm{\Gamma }_1_X\mathrm{\Gamma }_2,$$
as the generalized free product of two Bieberbach groups of dimension $`(n1),`$ amalgamated over a subgroup $`X`$ of index two. Moreover any GHW Bieberbach group ,with non-trivial center, is the semidirect product of $``$ with some lower dimensional GHW Bieberbach group. We have.
###### Conjecture 3
For any GHW Bieberbach group $`\mathrm{\Gamma }`$ with trivial center and dimension $`n`$, there exists a decomposition $`()`$ such that $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ are $`GHW`$ Bieberbach groups of dimension $`n1.`$
###### Remark 8
The conjecture is true for $`n4`$ (see \[14, page 30 and 38\]).
For the remaining problem we have: let $`M^n`$ be a flat oriented manifold with GHW fundamental Bieberbach group and dimension $`n`$. It is well known, see , that $`M^n`$ is a rational homology sphere.
###### Problem 12
What are the topological (geometrical) properties of $`M^n`$ ?
## 4 Flat manifolds and other geometries
Let us recall the following well known question. Question (Farrell-Zdravkovska) ) Let $`M^n`$ be a $`n`$-dimensional flat manifold. Is there a $`(n+1)`$-dimensional hyperbolic manifold $`W`$ the boundary of which equals $`M^n`$ ? We would like to mention that the above question is very close related to the following one. Let $`V^{n+1}`$ be a hyperbolic Riemannian manifold (constant negative curvature) of finite volume. It is well known that $`V^{n+1}`$ has finite number of cusps and each cusp is topologicaly $`M^n\times ^0,`$ where $`M^n`$ is $`n`$-dimensional flat manifold. Is there some $`V^{n+1}`$ with one cusp homeomorphic to $`M^n\times ^0`$ ? D.D.Long and A.Reid proved () that the answer to the problem above is negative if the $`\eta `$-invariant of the signature operator of $`M^n`$ is not an integer. To prove it they used the Atiyach, Patodi, Singer formula, where signature operator $`D`$ means some differential elliptic operator and its $`\eta `$-invariant measure the symmetry of the spectrum of $`D.`$ They also proved that there exists, already in dimension three, a flat manifold which signature operator has the $`\eta `$-invariant $`.`$ By our assumptions this method works only in dimension $`4k1.`$
###### Problem 13
Classify flat manifolds with non-integral signature $`\eta `$-invariant.
They also proved () that any flat manifold has a ”realization” as a cusp of some hyperbolic orbifold of finite volume.
###### Remark 9
We have the following flat manifolds of dimension two: the torus and the Klein bottle. The first one has a ”realization” as a cusp of the eight knot complement and the second one as a cusp of the Gieseking manifold, (c.f. ). In dimension three we have ten flat manifolds and each one has ”realization” as a cusp of some four dimensional hyperbolic manifold $`W,`$ which is not necessarly one-cusped, (see ).
One major difficulty is the lack of examples of hyperbolic manifolds. In particular we have.
###### Problem 14
Give an example of a four dimensional hyperbolic manifold of finite volume with only one cusp.
The $`\eta `$-invariant of the signature operator for a flat manifold is an important quantity as we have observed above. Let us ask the same question about the Dirac operator. We shall need our orientable flat manifold $`M^n`$ to have a spin structure. It turns out that this is equivalent to the existence of a homomorphism $`ϵ:\mathrm{\Gamma }Spin(n)`$ such that the following diagram is commutative
$$Spin(n)$$
$$ϵ\lambda _n$$
$$\mathrm{\Gamma }\stackrel{r}{}SO(n),$$
where $`\lambda _n`$ is the universal covering and $`r`$ is the projection onto the linear part, .
###### Problem 15
Classify the holonomy groups of flat manifolds which admit a spin structure
###### Remark 10
It is known \[8, Proposition 1, Corollary 1\] that any flat manifold with holonomy group of odd order or of order $`2k`$ , $`k`$ an odd number, has a spin structure. However, see \[21, Theorem 3.2\], for the holonomy group $`_2\times _2,`$ there exist orientable flat manifolds $`M_1,M_2`$ of dimension $`6`$ with different holonomy representations $`h_1,h_2:_2\times _2GL(6,),`$ only one of which has a spin structure.
Question Is there an example of an oriented flat spin - manifold for which the Dirac $`\eta `$-invariant is not equal to the signature $`\eta `$-invariant modulo $``$ ? For the methods of calculation of the Dirac $`\eta `$-invariant see .
A six dimensional flat manifold is Calabi-Yau if and only if its holonomy representation has the property that each $``$-irreducible summand, which is also $``$-irreducible, occurs with even multiplicity. Hence the classification of such manifolds is possible, see .
###### Problem 16
What are the properties of the Calabi-Yau flat manifolds. Do they have some ”mirror-symmetry” ?
Acknowledgements This work was carried out in the framework of the project BIL01/C-31 for Bilateral Scientific Cooperation (Flanders–Poland) and was supported by University of Gdańsk BW - 5100-5-0149-4.
Institute of Mathematics
University of Gdańsk
ul. Wita Stwosza 57
80-952 Gdańsk
Poland
E-mail: matas @ paula.univ.gda.pl |
warning/0507/cond-mat0507539.html | ar5iv | text | # Actively Tuned and Spatially Trapped Polaritons
## Abstract
Abstract. We report active tuning of the polariton resonance of quantum well excitons in a semiconductor microcavity using applied stress. Starting with the quantum well exciton energy higher than the cavity photon mode, we use stress to reduce the exciton energy and bring it into resonance with the photon mode. At the point of zero detuning, line narrowing and strong increase of the photoluminescence are seen. By the same means, we create an in-plane harmonic potential for the polaritons, which allows trapping, potentially making Bose-Einstein condensation of polaritons analogous to trapped atoms possible. We demonstrate drift of the polaritons into this trap.
Microcavity polaritons have in the past decade been the object of great interest for many scientistsyama ; deveaud ; kav ; dang ; baum ; ciuti ; little ; gia ; science interested in the study of Bose-Einstein condensation (BEC). These particles, which are mixed states of photons and excitons, have a very light mass, which in principle allows them to condense at critical temperatures near room temperature. For two-dimensional bosonic systems to truly condense at finite temperatures, the application of potential traps or confinement in a region of finite size is essential.mullin Here we present a method to actively couple the exciton mode to the cavity mode at fixed $`k_{||}=0`$ and at the same time create an in-plane spatial trap for both the lower and upper polaritons. The polariton photoluminescence (PL) jumps up dramatically at resonance, and both the PL and the reflectivity show line narrowing as the system approaches resonance.
The sample studied consists of three sets of four GaAs/AlAs quantum wells embedded in a GaAs/AlGaAs microcavity, with each set of quantum wells at an antinode of the confined mode, similar to the structure used in previous work.yama The cavity is designed in such a way that it is initially negatively detuned, with $`\delta 20`$ meV ($`\delta =E_{\mathrm{cav}}E_{\mathrm{ex}})`$. A force is applied on back side of the 150 $`\mu `$m thick substrate with a rounded-tip pin, with approximately 50 $`\mu `$m tip radius, as shown in Figure 1. This pushes the exciton energy down toward the cavity mode, at the same time creating a harmonic potential, following the method published previouslyapl ; ssc for excitons. The harmonic potential is centered in the plane of at the point of pin-sample contact.
Figure 2 shows the reflectivity spectrum as a function of position on the sample, showing the anticrossing of the upper and lower polariton branches as the cavity length is varied, due to the thinning of the layer thickness by about 10% toward the edge of the wafer, which is part of the growth process. The pin stress point is chosen several millimeters to the right of the crossover point, the point of strongest coupling. Figures 3 and 4 show photoluminescence and reflectivity data for a sequence of increasing stresses applied to this sample. For the photoluminescence, a helium-neon laser source (633 nm) is used to excite the sample off-resonantly, well above the band gap, at $`\theta =12^{}`$ incidence, and defocused to a spot size of several millimeters to cover the entire region of observation. Photoluminescence emission collected normal to the sample is directed to a spectrometer and captured with a Photometrics back-illuminated CCD camera. For the sample reflectivity, a collimated light beam (750 nm$``$1000 nm) is directed normal to the sample. The reflected light is also collected normal to the sample. The mirror placed in the same plane as the sample is used to normalize the sample reflectance. For all the experiments, the sample was maintained at the temperature of 4.2 K. At this low temperature, no luminescence is seen from the upper polariton. Upper polariton emission for this sample starts to appear at about 40 K.
As seen in these figures, a harmonic potential for both the upper and lower polaritons is created. The polaritons are clearly in the strong coupling regime, since if they were not, only the exciton states would respond to stress; the stress has negligible effect on the dielectric constants of the materials and therefore negligible effect on the cavity mode in the weak coupling limit. The energy gap between the upper and lower polariton branches decreases, while the overall energy shifts lower due to the band gap reduction.
In addition to the energy shift of the bands, a striking increase of the photoluminescence occurs, as seen in Figure 4. This is similar to the increase of photoluminescence at resonance seen by tuning of the resonance using a wedge of varying cavity thickness,stanley but the increase in the present case is dramatic, a factor of about 100. The increase of the total photoluminescence emitted from the front surface is consistent with an increase of the coupling constant at resonance. Consistent with the strong coupling, one can see in Figure 4 the narrowing of the reflectivity spectra as the bare excitons and bare photon modes approach resonanceres during stress tuning.
Since there is an energy gradient for the polaritons, one expects that they will undergo drift. Figure 5 shows spatially resolved photoluminescence when the laser is tightly focused and moved to one side of the potential minimum in the lower polariton branch. Drift is clearly observed over a distance of more than 100 $`\mu `$m, similar to the drift seen earlierdrift for polaritons in an energy gradient created by a wedge of the cavity thickness. The increase of the PL intensity seen in Figure 4 may be partly related to this effect, since polaritons will concentrate at the bottom of the well instead of diffusing away from the excitation spot.
This method of trapping opens a wide variety of possibilities and promise in the area of microcavity research and BEC of polaritons. Typically, only a tiny region of a wafer is in the strong coupling regime, due to the wedge of the layer thicknesses in standard growth processes. By using stress, one is no longer limited to this small region; the method allows the freedom to use nearly any part of the wafer and tune the bands to the region of strong coupling. Using electric field to tune the resonanceefield has the drawback that the oscillator strength of the exciton changes strongly with electric field. Also, as discussed above, a harmonic potential minimum is essential for Bose-Einstein condensation of polaritons or any other particles in two dimensions. The point of high stress becomes a confining point for carriers, which can be used in a polariton laser. In previous experiments,yama the carriers were in free expansion with diffusion, with energy shifts which depended on the local density.little The present experiments allow theory to treat a quasiequilibrium gas with a known confining potential.
Acknowledgements. We wish to thank V. Hartwell, Z. Vörös, and A. Heberle for the invaluable comments and discussions, and H. Deng, G. Weihs, and Y. Yamamoto for contributions to the design of this sample. This material has been supported by the National Science Foundation under Grant No. 0404912 and by DARPA under Army Research Office Contract No. W911NF-04-1-0075. |
warning/0507/math-ph0507042.html | ar5iv | text | # Polygamma theory, the Li/Keiper constants, and validity of the Riemann Hypothesis
(March 6, 2005)
## Abstract
The Riemann hypothesis is equivalent to the Li criterion governing a sequence of real constants $`\{\lambda _k\}_{k=1}^{\mathrm{}}`$, that are certain logarithmic derivatives of the Riemann xi function evaluated at unity. We investigate a related set of constants $`c_n`$, $`n=1,2,\mathrm{}`$, showing in detail that the leading behaviour $`(1/2)\mathrm{ln}n`$ of $`\lambda _n/n`$ is absent in $`c_n`$. Additional results are presented, including a novel explicit representation of $`c_n`$ in terms of the Stieltjes constants $`\gamma _j`$. We conjecture as to the large-$`n`$ behaviour of $`c_n`$. Should this conjecture hold, validity of the Riemann hypothesis would follow.
Key words and phrases
Li/Keiper constants, Riemann zeta function, Riemann xi function, logarithmic derivatives, Riemann hypothesis, Li criterion, Laurent expansion, Stieltjes constants
Introduction
The Riemann hypothesis is equivalent to the Li criterion governing the sequence of real constants $`\{\lambda _k\}_{k=1}^{\mathrm{}}`$, that are certain logarithmic derivatives of the Riemann xi function evaluated at unity. This equivalence results from a necessary and sufficient condition that the logarithmic derivative of the function $`\xi [1/(1z)]`$ be analytic in the unit disk, where $`\xi `$ is the Riemann xi function. The Li equivalence states that a necessary and sufficient condition for the nontrivial zeros of the Riemann zeta function to lie on the critical line Re $`s=1/2`$ is that $`\{\lambda _k\}_{k=1}^{\mathrm{}}`$ is nonnegative for every integer $`k`$.
This paper is a further contribution to our research program to characterize the Li (Keiper ) constants . We have previously rederived an arithmetic formula for these constants, and described how it could be used to estimate them. Elsewhere, among several other results, we have examined summatory properties of the Li and Stieltjes constants, and investigated the $`\eta _j`$ coefficients appearing in the logarithmic derivative of the zeta function about $`s=1`$ . In particular, a key feature of the sequence $`\{\eta _j\}_{j=0}^{\mathrm{}}`$ is now known: it possesses strict sign alternation .
In this paper, we investigate a related set of constants $`c_n`$, $`n=1,2,\mathrm{}`$, that might be thought of as reduced Li/Keiper constants. We show in detail that the leading behaviour $`(1/2)\mathrm{ln}n`$ of $`\lambda _n/n`$ is absent in $`c_n`$. What remains in $`c_n`$ is a direct manifestation of fundamental properties of the zeta function. Thus, $`c_n`$ can be variously interpreted as reflecting the nontrivial zeros, or the $`\eta _j`$ constants. We present additional analytic results, including an explicit representation of $`c_n`$ in terms of the Stieltjes constants $`\gamma _k`$. We conjecture on the precise order of $`c_n`$ in $`n`$, we comment on the nature of the logarithmic derivative of the zeta function, and we briefly discuss possible interpretations of some of our results.
The Gamma function is important in the theory of the Riemann zeta function– for instance it is needed to complete $`\zeta `$ to the $`\xi `$ function. The Gamma function figures prominently in the functional equation for the zeta function, thereby largely determining the location of the trivial zeros and other analytic properties. Hence the digamma function appears in the logarithmic derivative of the xi function, and higher derivatives introduce the polygamma functions $`\psi ^{(j)}`$ . Portions of the theory of this family of functions are very important in much of this paper. We have also made extensive use of the properties of $`\psi ^{(j)}`$ in previous works .
Our approach is explicit and very much in the spirit of constructivistic mathematics. Indeed, our work may be much more explicit than what would have been thought possible just a few years ago.
From improved numerical calculation to height $`T2.38\times 10^{12}`$ , it is now known that at least the first ten trillion complex zeros of the zeta function lie on the critical line. For our purposes, this effectively ensures that approximately the first $`10^{26}`$ $`\lambda _k`$’s are nonnegative, and this fact may have significant implications for our investigations. For instance, it may very well turn out that working asymptotically in $`k`$ will suffice in our research program.
The Li equivalence is by itself a qualitative reformulation of the Riemann hypothesis. The Riemann hypothesis does not of itself dictate the exact nature of the Li/Keiper constants. In fact, one can easily formulate conjectures on the nature and order of the Li/Keiper constants that are then stronger than the Riemann hypothesis. These observations indicate that the Riemann hypothesis may be verifiable without knowing the optimal order or other properties of the Li/Keiper constants that would more fully characterize them.
It is possible to use our approach also in pursuit of confirmation of the extended and generalized Riemann hypotheses. The corresponding $`\lambda `$ constants have been defined for Dirichlet and Hecke L-functions and other zeta functions , and the same leading behaviour $`O(j\mathrm{ln}j)`$ has been found . Our attention here is strictly with the classical zeta function.
Preliminary Relations
We first recall some notation, introduce some definitions, and present an important Lemma for subsequent developments. We introduce the function
$$F(z)=\mathrm{ln}\left[\frac{z}{1z}\zeta \left(\frac{1}{1z}\right)\right],$$
$`(1)`$
whose analyticity in the unit disc $`|z|<1`$ in the complex plane is equivalent to the Riemann hypothesis. Then, should the power series
$$F(z)=\underset{n=0}{\overset{\mathrm{}}{}}c_nz^n$$
$`(2)`$
converge for $`|z|<1`$, the Riemann hypothesis follows. In short order, one verifies that the $`c_n`$ are real constants, $`c_0=0`$, and $`c_1=\gamma `$, the Euler constant. Therefore, we may write
$$F(z)=\gamma z+\underset{n=2}{\overset{\mathrm{}}{}}c_nz^n.$$
$`(3)`$
This paper investigates the behaviour of the constants $`c_n`$. The importance of this subject is clear: from any subexponential bound on them the Riemann hypothesis follows . Indeed, in this paper we additionally conjecture the true order of the constants $`|c_n|`$, such that this conjecture is stronger than the Riemann hypothesis itself.
Suppose we knew that $`F(z)`$ is analytic and univalent within the unit disc. Then the function $`F(z)/\gamma `$, satisfying $`F(0)/\gamma =0`$ and $`F^{}(0)/\gamma =1`$, is schlicht, fulfills the conditions for the Bieberbach conjecture to hold, and thus we would have $`|c_n|\gamma n`$ for all $`n=1,2,\mathrm{}`$. This shows the self consistency of the complex analysis involved. In fact, any direct application to the Bieberbach conjecture is thwarted due to the essential singularity in $`F`$, making it highly non-univalent.
Related to a prefactor in Eq. (1), $`k(z)=z/(1z)^2`$ is the Koebe function, and $`k_\alpha (z)=z/(1\alpha z)^2`$ with $`|\alpha |=1`$ are rotations of it. These are the only functions for which equality holds in the conclusion of the Bierberbach conjecture. The history of the Bieberbach conjecture shows that it is easier to obtain results about the logarithmic coefficients of a univalent function rather than for the coefficients of the function itself , and the approach of Smith seems to fit within this framework.
The classical Laurent expansion of the Riemann zeta function about the unique pole at $`s=1`$ introduces the Stieltjes constants $`\gamma _k`$ , with $`\gamma _0=\gamma `$. We have
$$\zeta (s)=\frac{1}{s1}+\underset{n=0}{\overset{\mathrm{}}{}}\frac{(1)^n}{n!}\gamma _n(s1)^n,$$
$`(4)`$
where the Stieltjes constants can be written in the form
$$\gamma _k=\underset{N\mathrm{}}{lim}\left(\underset{m=1}{\overset{N}{}}\frac{1}{m}\mathrm{ln}^km\frac{\mathrm{ln}^{k+1}N}{k+1}\right).$$
$`(5)`$
and several other forms have been given . It is clear that the $`c_n`$’s are multinomials in the Stieltjes constants, and that $`c_n`$ contains terms $`(1)^n\gamma ^n/n`$ and $`(1)^n\gamma _{n1}/(n1)!`$. Indeed, we are able to write much more, giving, for instance, an explicit formula for $`c_n`$ in terms of the $`\gamma _k`$’s. Once again, we may therefore observe that sufficient estimation of the Stieltjes constants would provide verification of the Riemann hypothesis.
We discuss these connections with the Stieltjes constants in a later section. For the moment, we simply point out that Appendix A contains explicit formulae for the first few $`c_n`$’s in terms of them.
The conformal map introducing the Li/Keiper constants and the coefficients $`c_n`$ is a natural one. This map $`z=11/ss=1/(1z)`$ takes the right-half plane Re $`s>1/2`$ to the interior of the unit circle in the complex $`z`$-plane. Just this sort of mapping arises in the theory of finite fields, whose zeta functions have zeros on a circle in the complex plane . We recall that Weil proved the Riemann hypothesis holds for nonsingular curves over a finite field , while Deligne established the validity of Weil’s conjectures for generalized hypersurfaces that may include intersections of hypersurfaces .
We first present the explicit connection between the Li/Keiper constants and the constants $`c_n`$. We have
Lemma 1
$$\frac{\lambda _n}{n}=c_n+\frac{1}{n}\frac{1}{2}\mathrm{ln}\pi +d_n,n1,$$
$`(6)`$
where $`d_n`$ is the coefficient of $`z^n`$ of the function $`\mathrm{ln}\mathrm{\Gamma }[1/2(1z)]`$, and $`\mathrm{\Gamma }`$ is the Gamma function. That is,
$$d_n=\frac{1}{n!}\frac{d^n}{dz^n}\mathrm{ln}\mathrm{\Gamma }\left[\frac{1}{2(1z)}\right]_{z=0}.$$
$`(7)`$
We will present not only $`d_n`$, but the derivatives themselves in this equation. We then estimate $`d_n`$ in $`n`$ and demonstrate that the leading behaviour of $`\lambda _n/n`$ in Eq. (6) exactly cancels it. This is highly supportive of a decrease of $`|c_n|`$ with $`n`$, and of the conclusion that the Riemann hypothesis should hold. Since we have previously conjectured as to the subdominant behaviour of the Li constants , we thereby have an immediate conjecture for $`|c_n|`$.
Before proving the Lemma and going on to expressions for $`d_n`$, we give some brief background on the Li (or Keiper) constants. The function $`\xi `$ is determined from $`\zeta `$ by the relation
$$\xi (s)=\frac{s}{2}(s1)\pi ^{s/2}\mathrm{\Gamma }\left(\frac{s}{2}\right)\zeta (s),$$
$`(8)`$
and satisfies the functional equation $`\xi (s)=\xi (1s)`$. The sequence $`\{\lambda _n\}_{n=1}^{\mathrm{}}`$ is defined by
$$\lambda _n=\frac{1}{(n1)!}\frac{d^n}{ds^n}[s^{n1}\mathrm{ln}\xi (s)]_{s=1}.$$
$`(9)`$
The $`\lambda _j`$’s are connected to sums over the nontrivial zeros of $`\zeta (s)`$ by way of
$$\lambda _n=\underset{\rho }{}\left[1\left(1\frac{1}{\rho }\right)^n\right],$$
$`(10)`$
and
$$\lambda _1=\frac{\mathrm{ln}\pi }{2}+\frac{\gamma }{2}+1\mathrm{ln}2.$$
$`(11)`$
In the representation
$$\frac{\lambda _n}{n}=\frac{1}{n}S_1(n)+\frac{1}{n}S_2(n)\frac{1}{2}(\gamma +\mathrm{ln}\pi +2\mathrm{ln}2),$$
$`(12)`$
the sum
$$S_1(n)\underset{m=2}{\overset{n}{}}(1)^m\left(\genfrac{}{}{0pt}{}{n}{m}\right)(12^m)\zeta (m),n2,$$
$`(13)`$
has been characterized :
$$\frac{n}{2}\mathrm{ln}n+(\gamma 1)\frac{n}{2}+\frac{1}{2}S_1(n)\frac{n}{2}\mathrm{ln}n+(\gamma +1)\frac{n}{2}\frac{1}{2}.$$
$`(14)`$
Further bounds on $`S_1(n)`$ have been developed by applying Euler-Maclaurin summation to all orders .
For the sum
$$S_2(n)\underset{m=1}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{m}\right)\eta _{m1},$$
$`(15)`$
the constants $`\eta _j`$ can be written as
$$\eta _k=\frac{(1)^k}{k!}\underset{N\mathrm{}}{lim}\left(\underset{m=1}{\overset{N}{}}\frac{1}{m}\mathrm{\Lambda }(m)\mathrm{ln}^km\frac{\mathrm{ln}^{k+1}N}{k+1}\right),$$
$`(16)`$
and $`\mathrm{\Lambda }`$ is the von Mangoldt function , such that $`\mathrm{\Lambda }(k)=\mathrm{ln}p`$ when $`k`$ is a power of a prime $`p`$ and $`\mathrm{\Lambda }(k)=0`$ otherwise. The constants $`\eta _j`$ enter the expansion around $`s=1`$ of the logarithmic derivative of the zeta function,
$$\frac{\zeta ^{}(s)}{\zeta (s)}=\frac{1}{s1}\underset{p=0}{\overset{\mathrm{}}{}}\eta _p(s1)^p,|s1|<3,$$
$`(17)`$
and the corresponding Dirichlet series valid for Re $`s>1`$ is
$$\frac{\zeta ^{}(s)}{\zeta (s)}=\underset{n=1}{\overset{\mathrm{}}{}}\frac{\mathrm{\Lambda }(n)}{n^s}.$$
$`(18)`$
The constants $`\eta _j`$, with $`\eta _0=\gamma `$, have been written explicitly in terms of the Stieltjes constants , a point on which we return later. Additionally, we recently proved the strict sign alternation of the sequence $`\{\eta _j\}_{j=0}^{\mathrm{}}`$ .
Now that the sum $`S_2(n)`$ has been introduced, we may present the exact relation
Lemma 2
$$\frac{S_2(n)}{n}=c_n.$$
$`(19)`$
For, from Eq. (2) it follows that $`c_n=(1/n!)(d^n/dz^n)F(z)|_{z=0}`$, and from Eqs. (8) and (9) we have
$$S_2(n)=\frac{1}{(n1)!}\frac{d^n}{ds^n}[s^{n1}\mathrm{ln}[(s1)\zeta (s)]]_{s=1}$$
$$=\underset{m=1}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{m}\right)\frac{1}{(m1)!}\frac{d^m}{ds^m}\mathrm{ln}\left[(s1)\zeta (s)\right]_{s=1}.$$
$`(20)`$
Under the mapping $`s(z)=1/(1z)`$, derivatives transform as $`d/dz=s^2d/ds`$, and the Lemma follows.
Proof of Lemma 1 and expressions for $`d_n`$
From Eqs. (1) and (8) we have
$$\mathrm{ln}\xi \left(\frac{1}{1z}\right)=F(z)\mathrm{ln}(1z)+\frac{1}{2(z1)}\mathrm{ln}\pi +\mathrm{ln}\mathrm{\Gamma }\left[\frac{1}{2(1z)}\right].$$
$`(21)`$
Since we have
$$\mathrm{ln}\xi \left(\frac{1}{1z}\right)=\underset{n=1}{\overset{\mathrm{}}{}}\frac{\lambda _n}{n}z^n,$$
$`(22)`$
the expansion of Eq. (21) in powers of $`z`$ readily yields Eq. (6). (For the latter equation, we have used the convention $`\xi _{Li}(z)=2\xi (z)`$, such that $`\xi _{Li}(0)=1`$, in place of $`\xi (z)`$. Otherwise, Eq. (6) will have another minor term $`\mathrm{ln}2`$.)
From Lemmas 1 and 2 it follows that
Corollary
$$\frac{S_1}{n}\frac{\gamma }{2}=\frac{1}{n}+d_n.$$
$`(23)`$
We next have
Lemma 3
$$\frac{d^n}{dz^n}\mathrm{ln}\mathrm{\Gamma }\left[\frac{1}{2(1z)}\right]=\frac{1}{2}\frac{d^{n1}}{dz^{n1}}\frac{1}{(1z)^2}\psi \left[\frac{1}{2(1z)}\right]$$
$$=\frac{1}{2}\underset{j=0}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{n1}{j}\right)\frac{(nj)!}{(1z)^{nj+1}}\frac{d^j}{dz^j}\psi \left[\frac{1}{2(1z)}\right]$$
$$=\frac{1}{2}\left\{\frac{n!}{(1z)^{n+1}}\psi \left[\frac{1}{2(1z)}\right]+\underset{j=1}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{n1}{j}\right)\frac{(nj)!}{(1z)^{n+1}}\underset{\mathrm{}=1}{\overset{j}{}}\left(\genfrac{}{}{0pt}{}{j}{\mathrm{}}\right)\frac{(j1)!}{(\mathrm{}1)!}\frac{1}{2^{\mathrm{}}(1z)^{\mathrm{}}}\psi ^{(\mathrm{})}\left[\frac{1}{2(1z)}\right]\right\},$$
$`(24)`$
where $`\psi `$ is the digamma function and $`\psi ^{(j)}`$ is the polygamma function.
We mention two proofs of this Lemma. The key ingredient is knowing how to write the successive derivatives of the digamma factor. This can be accomplished by applying the Faa di Bruno formula for the $`n`$th derivative of a composite function. We have
$$\frac{d^n}{dz^n}\psi \left[\frac{1}{2(1z)}\right]=\frac{1}{(1z)^n}\underset{j=1}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{j}\right)\frac{(n1)!}{(j1)!}\frac{1}{2^j(1z)^j}\psi ^{(j)}\left[\frac{1}{2(1z)}\right].$$
$`(25)`$
This equation is just a slight extension of a formula for the derivatives of a function $`\theta (1/x)`$ . Equation (25), when used with the product rule, completes the Lemma.
Another method can be based upon the expansion
$$\mathrm{\Gamma }(z)=\frac{1}{z}\mathrm{exp}\left[\gamma z+\underset{k=2}{\overset{\mathrm{}}{}}\frac{(1)^k\zeta (k)}{k}z^k\right],$$
$`(26)`$
giving
$$\mathrm{ln}\mathrm{\Gamma }\left[\frac{1}{2(1z)}\right]=\mathrm{ln}2\underset{n=1}{\overset{\mathrm{}}{}}\frac{z^n}{n}\frac{\gamma }{2}\underset{n=0}{\overset{\mathrm{}}{}}z^n+\underset{k=2}{\overset{\mathrm{}}{}}\frac{(1)^k\zeta (k)}{k2^k}\left(\underset{j=0}{\overset{\mathrm{}}{}}z^j\right)^k.$$
$`(27)`$
Expanding the powers of power series on the right side returns us to the Lemma. Afterall, we have the relation
$$\psi ^{(j)}\left(\frac{1}{2}\right)=(1)^{j+1}j!(2^{j+1}1)\zeta (j+1),j1,$$
$`(28)`$
and this is very helpful in rewriting the constants $`d_n`$ below.
From Eqs. (25) and (28) we have an exact reformulation of digamma derivatives of interest:
Lemma 4
$$\frac{d^n}{dz^n}\psi \left[\frac{1}{2(1z)}\right]_{z=0}=n!\underset{m=1}{\overset{\mathrm{}}{}}\frac{1}{m^2}\left[2\left(1\frac{1}{m}\right)^{n1}\frac{1}{2}\left(1\frac{1}{2m}\right)^{n1}\right].$$
$`(29)`$
This Lemma follows by setting $`z=0`$ in Eq. (25), inserting Eq. (28), using the Dirichlet series for the zeta function, and reordering the double sum. The use of a derivative relation of a binomial sum completes the work.
We are in position to write compact, yet exact, expressions for the constants $`d_n`$:
Lemma 5
$$d_n=\frac{1}{2}\psi \left(\frac{1}{2}\right)+\frac{1}{2n}\underset{j=1}{\overset{n1}{}}(nj)\underset{m=1}{\overset{\mathrm{}}{}}\frac{1}{m^2}\left[2\left(1\frac{1}{m}\right)^{j1}\frac{1}{2}\left(1\frac{1}{2m}\right)^{j1}\right]$$
$$=\frac{1}{2}\psi \left(\frac{1}{2}\right)+\frac{1}{2n}\underset{m=1}{\overset{\mathrm{}}{}}\left[2\left(1\frac{1}{m}\right)^n2\left(1\frac{1}{2m}\right)^n+\frac{n}{m}\right],n1,$$
$`(30)`$
where $`\psi (1/2)/2=\gamma /2\mathrm{ln}2`$. The first line of the Lemma follows from the combination of the results of Lemmas 2 and 3, and the second line follows from application of finite geometric series.
High order approximation for the constants $`d_n`$
There are many ways in which to obtain highly accurate approximations to $`d_n`$ for large values of $`n`$. The upshot is
Lemma 6
$$\mathrm{ln}\mathrm{\Gamma }\left[\frac{1}{2(1z)}\right]\frac{1}{2}\mathrm{ln}\pi +\frac{1}{2}\underset{j=1}{\overset{\mathrm{}}{}}\left[\psi (j)+\gamma \mathrm{ln}21\right]z^j.$$
$`(31)`$
That is, for $`j>>1`$ we have
$$d_j=\frac{1}{2}\left[\mathrm{ln}j\frac{1}{2j}\frac{1}{12j^2}+\gamma \mathrm{ln}21+O\left(\frac{1}{j^4}\right)\right].$$
$`(32)`$
In connection with Eq. (31), we recall the value of the digamma function at integer argument in terms of harmonic numbers $`H_n`$: $`\psi (n)=H_{n1}\gamma `$.
We indicate a couple of approaches for obtaining Lemma 6. One is based upon using the integral corresponding to the summation on the second line of Eq. (30). It turns out that this integral, $`I_1(n)`$, was extensively studied in Appendix A of Ref. , and we have taken over the results.
In another method, we apply Euler-Maclaurin summation to the sum over $`m`$ on the first line of Eq. (30). In doing so, we put
$$f(m)\frac{1}{m^2}\left[2\left(1\frac{1}{m}\right)^{n1}\frac{1}{2}\left(1\frac{1}{2m}\right)^{n1}\right],$$
$`(33)`$
such that $`f(1)=2^n`$, $`f(\mathrm{})=0`$, and we have the elementary integral
$$_1^{\mathrm{}}f(k)𝑑k=\frac{1+2^n}{n}.$$
$`(34)`$
Therefore we obtain
$$\frac{d^n}{dz^n}\psi \left[\frac{1}{2(1z)}\right]_{z=0}n!\left(\frac{1+2^n}{n}2^{n1}\right).$$
$`(35)`$
Then the sum over $`j`$ can be performed in Eq. (30). In either approach, we discard any terms in the final result that are exponentially small in $`n`$, such as $`2^{1n}/n`$.
Of note, all terms in $`d_n`$ beyond the leading logarithmic dependence are in terms of integral powers of $`1/n`$–there is no algebraic dependence upon $`n`$.
Relations and formulae for $`c_n`$
With the aid of Cauchy’s integral formula, the $`c_n`$’s can be written as
$$c_n=\frac{1}{2\pi i}_C\frac{F(z)}{z^{n+1}}𝑑z,$$
$`(36)`$
where $`C`$ is a simple closed contour about the origin, and this can serve as the basis of a numerical method . If $`C`$ is a circle of radius $`r`$ about the origin, then $`c_n=r^n_0^1F(re^{2\pi i\varphi })e^{2\pi in\varphi }𝑑\varphi `$, and this invites the use of fast Fourier transform for evaluation.
The combination of the result of Lemma 5 with Eqs. (6), (12), and (14) shows that $`c_n=S_2(n)/n+O(1/n)`$, and Lemma 2 gives the strengthening to $`c_n=S_2(n)/n`$. Since we have previously conjectured that $`|S_2(n)|=O(n^{1/2+\epsilon })`$ for $`\epsilon >0`$, we have
Conjecture
$$|c_n|=O\left(\frac{1}{n^{1/2\epsilon }}\right),$$
$`(37)`$
where $`\epsilon >0`$ but is otherwise arbitrary. That is, we anticipate that the magnitudes $`|c_n|`$ decrease nearly as the square root of $`n`$ for large $`n`$. In Figure 1 we compare such a decrease with available numerical evidence . Figure 1 contains a semilogarithmic plot of $`|c_n|`$ versus $`n`$, together with a curve corresponding to $`6/\pi ^2\sqrt{n}`$. For this limited set, after a few initial values, the latter curve appears to provide a consistent upper bound. In light of the von Koch result on the Riemann hypothesis that $`\psi (x)=x+O(x^{1/2}\mathrm{ln}^2x)`$ , where $`\psi `$ is the Chebyshev function, we suspect that the optimal order of $`|c_n|`$ is very close to $`O(\mathrm{ln}n/n^{1/2})`$. Figure 2 shows an example plot of $`c_n`$ versus $`n`$, that illustrates the oscillatory behaviour of these constants. In addition, Smith has now numerically confirmed our conjecture for the first approximately $`10^5`$ values of $`c_n`$ .
In Figure 3 we have plotted the differences $`\delta _n=c_n^2c_{n1}c_{n+1}`$ versus $`n`$, that appear to support a decrease in $`|c_n|`$ with increasing $`n`$. In addition, the behaviour may indicate a correlation in the sign or other properties of the $`c_n`$’s. Figure 4 plots the magnitude of the discrete Fourier transform applied to this sequence. This plot indicates underlying structure.
Our conjecture suggests that it may be worthwhile to study in detail the properties of the particular polylogarithm
$$L_{1/2}(z)\underset{n=1}{\overset{\mathrm{}}{}}\frac{z^n}{n^{1/2}},$$
$`(38)`$
such that $`L_{1/2}(1)=(\sqrt{2}1)\zeta (1/2)`$.
Previously, we obtained an expression for the Li/Keiper constants explicitly in terms of the Stieltjes constants . We recall this and related results :
Theorem
$$\lambda _n=1\frac{n}{2}(\mathrm{ln}\pi +2\mathrm{ln}2\gamma )+S_1(n)\underset{j=2}{\overset{n}{}}(1)^j\left(\genfrac{}{}{0pt}{}{n}{j}\right)j\underset{h=1}{\overset{j}{}}\frac{1}{h}\underset{\stackrel{j_10,\mathrm{},j_h0}{j_1+\mathrm{}+j_h=jh}}{}\underset{b=1}{\overset{h}{}}\frac{\gamma _{j_b}}{j_b!},n2,$$
$`(39)`$
$$S_2(n)=n\gamma \underset{j=2}{\overset{n}{}}(1)^j\left(\genfrac{}{}{0pt}{}{n}{j}\right)j\underset{h=1}{\overset{j}{}}\frac{1}{h}\underset{\stackrel{j_10,\mathrm{},j_h0}{j_1+\mathrm{}+j_h=jh}}{}\underset{b=1}{\overset{h}{}}\frac{\gamma _{j_b}}{j_b!},n2,$$
$`(40)`$
and
Theorem
$$\eta _{k1}=(1)^kk\underset{h=1}{\overset{k}{}}\frac{1}{h}\underset{\stackrel{j_10,\mathrm{},j_h0}{j_1+\mathrm{}+j_h=kh}}{}\underset{b=1}{\overset{h}{}}\frac{\gamma _{j_b}}{j_b!},k2.$$
$`(41)`$
From Lemma 2 we obtain the exact relation
Theorem
$$c_n=\gamma \frac{1}{n}\underset{j=2}{\overset{n}{}}(1)^j\left(\genfrac{}{}{0pt}{}{n}{j}\right)j\underset{h=1}{\overset{j}{}}\frac{1}{h}\underset{\stackrel{j_10,\mathrm{},j_h0}{j_1+\mathrm{}+j_h=jh}}{}\underset{b=1}{\overset{h}{}}\frac{\gamma _{j_b}}{j_b!},n2.$$
$`(42)`$
On the right side of Eq. (42), the constrained sum over the indices $`j_{\mathrm{}}`$ means that we have a partition of $`kh`$ over the nonnegative integers. All such partitions are considered, meaning that their order does not matter. The number of such partitions is $`\left(\genfrac{}{}{0pt}{}{n1}{h1}\right)`$ in $`\eta _{n1}`$ or $`c_n`$.
From Eq. (4) we have
$$\frac{z}{1z}\zeta \left(\frac{1}{1z}\right)=1+\underset{n=0}{\overset{\mathrm{}}{}}\frac{(1)^n}{n!}\gamma _n\left(\frac{z}{1z}\right)^{n+1}.$$
$`(43)`$
Then from the definition (1) and performing various expansions, we have
$$F(z)=\underset{n=1}{\overset{\mathrm{}}{}}\frac{(1)^n}{n}\left[\underset{j=0}{\overset{\mathrm{}}{}}\frac{(1)^j}{j!}\gamma _jz^{j+1}\left(\underset{m=0}{\overset{\mathrm{}}{}}z^m\right)^{j+1}\right]^n.$$
$`(44)`$
Carrying out the expansion in powers of $`z`$ in this equation must necessarily return us to Eq. (42) for the coefficients $`c_n`$.
From the Hadamard product formula for the zeta function (e.g., ),
$$\zeta (s)=\frac{\mathrm{exp}(\mathrm{ln}2\pi 1\gamma /2)s}{2(s1)\mathrm{\Gamma }(s/2+1)}\underset{\rho }{}\left(1\frac{s}{\rho }\right)e^{s/\rho },$$
$`(45)`$
we obtain
$$F(z)=\frac{\mathrm{ln}2\pi 1\gamma /2}{1z}\mathrm{ln}2\mathrm{ln}\mathrm{\Gamma }\left[\frac{32z}{2(1z)}\right]+\underset{\rho }{}\left\{\mathrm{ln}\left[1\frac{1}{\rho (1z)}\right]+\frac{1}{\rho (1z)}\right\}.$$
$`(46)`$
Since by Eq. (10), $`_\rho 1/\rho \sigma _1=\lambda _1`$, we have
$$F(z)=\frac{\mathrm{ln}\pi }{2(1z)}\mathrm{ln}2\mathrm{ln}\mathrm{\Gamma }\left[\frac{32z}{2(1z)}\right]+\underset{\rho }{}\mathrm{ln}\left[1\frac{1}{\rho (1z)}\right].$$
$`(47)`$
The value $`F(0)=_\rho \mathrm{ln}[(\rho 1)/\rho ]=0`$ obtains because the sum over all the complex zeros of $`\zeta `$ contains the pairs of $`\rho `$ with $`1\rho `$.
From the functional equation in the form $`\zeta (z)=\pi ^{z1}2^z\mathrm{\Gamma }(1z)\zeta (1z)\mathrm{sin}(\pi z/2)`$, we obtain
$$F\left(\frac{1}{z}\right)=F(z)\mathrm{ln}z+\frac{\mathrm{ln}(\pi )}{z1}+\frac{z}{z1}\mathrm{ln}2+\mathrm{ln}\mathrm{\Gamma }\left(\frac{1}{1z}\right)+\mathrm{ln}\mathrm{sin}\left[\frac{\pi }{2}\frac{z}{(z1)}\right].$$
$`(48)`$
This equation should be very useful in obtaining results on the boundedness of the $`c_n`$’s.
Discussion of the logarithmic derivative of $`\zeta `$
This function has proved to be central in analytic number theory. Here we recall some known results and relate them to Eqs. (17), (18), and others.
We have (e.g., ) in terms of the prime counting function $`\pi (x)`$
$$\mathrm{ln}\zeta (s)=s_2^{\mathrm{}}\frac{\pi (x)}{x(x^s1)}𝑑x,\text{Re}s>1.$$
$`(49)`$
Then
$$\frac{\zeta ^{}(s)}{\zeta (s)}=_2^{\mathrm{}}\frac{\pi (x)}{x(x^s1)}𝑑xs_2^{\mathrm{}}\frac{x^{s1}\pi (x)\mathrm{ln}x}{(x^s1)^2}𝑑x,\text{Re}s>1.$$
$`(50)`$
Since the function $`\pi `$ has steps, these induce changes in the coefficients $`\eta _j`$, hence in $`S_2(n)`$ or the $`c_n`$ constants. In the common region of validity Re $`s>1`$ $``$ $`|s1|<3`$, we have from Eqs. (17) and (49)
$$[(s1)+1]_2^{\mathrm{}}\frac{\pi (x)}{x(x^s1)}𝑑x=\mathrm{ln}(s1)\underset{p=1}{\overset{\mathrm{}}{}}\frac{\eta _{p1}}{p}(s1)^p.$$
$`(51)`$
That is, we could write for instance
$$\mathrm{ln}(s1)\underset{p=1}{\overset{\mathrm{}}{}}\frac{\eta _{p1}}{p}(s1)^p=[(s1)+1]_2^{\mathrm{}}dx\frac{\pi (x)}{x}\{\frac{1}{x1}\frac{x\mathrm{ln}x}{(x1)^2}(s1)$$
$$+[\frac{x}{2}\frac{\mathrm{ln}^2x}{(x1)^2}+\frac{x^2\mathrm{ln}^2x}{(x1)^3}](s1)^2+O[(s1)^3]\},$$
$`(52)`$
where of course by the prime number theorem $`\pi (x)x/\mathrm{ln}x`$ as $`x\mathrm{}`$. This equation in powers of $`s1`$ gives in principle an integral representation for each of the coefficients $`\eta _j`$.
We mention an important occurrence of the logarithmic derivative in numerical analysis. The reciprocal of this function is key in the classical Newton iteration for root finding, bringing in connections with discrete dynamical systems. Then one seeks the attracting fixed points of the associated Newtonian mapping. Therefore, from this point of view it is not unexpected that the logarithmic derivative should play an important role in determining zeros.
Summary and Brief Discussion
By way of Lemma 2, or equivalently from the use of the theory of polygamma functions, we have shown that the constants $`c_n`$ of Eqs. (2) and (3) omit the leading growth $`(1/2)\mathrm{ln}n`$ of $`\lambda _n/n`$. We have presented additional analytic arguments, a conjecture, and partial numerical results that point to the decrease of $`|c_n|`$ with $`n`$. We have pointed out the limited possibility of directly applying the Bieberbach conjecture because the function $`F`$ of Eq. (1) is not univalent within the unit disc.
The quantity $`S_2(n)`$ is formed as the binomial sum of the alternating $`\eta _j`$ values of Eq. (17). The latter is a correlated sequence. For, we have previously exhibited the explicit summatory relation imposed upon the $`\eta `$’s by the functional equation of either the zeta or xi functions. This relation implies that a given $`\eta _j`$ is connected to all the other values $`\eta _{j+1}`$, $`\eta _{j+2}`$, $`\mathrm{}`$.
In Appendix B we call out an integral representation of the alternating zeta function that may permit a joining of probabilisitic interpretation of the zeta function with Krein spectral shift functions. In turn, this may provide a useful link between Hardy space theory and inverse scattering theory. Though fairly independent of the approach of this paper, we believe it may be worth pointing this out to other investigators.
The importance of an explicit formula for $`S_2(n)`$ or $`c_n`$ should not be overlooked. For instance, in principle, only improved estimation of the Stieltjes constants prevents verification of the Riemann hypothesis by way of either the Li criterion or by way of Criterion (c) of Ref. . Concerning the magnitudes $`|c_n|`$, any subexponential bound would serve to verify the Riemann hypothesis.
Acknowledgements
I thank W. van Dam for useful discussion. I thank K. Maślanka for useful correspondence and the numerical values of $`S_2(n)`$ used for the figures. I thank W. D. Smith for useful correspondence. This work was partially supported by a SPARC grant from Regis University.
Figure Captions
FIG. 1. In this semilogarithmic plot, the upper curve corresponds to values of $`6/\pi ^2\sqrt{n}`$ versus $`n`$, and the lower to values of $`|c_n|`$ versus $`n`$.
FIG. 2. Plot of $`c_n`$ versus $`n`$.
FIG. 3. Plot of the differences $`\delta _n=c_n^2c_{n1}c_{n+1}`$ versus $`n`$.
FIG. 4. Plot of the magnitude of the discrete Fourier transform of the $`c_n`$ sequence.
Appendix A: Examples of $`c_n`$ in terms of the Stieltjes constants
Here, $`\gamma `$ is the Euler constant and $`\gamma _k`$ are the Stieltjes constants appearing in Eq. (4). We have
$$c_1=\gamma ,c_2=\gamma \gamma ^2/2\gamma _1,$$
$`(A.1)`$
$$c_3=\gamma \gamma ^2+\frac{1}{3}\gamma ^32\gamma _1+\gamma \gamma _1+\frac{1}{2}\gamma _2,$$
$`(A.2)`$
$$c_4=\gamma ^3\frac{1}{4}\gamma ^4\frac{1}{2}\gamma ^2(3+2\gamma _1)+\gamma (1+3\gamma _1\frac{1}{2}\gamma _2)+\frac{1}{6}[3\gamma _1(6+\gamma _1)+9\gamma _2\gamma _3],$$
$`(A.3)`$
and
$$c_5=\gamma ^4+\frac{1}{5}\gamma ^5+\gamma ^3(2+\gamma _1)+\frac{1}{2}\gamma ^2(48\gamma _1+\gamma _2)+\gamma [1+\gamma _1(6+\gamma _1)2\gamma _2+\frac{1}{6}\gamma _3]$$
$$+\frac{1}{24}[72\gamma _2+12\gamma _1(84\gamma _1+\gamma _2)16\gamma _3+\gamma _4].$$
$`(A.4)`$
The first few $`d_k`$’s are given by
$$d_0=\frac{1}{2}\mathrm{ln}\pi ,$$
$`(A.5)`$
$$d_1=\gamma /2\mathrm{ln}2,$$
$`(A.6)`$
$$d_2=\gamma +\frac{1}{8}\pi ^22\mathrm{ln}2,$$
$`(A.7)`$
$$d_3=3\gamma +\frac{3}{4}\pi ^26\mathrm{ln}2\frac{7}{4}\zeta (3),$$
$`(A.8)`$
$$d_4=12\gamma +\frac{9}{2}\pi ^2+\frac{\pi ^4}{16}24\mathrm{ln}221\zeta (3),$$
$`(A.9)`$
and
$$d_5=60\gamma +30\pi ^2+\frac{5}{4}\pi ^4120\mathrm{ln}2210\zeta (3)\frac{93}{4}\zeta (5).$$
$`(A.10)`$
Appendix B: The alternating zeta function in inverse spectral theory
We first recall the alternating zeta function
$$\underset{n=1}{\overset{\mathrm{}}{}}\frac{(1)^{n1}}{n^s}=(12^{1s})\zeta (s),\text{Re}s>0,s1,$$
$`(B.1)`$
this being one of the many analytic continuations of the Dirichlet series for the Riemann zeta function. Without going into the details, it turns out that this function can be written as an integral representation with the Krein spectral shift function associated with the harmonic oscillator Hamiltonian on the line, with a Dirichlet boundary condition at the origin. As a Corollary, one may write
$$(12^{1s})\zeta (s)=s_0^{\mathrm{}}e^{sx}\varphi (x)𝑑x,\text{Re}s>0,$$
$`(B.2)`$
where
$$\varphi (x)=\underset{n=1}{\overset{\mathrm{}}{}}\chi _{[\mathrm{ln}(2n1),\mathrm{ln}2n]}(x),$$
$`(B.3)`$
and $`\chi `$ is the characteristic function of an interval. More generally, if the Dirichlet boundary condition is enforced at any other point $`x`$, a family of functions $`\zeta (x,s)`$ is generated. The points of discontinuity of $`\zeta (x,s)`$ satisfy a differential equation in $`x`$ called the Dubrovin equation. This differential equation gives a curve in the space of analytic functions with the alternating zeta function divided by $`s`$ as the initial value .
Now it is possible to represent the function
$$\eta (s)\frac{(s1)}{s^2}\zeta (s)=_0^{\mathrm{}}e^{xs}\varphi _1(x)𝑑x,\text{Re}s>0,$$
$`(B.4)`$
with the real-valued function
$$\varphi _1(x)=\underset{1ne^x}{}(1+\mathrm{ln}nx).$$
$`(B.5)`$
Then it is possible to introduce a family of probability densities with $`x[0,\mathrm{})`$ as $`p_\sigma (x)=\varphi _1(x)\mathrm{exp}(\sigma x)/\eta (\sigma )`$ for $`\sigma >0`$ . The cumulants of $`p_\sigma `$ can be written either in terms of the Stieltjes constants or the Li/Keiper constants at $`\sigma =1`$ . Comparing Eqs. (B.2) with (B.4) and (B.3) with (B.5), it appears that it should be possible to combine a probabilistic setting for the zeta function with an inverse spectral theory. This point of view offers a connection between quantum dynamics and stochastic processes. |
warning/0507/cond-mat0507558.html | ar5iv | text | # Multiscale Modeling of Materials - Concepts and Illustration
## I Introduction
The field of multiscale modeling has opened a new era to computational science for studying complex phenomena like fracture, hydrolysis, enzymatic reactions, solute-solvent studies, hydrogen embrittlement, and many other chemo-mechanical processes in macroscopic samples. Often these phenomena require a very accurate description at one scale, while at another scale a much coarser method can be applied to get satisfactory results. In fact, the coarser description for the bulk sample is required because the more fundamental methods are computationally too intensive to be applied to the entire system. These scales may be length or time, or a combination of both. This scheme of combining different models at different scales to achieve a balance of accuracy, efficiency, and realistic description is known as multiscale modeling. It is accomplished by applying the high accuracy method only in a small domain where it is needed the most and more approximate methods for the rest of the bulk where they are appropriate. For example, in crack propagation one has to apply detailed quantum mechanics at the tip of the crack. Here the bonds are breaking between the atoms, causing a marked deformation of the electron cloud and charge transfer between ions. But far from the crack tip where the atoms are less deformed they can be described by classical mechanics with appropriate potentials. Thus, the problem arises of linking quite different descriptions for the different length scales Broughton1999 .
There is a broad diversity of other examples requiring multiscale modeling: chemo-mechanical polishing Singh2002 , tidal wave prediction Clementi1988 , atmospheric sciences, embrittlement of nuclear reactors Odette2001 , and many biological systems Gogonea2002 . Two classes can be distinguished Rudd2000 , “serial multiscale modeling” where the various scales are weakly coupled but the computation of parameters at smaller scales is required for its use in more phenomenological models at a larger scale, and “concurrent multiscale modeling” where the various descriptions applied on different scales should all be nested with proper boundary conditions. The multiscale modeling discussed here is of the “concurrent” type, although there are components that can be considered as the “serial” type (e.g., parameterization of the underlying approximate quantum method used).
Three different length scale levels are distinguished, the nanoscale where the details of electron structure and quantum chemistry are important, the atomic or microscale where appropriate classical pair potentials capture the structure accurately, and the macroscale where a continuum mechanical description applies. The present work addresses only one particular subclass of multiscale modeling, known as hybrid quantum mechanical (QM)/classical mechanical (CM) problems, and their application to solid insulators - silica in particular. The further embedding of the atomic scale in a macroscale model is not considered here. Early work on this QM/CM topic began with Warshel and Levitt Warshel1976 and has accelerated over the past decadeGao1996 . Applications of QM/CM include biological systems where this scheme is sometimes referred to as QM/MM (molecular mechanics) modeling (enzymes, DNAs and proteins) Monard1999 ; Amara1999 ; Cui2002 ; Titmuss2000 ; Murphy2002 , vibrational spectroscopy Chabana2001 , electronic excitations Thompson1996 ; Gao1996 ; Vries1996 , hydrolysis of silica Jung2001 ; Du2003 , and solute-solvent problems Eichinger1999 ; Guo1992 ; Gao1997 .
The challenge of the QM/CM simulation is to move from one length scale to another, as smoothly as possible. A solid or large cluster behaves as a single ”molecule” so the partitioning becomes more complicated due to covalent bonds that are cut between the QM and CM regions. These dangling bonds give rise to incorrect charge densities in the QM domain and other pathologies Sauer2000 ; Ogata2004 . A multiscale modeling scheme is proposed here that addresses this and other problems of the quantum embedding to provide a means for faithful coupling between the QM and CM regions. There are many tests for fidelity across the QM/CM interface, such as bond angles, bond lengths, and proton affinities. Instead, the criteria here are preservation of accurate electronic charge densities and total forces in the QM sub domain. These are the physically significant properties based more closely in the theoretical structure being modeled, as described in the next section. The approach includes the following two components:
(i) Modeling environmental effects for the quantum domain in an accurate and simple manner. Here, accurate means that the forces and charge densities of the QM domain are reproduced to within a few percent. This entails a chemically correct treatment of the dangling bonds, and an account of longer range Coulomb influences from the classical domain. Simplicity, means that the approximations made are transparent and that the difficulty of the quantum calculation is not increased. Here, the valency problem is solved by pseudo-atoms and the longer range environmental effects are included through low order polarization effects Mallik2004 .
(ii) Developing a new classical potential for the classical region tailored to the properties of interest. A new potential for the entire CM region is designed, based on selected force data calculated from the quantum method used in the QM domain. This data is generated for both equilibrium and near equilibrium states. The test for a suitable potential, in addition to producing the right structure, is the accurate reproduction of the chosen property to be studied, for near equilibrium states (typically some linear response characteristic) Mallik20042 . The use of this new potential prevents any mismatch of that response property across the QM/CM boundary. Hence different problems (fracture or corrosion) might require the construction of different potentials for the same system.
The criterion for the success of (i) and (ii) is that calculations of the desired properties should be indistinguishable for the quantum system, purely classical system, and the composite QM/CM system for near equilibrium states. Only then are the desired applications to far from equilibrium states reliable.
In the next four sections, the theoretical basis for this modeling scheme is described in some detail. The object is to provide a means to assess the approximations made more directly, and also to provide the basis for their extension to other systems. The complete formalism is then tested for a model system - a silica nanorod Zhu2003 (Figure 1) - in section VI. This silica nanorod containing 108 nuclei has been chosen since a quantum treatment of the whole sample is possible to assess quantitatively all aspects of the proposed scheme. Still, the system is big enough to yield bulk properties like stress-strain response (also shown in Figure 1). The nanorod has the proper stoichiometric ratio of silicon to oxygen observed in real silica (1:2) and is considered a viable model for studying silica. One of the rings near the center of the rod is chosen to be the QM domain and rest of the rod is treated classically (see Figure 6 below).
The system is studied under uniaxial strain. The quantum mechanical method used is the Transfer Hamiltonian (TH) method Taylor2003 using an NDDO type theory parameterized by coupled cluster data (described more completely in section VI). The construction of the classical pair potential is described and shown to yield the same Young’s constant as obtained from the QM calculation, to within a few percent. Next, the training of the pseudo atoms and validity of polarization representation for the environment of the QM system is tested. The resulting force and charge densities for the QM domain are accurately reproduced to within a few percent. Finally, the composite QM/CM rod is shown to have the same structure and elastic properties as both the QM and CM rods for near equilibrium states, as required. The entire analysis is repeated using density functional theory (DFT) as the underlying quantum theory, instead of the TH quantum mechanics. The same level of accuracy for the modeling is found to hold in this case as well.
## II The Idealized Quantum Solid
In this section the problem of multiscale modeling is introduced by first stating clearly what is the system being modeled. While much of this introductory matter is familiar, it provides the context for assessing *all* approximations involved in the final model material. First, the quantum description in terms of ions and electrons is given and the limitations for practical implementation are noted. Next, reduced self-consistent descriptions are given for the ions and electrons, each coupled to the other through their average charge density. In this form the classical limit for the ions can be taken, allowing the use of MD simulation methods for their dynamics. The equation for the electrons is simplified in a different direction, exploiting the very different time scales for electron and nuclear motion (the Born-Oppenheimer limit). This final description constitutes the idealized quantum solid for which the subsequent multiscale modeling scheme is proposed.
At the fundamental level a simple atomic or molecular solid can be described in terms of $`N_i`$ ”ions” with charge number $`Z_i`$ for the corresponding atoms of species $`i`$ and a set of $`N_e`$ electrons, with overall charge neutrality ($`_iN_iZ_i=N_e`$). To introduce the various levels of description it is useful to start with the density operator $`D`$ for the system as a whole. Properties of interest $`A`$ are given by the expectation value
$$A=Tr_{e,i}DA,$$
(1)
where the trace is taken over all electron and ion degrees of freedom. The density operator obeys the Liouville - von Neumann equation
$$_tD+\frac{i}{\text{h}\text{ }\text{ }}[H,D]=0.$$
(2)
The Hamiltonian operator $`H`$ is comprised of the Hamiltonians $`H_i`$ and $`H_e`$ for the isolated systems of ions and electrons, respectively, and their interaction $`U_{ie}`$
$$H=H_i+H_e+U_{ie},$$
(3)
$$U_{ie}=𝑑𝐫𝑑𝐫^{}\frac{\rho _i\left(𝐫\right)\rho _e\left(𝐫^{}\right)}{\left|𝐫𝐫^{}\right|}.$$
(4)
Here $`\rho _i\left(𝐫\right)`$ and $`\rho _e\left(𝐫\right)`$ are the ion and electron charge density operators. This is the most fundamental level for a quantum description of the system. For a pure state, $`D`$ is a projection onto a state $`\mathrm{\Psi }`$ which is determined from the corresponding Schroedinger equation. For small systems, this can be solved numerically by quantum Monte Carlo methods. At this level, the approximations involved (e.g., electron nodal location) can be considered mild and under control. However, quantum Monte Carlo methods are restricted to relatively small systems and the above fundamental description is not a practical method for application to larger systems, particularly if repeated calculation is required to follow the dynamics.
In many cases the properties of interest are functions only of the ion degrees of freedom (e.g. structure), $`AA_i`$, or they are purely electronic (e.g., optical), $`AA_e`$. Then a description is possible in terms the reduced density operators for the ions and for the electrons, $`D_i`$ and $`D_e`$ , resulting from appropriate partial traces over all the electrons or ions, respectively
$$A_i=Tr_iD_iA_i,A_e=Tr_eD_eA_e,$$
(5)
$$D_iTr_eD,D_eTr_iD.$$
(6)
Their equations follow directly from (2)
$$_tD_i+\frac{i}{\text{h}\text{ }\text{ }}[H_i,D_i]+\frac{i}{\text{h}\text{ }\text{ }}\left(U_iD_iD_iU_i^{}\right)=0,$$
(7)
$$_tD_e+\frac{i}{\text{h}\text{ }\text{ }}[H_e,D_e]+\frac{i}{\text{h}\text{ }\text{ }}\left(U_eD_eD_eU_e^{}\right)=0.$$
(8)
where the potential energy operators coupling the ion and electronic degrees of freedom are
$$U_i=𝑑𝐫𝑑𝐫^{}\frac{1}{\left|𝐫𝐫^{}\right|}\rho _i\left(𝐫\right)\stackrel{~}{\rho }_e\left(𝐫^{}\right),\stackrel{~}{\rho }_e\left(𝐫\right)=\left(Tr_e\rho _e\left(𝐫\right)D\right)\left(Tr_eD\right)^1,$$
(9)
$$U_e=𝑑𝐫𝑑𝐫^{}\frac{1}{\left|𝐫𝐫^{}\right|}\stackrel{~}{\rho }_i\left(𝐫\right)\rho _e\left(𝐫^{}\right),\stackrel{~}{\rho }_i\left(𝐫\right)=\left(Tr_i\rho _i\left(𝐫\right)D\right)\left(Tr_iD\right)^1$$
(10)
This is similar to the microscopic ion - electron coupling of (4) except that now the electron charge density operator $`\rho _e\left(𝐫\right)`$ is replaced by its conditional average, $`\stackrel{~}{\rho }_e\left(𝐫\right)`$, in the equation for $`D_i`$, and the ion charge density $`\rho _i\left(𝐫\right)`$ is replaced by its conditional average, $`\stackrel{~}{\rho }_i\left(𝐫\right)`$ in the equation for $`D_e`$. The description (7) and (8) is still exact but formal since these average charge densities are not determined by these equations alone. The simplest realistic approximation (mean field) is to neglect the direct correlations in the charge densities $`\stackrel{~}{\rho }_i\left(𝐫\right)`$ and $`\stackrel{~}{\rho }_e\left(𝐫\right)`$, i.e. replace $`DD_eD_i`$ in (7) and (10) to get
$$\stackrel{~}{\rho }_i\left(𝐫\right)Tr_i\rho _i\left(𝐫\right)D_i\overline{\rho }_i\left(𝐫\right),\stackrel{~}{\rho }_e\left(𝐫\right)Tr_e\rho _e\left(𝐫\right)D_e\overline{\rho }_e\left(𝐫\right).$$
(11)
As a consequence, the potentials $`U_i`$ and $`U_e`$ become Hermitian and the average charge densities are now self-consistently determined by (7) and (8). Self-consistency is required since $`D_e`$ and $`D_i`$ are functionals of the average charge densities $`\overline{\rho }_i\left(𝐫\right)`$ and $`\overline{\rho }_e\left(𝐫\right)`$, respectively.
The advantage of this reduced description is that the ions and electrons are described by separate equations, reducing the degree of difficulty of each. More importantly, this separation allows the introduction of appropriate approximations for each. The large differences in electron and ion masses imply corresponding differences in time scales and thermal de Broglie wavelengths. Consequently, the equation for the ions admits a classical limit for the conditions of interest, while that for the electrons does not. The classical limit of equation (7) becomes
$$_tD_i+\{\left(H_i+U_i\left[\overline{\rho }_e\right]\right),D_i\}=0,$$
(12)
where $`\{,\}`$ now denotes a Poisson bracket operation, and $`D_i`$ is a function of the ion positions and momenta, $`D_iD_i(\{𝐑_{i\alpha }\},\{𝐏_{i\alpha }\},t)`$ (here $`\alpha `$ denotes a specific ion). This equation now can be solved accurately and efficiently by molecular dynamics simulation methods, even for large systems. This is an essential step in almost all practical descriptions of bulk materials whose importance cannot be overstated.
Implementation of (12) still requires calculation of the potential energy $`U_i\left[\overline{\rho }_e\right]`$, which in turn requires determination of the electronic charge density from $`D_e`$. However, the general solution to (8) is a formidable problem: determination of the dynamics of $`N_e`$ electrons self-consistently in the presence of a changing ion charge density. Two simplifications are made to bring this problem under control. First, it is recognized that $`\tau \left(_tD_i\right)D_i^1<<1`$, where $`\tau `$ is the time scale for changes in the electron distribution function $`D_e.`$ Second, it is assumed that only the lowest energy state contributes to $`D_e`$ at any given time. The first approximation constitutes the Born-Oppenheimer approximation while the second approximation is appropriate for most structural studies, but must be relaxed for optical studies involving electronic transitions. In principle, (12) is solved in time steps $`\mathrm{\Delta }t`$. At each time step the electronic charge density is computed for the ion configurations at that time step in order to recompute new forces for the next time step. The electronic charge density calculation is a ground state eigenvalue problem for the given ion configuration.
To summarize, the final description of this idealized solid, it consists of a set of point ions governed by the classical equation (12) for their probability distribution $`D_iD_i(\{𝐑_{i\alpha }\},\{𝐏_{i\alpha }\},t)`$
$$_tD_i+\{\left(H_i+U_i\left[\overline{\rho }_e\right]\right),D_i\}=0,$$
(13)
$$U_i=𝑑𝐫𝑑𝐫^{}\frac{1}{\left|𝐫𝐫^{}\right|}\rho _i\left(𝐫\right)\overline{\rho }_e(𝐫^{},t),\overline{\rho }_e(𝐫,t)=Tr_e\rho _e\left(𝐫^{}\right)D_e\left[\overline{\rho }_i\right]$$
(14)
The average electron density $`\overline{\rho }_e\left(𝐫\right)`$ is determined from the ground state solution to (8)
$$\frac{i}{\text{h}\text{ }\text{ }}[H_e+U_e,D_e]=0,$$
(15)
$$U_e=𝑑𝐫𝑑𝐫^{}\frac{1}{\left|𝐫𝐫^{}\right|}\overline{\rho }_i(𝐫,t)\rho _e\left(𝐫^{}\right),\overline{\rho }_i(𝐫,t)=Tr_e\rho _i\left(𝐫\right)D_i\left[\overline{\rho }_e\right]$$
(16)
The analysis proceeds stepwise. The classical equations (13) are solved analytically in discrete time steps for the atomic coordinates and momenta. At each step the electron problem (15) is solved for the electron ground state using the ion configuration at the previous time step. From this ground state the electron charge density $`\rho _e\left(𝐫\right)`$ is determined. This gives the potential energy function $`U_i`$ and consequently the forces required to change the ion positions and momenta at the next time step. The process is repeated with a new electron charge density calculated at each time step using the new ionic configurations. All electron correlations are accounted for quantum mechanically in the eigenvalue problem; all atomic correlations are determined classically through direct solution of Newton’s equations. All structural properties of interest can be calculated since the phase points for the ions are known at all times. The dynamics of the electrons is only coarse-grained as they are ”slaved” to the time dependence of the ions.
The above defines a ”quantum molecular dynamics” representation of the idealized quantum solid. The semi-classical approximation for the ions, and ground state Born-Oppenhiemer approximation for the electrons are relatively weak under most conditions of interest. The resulting description allows an accurate classical treatment of the ionic structure while retaining relevant quantum chemistry for the interatomic forces due to electrons. If it could be implemented in practice for bulk systems of interest over reasonable time intervals there would be little need for multiscale modeling. With MD simulation the solution to (13) once $`U_i`$ has been provided is straightforward. So the problem has been reduced to a determination of the electron charge density. Unfortunately, the solution to (15) using realistic quantum chemical methods for even a few hundred ions at each time step becomes prohibitively time intensive.
## III The Formal Partition and Composite Solid
### III.1 The Representative Classical Solid
In many cases of interest (e.g., equilibrium structure, thermodynamics) the computational limitations of quantum chemical methods can be avoided through a purely classical representation of the solid, avoiding the intensive electron charge density calculation at each time step. This entails an idealization that has many variants. It consists of the representation of the true potential energy function $`U_i\left(\{𝐑_{i\alpha }\}\right)`$ in (13) by a suitably *chosen* function $`U_c\left(\{𝐑_{i\alpha }\}\right)`$. Consequently, its form does not need to be computed at each time step and the speed and efficiency of classical molecular dynamics is not compromised.
The problem with this approach lies in the choice for $`U_c\left(\{𝐑_{i\alpha }\}\right)`$. In principle, an exact mapping for some fundamental property such as the free energy $`F`$ can be imposed
$$F[U_i]=F_c[U_c].$$
(17)
where $`F_c`$ is the corresponding classical functional. Generally, such exact methods can be inverted only perturbatively and lead to a sequence of effective many ion interactions involving increasingly more particles. A more practical method is to assume pairwise additivity for effective point ”atoms”
$$U_c\left(\{𝐑_{i\alpha }\}\right)\frac{1}{2}\underset{k,j}{}\underset{\alpha }{\overset{N_k}{}}\underset{\beta }{\overset{N_j}{}}V_{kj}\left(\left|𝐑_{k\alpha }𝐑_{j\beta }\right|\right),$$
(18)
where $`k,i`$ label the species (ion or electron). The exact determination of the pair potentials $`V_{kj}\left(\left|𝐑_{k\alpha }𝐑_{j\beta }\right|\right)`$ is now much more restrictive as not all properties of interest have such a representation. Nevertheless pair properties such as the radial distribution functions $`g_{kj}\left(\left|𝐫\right|\right)`$ might be used to determine the pair potentials. As the $`g_{kj}\left(\left|𝐫\right|\right)`$ are unknown and difficult to calculate, the inversion is again difficult and not practical. Instead, experimental data is often used to fit a parameterized functional form chosen for the pair potentials. At this stage control over the approximation is lost and the method becomes phenomenological.
This phenomenological approach has been and remains a valuable tool of materials sciences. However, in the context of multiscale modeling it must be reconsidered carefully. First, it is recognized that there are local domains far from equilibrium where a purely classical potential cannot apply because of the inherent quantum chemistry active there (e.g., charge transfer and exchange). Conversely, there are large complementary domains in near equilibrium states where representation by an appropriate classical potential is possible. Multiscale modeling constructs distinct classical and quantum models of these subsystems and then requires fidelity at their interface. This is a severe test of the modeling assumptions in each subsystem. The approach proposed here confronts this issue directly in the construction of appropriate pair potentials for the problem considered.
### III.2 The Composite Quantum/Classical Solid
It is presumed that there is some method for identifying domains within the solid for which quantum chemical effects should be treated in detail. The quantum solid is then represented as a composite of two domains, the larger bulk in which a classical representation is used and a smaller ”reactive” domain in which the original quantum description is retained. The objective of the modeling described here is therefore to construct the potential function $`U_i`$ such that it gives an accurate description in both the reactive and non-reactive domains. This has two components, the determination of a pair potential for the forces on ions in the classical domain, and an accurate calculation of the charge density for the forces on ions in the quantum domain.
The classical and quantum domains are defined by labelling all ions as either classical or quantum, and associating spatial domains with the coordinates of each, denoted $`\{𝐑_{ci\alpha }\}`$ and $`\{𝐑_{qi\alpha }\}`$ respectively. The quantum domains are assumed to be small, to allow practical calculation of the electronic structure. In principle there could be several disconnected quantum domains, but to simplify the discussion we consider only one. It is assumed that initially the two sets are contiguous with a smooth interface and that diffusion or migration between them is not significant over the times of interest. In addition to the ions in the quantum domain, there are $`m`$ electrons where $`m`$ is determined by a condition on the charge of the quantum domain, here taken to be neutral. The total average electron charge density is then decomposed as
$$\overline{\rho }_e(𝐫,t)=\overline{\rho }_e(𝐫,t)\left(\chi _𝒬\left(𝐫\right)+\chi _𝒞\left(𝐫\right)\right)\overline{\rho }_{eq}(𝐫,t)+\overline{\rho }_{ec}(𝐫,t),$$
(19)
where $`\chi _𝒬\left(𝐫\right)`$ and $`\chi _𝒞\left(𝐫\right)`$ are characteristic functions for $`𝒬`$ and $`𝒞`$. The boundaries of the quantum spatial domain $`𝒬`$ are constrained by the choice of ions and total electron charge
$$me=_𝒬𝑑𝐫\overline{\rho }_{eq}(𝐫,t),$$
(20)
where $`𝒬`$ encloses all $`\{𝐑_{qi\alpha }\}`$ (this leaves some flexibility to define smooth interfaces). The complement of this is the classical domain $`𝒞`$. This gives a corresponding decomposition of the total ion potential energy function for the ions $`V_i`$ (ion - ion Coulomb interactions plus $`U_i`$ of (14)) so that (13) becomes
$$_tD_i+\{\left(K_i+V_{ic}+V_{iq}\right),D_i\}=0,$$
(21)
where $`K_i`$ is the ion kinetic energy and
$$V_{ic}=\frac{1}{2}_𝒞𝑑𝐫\rho _i\left(𝐫\right)\left(_𝒞𝑑𝐫^{}\frac{1}{\left|𝐫𝐫^{}\right|}\left(\rho _i\left(𝐫^{}\right)\delta \left(𝐫𝐫^{}\right)+\overline{\rho }_{ec}(𝐫^{},t)\right)+_𝒬𝑑𝐫^{}\frac{1}{\left|𝐫𝐫^{}\right|}\left(\rho _i\left(𝐫^{}\right)+\overline{\rho }_{eq}(𝐫^{},t)\right)\right),$$
(22)
and
$$V_{iq}=\frac{1}{2}_𝒬𝑑𝐫\rho _i\left(𝐫\right)\left(_𝒬𝑑𝐫^{}\frac{1}{\left|𝐫𝐫^{}\right|}\left(\rho _i\left(𝐫^{}\right)\delta \left(𝐫𝐫^{}\right)+\overline{\rho }_{eq}(𝐫,t)\right)+_𝒞𝑑𝐫^{}\frac{1}{\left|𝐫𝐫^{}\right|}\left(\rho _i\left(𝐫^{}\right)+\overline{\rho }_{ec}(𝐫^{},t)\right)\right).$$
(23)
The first terms of the integrands on the right sides of (22) and (23) represent the interactions of the ions with a given subsystem in the presence of the average electronic charge density of that subsystem. The second terms represent the interactions of those ions with their complementary subsystem. Half of the ion - ion potential energy between the two subsystems has been associated with each potential in this decomposition so that the total force acting on the quantum domain by the classical domain is equal and opposite to that on the classical domain due to the quantum domain.
The potential $`V_{ic}`$ is due to ions. By definition these ions are in near equilibrium states and therefore this part of the potential should be represented well by classical pair potentials of the form (18), with appropriately chosen parameters as discussed below. The potential $`V_{iq}`$ is due to ions in the quantum domain. As this is the domain that can be far from equilibrium the average electronic charge densities must be calculated in detail from the quantum description (15). There are two parts to this charge density affecting the ions in the quantum domain, that due to the $`m`$ electrons of the quantum domain $`\overline{\rho }_{eq}(𝐫,t)`$, and that due to the surrounding classical domain $`\rho _i\left(𝐫\right)+\overline{\rho }_{ec}(𝐫,t)`$. The equation governing the $`m`$ electrons of the quantum domain is coupled to this same classical domain average charge density. The scheme proposed here consists of modeling this charge density $`\overline{\rho }_{ec}(𝐫,t)`$ as an accurate and practical representation for the environment of the quantum domain, allowing calculation of $`\overline{\rho }_{eq}(𝐫,t)`$ and therefore determining both $`V_{iq}`$ and $`V_{ic}`$.
In summary, the multiscale model is obtained by replacing the ion \- ion contribution in the classical domain (the first term on the right side of (22)) by suitably parameterized pair potentials, and constructing an accurate approximate calculation of the electron density in the quantum domain. The latter entails a representation of the effects of the electron density in the classical domain on charges in the quantum domain. In this way the potentials of (22) and (23) are entirely determined. The details of this construction are addressed in the following sections.
## IV Description of the Classical Subsystem
### IV.1 Construction of the pair potentials
The proposed method for constructing a classical pair potential for use in multiscale modeling has several components: 1) it should be constructed for accuracy of the specific properties to be studied, 2) it should be ”trained” on quantum data generated in the same way as for the quantum domain at its interface, 3) it should predict accurate equilibrium structure, but include training on appropriate near equilibrium states as well, and 4) simplicity of form for parameterization and implementation in MD codes should be maintained.
The steps in constructing a potential are the following. First the specific quantum method to be used in the multiscale modeling is identified (below, transfer Hamiltonian or density functional) and applied to a large cluster or representative sample of the solid to be modeled. The forces on atoms for both equilibrium and near equilibrium states are then calculated quantum mechanically. Next, a simple functional form for the pair potentials with the correct physical shape (e.g., one of the existing phenomenological forms) is chosen and the parameters controlling that shape selected for optimization. The forces on the ions are calculated for these chosen potentials at the configurations used in the quantum force calculations, and compared with those quantum forces. The parameters are adjusted for a best fit to the quantum force data. When a good fit has been obtained, the potential energy at equilibrium is tested for stability using a gradient algorithm to establish that a local minimum of the potential has been obtained. Finally, the property of interest (e.g., some linear response) is tested by comparing its calculation using MD simulation with the fitted potentials and that with the quantum forces. If necessary, the fitting procedure can be repeated with differing weights for the quantum force data in equilibrium and near equilibrium states, or other additional input from the quantum calculations.
A more detailed discussion of the flexibility in this procedure is described for the example of the next section. It entails some art as well as science in choosing the optimization methods and applying them. In that example, it was found that some repeated combination of genetic algorithms and scaling provided the accuracy required. The primary result of this approach is a classical pair potential which gives the properties of interest by design, that match those being calculated quantum mechanically across the interface - for near equilibrium states.
### IV.2 Effects of the environment
The above construction describes the modeling of the first term on the right side of (22) where both the ions and the average electron charge density refer to the classical domain, assumed in a near equilibrium state. The second term involves a coupling to ions and average charge density (now the ground state charge density) in the quantum domain $`\overline{\rho }_{eq}(𝐫,t)`$. This charge density is computed for the forces in the quantum domain (next section) and hence the second term of (22) is known as well. Thus, the potential energy for the ions of the classical domain is determined from synthesized pair potentials among the ions and electrons of the classical domain, plus an interaction with the electrostatic potential of the ions and average electron charge density of the quantum domain. However, in the examples discussed below the coupling of the classical domain ions to the quantum domain is simplified by using the same pair potentials as for ions within the classical domain. This is expected to be quite accurate if the quantum domain charge density is not distorted very much from the near equilibrium states, since the potentials are trained to be equivalent to the charge density in the near equilibrium domain.
## V Description of the Quantum Subsystem
The quantum domain is a subsystem of $`m`$ electrons localized about the designated ions defining that domain. Consider the reduced density operator for $`m`$ electrons defined by
$$D_e^{(m)}=Tr_e^{N_em}D_e$$
(24)
More specifically, this partial trace is defined in coordinate representation by
$$𝐫_1,..,𝐫_m\left|D_e^{(m)}\right|𝐫_1^{},..,𝐫_m^{}=d𝐫_{m+1}..d𝐫_{N_e}𝐫_1,..,𝐫_m,𝐫_m,..𝐫_{N_e}\left|D_e\right|𝐫_1^{},..,𝐫_m^{},𝐫_m,..𝐫_{N_e}$$
(25)
Clearly the full exchange symmetry among all $`N_e`$ electrons is preserved. However, this reduced density operator is not specific to the quantum domain defined above only. For example, the diagonal elements $`𝐫_1,..,𝐫_m\left|D_q\right|𝐫_1,..,𝐫_m`$ give the probability density to find $`m`$ electrons at the specified positions, and the latter can be chosen anywhere inside the system. Thus, only when the positions are restricted to the quantum domain $`𝒬`$ does this reduced density operator represent the electrons of that domain. Similarly, if $`\rho _e^{\left(m\right)}\left(𝐫\right)`$ is the charge density operator for $`m`$ electrons its average is
$$\overline{\rho }_e^{(m)}\left(𝐫\right)=Tr_e\rho _e^{\left(m\right)}\left(𝐫\right)D_e=Tr_e^m\rho _e^{\left(m\right)}\left(𝐫\right)D_e^{(m)}$$
(26)
where the trace in the second equality is over $`m`$ degrees of freedom. This average charge density represents the average contribution of $`m`$ electrons at any point $`𝐫`$. Both $`𝐫_1,..,𝐫_m\left|D_q\right|𝐫_1,..,𝐫_m`$ and $`\overline{\rho }_e^{(m)}\left(𝐫\right)`$ change with $`𝐫`$ since there is an absolute reference background set by the functional dependence on the ion charge density. Consequently, in all of the following discussion of this section $`𝐫_1,..,𝐫_m\left|D_q\right|𝐫_1,..,𝐫_m`$ and $`\overline{\rho }_e^{(m)}\left(𝐫\right)`$ are considered only for positions within the chosen quantum domain. Accordingly, this coordinate representation $`\left\{\right|𝐫_1,..,𝐫_m\}`$ defines an $`m`$ electron Hilbert space of functions defined over the quantum domain $`𝒬`$. In this context $`D_e^{(m)}`$ becomes the reduced density operator for the quantum domain and $`\rho _e^{\left(m\right)}\left(𝐫\right)`$ its charge density operator
$$D_e^{(m)}D_q,\rho _e^{\left(m\right)}\left(𝐫\right)\rho _{eq}\left(𝐫\right).$$
(27)
The equation determining the reduced density operator for the quantum subsystem follows directly from this definition and Eq. (15) for $`D_e`$
$$\frac{i}{\text{h}\text{ }\text{ }}[\left(K_e+V_{eq}\right),D_q]=0.$$
(28)
where $`K_e`$ is the kinetic energy for the $`m`$ electrons and
$$V_{eq}=\frac{1}{2}_𝒬𝑑𝐫\rho _{eq}\left(𝐫\right)\left(_𝒬𝑑𝐫^{}\frac{1}{\left|𝐫𝐫^{}\right|}\left(\rho _{eq}\left(𝐫^{}\right)\delta \left(𝐫𝐫^{}\right)+\overline{\rho }_i\left(𝐫^{}\right)\right)+_𝒞𝑑𝐫^{}\frac{1}{\left|𝐫𝐫^{}\right|}\left(\overline{\rho }_{ec}\left(𝐫^{}\right)+\overline{\rho }_i\left(𝐫^{}\right)\right)\right).$$
(29)
(The contribution from $`\overline{\rho }_{ec}\left(𝐫^{}\right)`$ in the second term actually should be a conditional average; the same mean field approximation as in (11) has been introduced here for consistency). The first term on the left side of (29) describes the isolated quantum subsystem, composed of the Coulomb interactions among the $`m`$ electrons and their coupling to the average charge density of the ions in the quantum domain. The second term is the interaction of these electrons with their environment, the total average charge density of the classical domain.
There are two distinct types of contributions from this charge density of the classical domain. The first is associated with a subset of ions at the border of the QM/CM domains which describe chemical bonds in the full quantum solid. These ions locate regions where there is a highly localized electron charge density shared with the quantum domain, including both strong correlation and exchange effects, localized and non-uniform. The second type of contribution is associated with the remaining border ions and those charges more distant. In this second case the quantum correlation and exchange effects of the first type are much weaker, and the dominant effect is that of a polarized charge density due to the ions and electrons.
These two types of effects of the environment can be identified by separating the total charge density for the classical domain in (29) into domains centered on the border ions responsible for bonding, and their complement
$$\overline{\rho }_c\left(𝐫\right)\overline{\rho }_{ec}\left(𝐫\right)+\overline{\rho }_i\left(𝐫\right)=\underset{\alpha 𝒱}{}\chi \left(\left|𝐫𝐑_\alpha \right|\right)\overline{\rho }_c\left(𝐫\right)+\mathrm{\Delta }\overline{\rho }_c\left(𝐫\right)$$
(30)
where it is understood that $`𝐫`$ is in the classical domain. The set of border ions for which a bond has been broken in identifying the quantum domain is denoted by $`𝒱`$. Also, $`\chi \left(\left|𝐫𝐑_\alpha \right|\right)`$ is a characteristic function specifying a domain centered on $`𝐑_\alpha `$ such that it does not overlap neighboring ions. Its size is taken large enough to incorporate the bound electrons forming the ”atom” for this ion in the quantum solid. The second term $`\mathrm{\Delta }\overline{\rho }_c\left(𝐫\right)`$ is the charge density for all remaining ions and electrons of the classical environment.
### V.1 Coulomb effects of environment
By its definition, the averaged charge density $`\mathrm{\Delta }\overline{\rho }_c\left(𝐫\right)`$ does not include the contributions to chemical bonding with the quantum domain required for its valency. Hence the electrostatic potential associated with it can be expected to have a regular multipole expansion
$$_𝒬𝑑𝐫^{}\frac{1}{\left|𝐫𝐫^{}\right|}\mathrm{\Delta }\overline{\rho }_c\left(𝐫\right)\frac{\widehat{𝐫}𝐝}{r^2}+..$$
(31)
The leading monopole term is zero due to charge neutrality of the classical subsystem, and the dipole moment $`𝐝`$ for the entire environment is
$$𝐝=𝑑𝐫^{}𝐫^{}\mathrm{\Delta }\overline{\rho }_c\left(𝐫\right).$$
(32)
These results are still quite formal but they provide the basis for the phenomenology proposed for this part of the modeling: Replace all effects of the classical environment on the quantum system, exclusive of bonding, by an effective dipole representing the polarization of the medium by the quantum domain. The origin of this dipole in the above analysis shows that in general it will depend on the geometry and the state of both the classical and quantum domain (e.g., it will change under conditions of strain). In the phenomenological application of this prescription the dipole must be supplied by some simpler means since the electronic contribution to $`\mathrm{\Delta }\overline{\rho }_c\left(𝐫\right)`$ is not known (for example, see reference Mallik2004 ).
### V.2 Electron exchange with environment
The contributions from $`\overline{\rho }_c\left(𝐫\right)`$ in the regions where bonds have been cut require a more detailed treatment. Clearly, a necessary condition is that valence saturation should be restored in the quantum domain. However, that is not sufficient to assure the correct charge density there nor the correct forces within that domain. Instead, the charge density near the border ions responsible for bonding should induce a realistic charge density within the quantum domain. To see how this can be done first write the contribution from one such border ion to (30) as
$$𝑑𝐫^{}\frac{1}{\left|𝐫𝐫^{}\right|}\chi \left(\left|𝐫^{}𝐑_\alpha \right|\right)\overline{\rho }_c\left(𝐫\right)\varphi _{p\alpha }\left(𝐫\right)$$
(33)
The potential $`\varphi _{p\alpha }\left(𝐫\right)`$ represents the actual electrostatic and exchange effects of the nucleus of the border of the quantum domain. This is comprised largely of the ion at the site plus its closed shell electrons distorted by exchange and correlation effects of the site in the quantum domain with which it is bonding. Consequently, an appropriate pseudo-potential is introduced at each such border ion whose behavior in the direction of the quantum system is the same as that of the actual ion. These ions plus their pseudo-potentials will be called ”pseudo atoms”.
To accomplish this, an appropriate candidate for the pseudo atom is chosen based on valency of the particular pair of border ion and its neighbor in the quantum domain. Next, a large molecule or small cluster containing that bonding pair is chosen for training the pseudo atom. The bond is then broken and the relevant member of the pair replaced by the pseudo atom. The training consists of parameterizing the pseudo atoms’ effective potential to give the same forces as in the original cluster. For example, the pseudo atom might consist of an ion plus its closed shell electrons chosen to satisfy the valency for the bond broken. The adjustable parameters could refer to a characterization of the closed shell electron distribution. In this way, it is assured that the pseudo atom not only gives the correct saturation of the dangling bond but also reproduces the forces within the cluster and hence gives a realistic representation of the charge distribution between the chosen pair. For the small training cluster (e.g., Figure 6), forces are determined both at equilibrium and under strain to within about one tenth percent. The primary assumption is that the bonding of interest is a local effect, so that the pseudo atom trained in the cluster will have a similar accuracy when used in the bulk solid . The training depends on the particular quantum method used to describe the solid, and is illustrated in more detail for the specific example considered in the later sections.
In summary, the environmental effects on the quantum system are accounted for approximately by a dipole representing its polarization and a collection of pseudo atoms located at the sites of ions where bonds have been cut
$$V_{eq}\frac{1}{2}_𝒬𝑑𝐫\rho _{eq}\left(𝐫\right)\left(_𝒬𝑑𝐫^{}\frac{1}{\left|𝐫𝐫^{}\right|}\left(\rho _{eq}\left(𝐫^{}\right)\delta \left(𝐫𝐫^{}\right)+\overline{\rho }_i\left(𝐫^{}\right)\right)+\frac{\widehat{𝐫}𝐝}{r^2}+\underset{\alpha 𝒱}{}\varphi _{p\alpha }\left(𝐫\right)\right).$$
(34)
This model now allows solution to (28) for the electron distribution the quantum domain, including its coupling to the classical environment. Then $`\overline{\rho }_{eq}\left(𝐫\right)`$ is calculated from (26) and (27). Finally, the desired potentials $`V_{ic}`$ and $`V_{iq}`$ of (22) and (23) are fully determined
$$V_{ic}=\frac{1}{2}\underset{i,j}{}\underset{\alpha }{\overset{N_i}{}}\underset{\beta }{\overset{N_j}{}}V_{ij}\left(\left|𝐑_{i\alpha }𝐑_{j\beta }\right|\right)+\frac{1}{2}_𝒞𝑑𝐫\rho _i\left(𝐫\right)_𝒬𝑑𝐫^{}\frac{1}{\left|𝐫𝐫^{}\right|}\left(\rho _i\left(𝐫^{}\right)+\overline{\rho }_{eq}(𝐫^{},t)\right),$$
(35)
and
$$V_{iq}=\frac{1}{2}_𝒬𝑑𝐫\rho _i\left(𝐫\right)\left(_𝒬𝑑𝐫^{}\frac{1}{\left|𝐫𝐫^{}\right|}\left(\rho _i\left(𝐫^{}\right)\delta \left(𝐫𝐫^{}\right)+\overline{\rho }_{eq}(𝐫,t)\right)+\frac{\widehat{𝐫}𝐝}{r^2}+\underset{\alpha 𝒱}{}\varphi _{p\alpha }\left(𝐫\right)\right).$$
(36)
The solution to the classical ion motion (21) can then proceed via MD simulation. This completes the proposed scheme for multiscale modeling of the idealized quantum solid. Rather than critique further the assumptions made in this abstract context, the scheme is discussed in more detail for a specific application to an $`Si0_2`$ nanorod.
## VI $`Si0_2`$ Nanorod - A Critical Test
The modeling scheme of the previous sections is now illustrated in detail and tested critically for a model solid consisting of a silica nanorod Zhu2003 , shown in Figure 1. It consists of 108 Si and O ions with the stoichiometric ratio of one silicon to two oxygens, a stack of six $`Si_6O_6`$ rings sharing both above and below a ring of oxygen atoms. To saturate overall valency the nanorod is terminated with end cap rings whereby each silicon atom of the terminating ring is connected by bridging oxygen or two interstitial oxygens. The size of this nanorod can be readily adjusted by adding or removing $`Si_6O_6`$ planar rings with corresponding O atoms. Many of the results reported here have been studied for larger rods as well Zhu2003 . As noted in the Introduction, this system exhibits strain to fracture (Figure 1), yet it is small enough for practical application of the chosen quantum mechanics to the whole rod. Consequently, the results of modeling it as a composite rod can be tested quantitatively against the ”exact” results. That is the objective of this section.
The decomposition of the rod into classical and quantum domains is accomplished by identifying the ions of one of the rings near the center as the QM domain. The rest of the ions define the classical domain. The actual geometry of the interface is chosen by charge neutrality of the quantum domain, as given by Eq.(20).
### VI.1 Transfer Hamiltonian Quantum Mechanics
At the quantum level, each possible choice for the quantum chemical method entails different approximations used to optimize the accuracy and speed of the calculations. Each approximation has its own advantages and limitations, but a general embedding scheme should be insensitive to these. Much of the previous work on multiscale modeling is based on tight binding models Broughton1999 ; Rudd2000 because of their ability to treat hundreds of atoms in the QM domain. However, the Hamiltonian in these models is oversimplified and cannot account for charge transfer, which is essential to bond breaking. A more complete quantum description is considered here. It is appropriate to stress at this point that the objective of the following sections is to demonstrate and test the modeling scheme for a given quantum description. That is a separate issue from the question of how accurate that description may be. To test how robust the proposed embedding scheme is, two different methods are chosen for the underlying quantum mechanical description. The first method is the Transfer Hamiltonian (TH) - Neglect of Diatomic Differential Overlap (NDDO) Taylor2003 . The second method is the Born-Oppenheimer local spin density Barnett1993 within density functional theory using a generalized gradient approximation. The results obtained from the TH are described first.
The goal of the TH strategy is to provide forces for realistic MD simulations of the quality of the coupled cluster singles and doubles method Bartlett1995 in computationally accessible times Taylor2004 . The TH is a low rank, self - consistent, quantum mechanical single particle operator whose matrix elements are given in terms of parameterizable functions. The functional form of the TH can be chosen to be of any semi-empirical functional form, but here the NDDO form is chosen because of its better qualitative description over Intermediate Neglect of Differential Overlap (INDO) Hamiltonians Hsiao2001 . The NDDO approximation restricts charge distribution to basis functions on the same center, so this has at most two-atom interactions. Because of its NDDO form it is anticipated that the parameters fit to coupled cluster data are saturated for small cluster size and that these parameters can then be transferred to extended systems. This transfer might be considered an example of serial multiscale-modeling.
To apply the TH to silica systems for studying complex processes such as fracture or hydrolysis where the role of electrons is important, Taylor et al. Taylor2003 parameterized the TH for the following reactions $`Si_2H_62SiH_3`$, $`Si_2O_7H_6Si(OH)_3+SiO(OH)_3`$ (radical mode of dissociation), and $`Si_2O_7H_6`$ $`{}_{}{}^{+}Si(OH)_3+^{}OSi(OH)_3`$ (ionic mode of dissociation). By describing the two different dissociation pathways for $`Si_2O_7H_6`$, electron state specificity is systematically introduced into the TH model. Since fracture is expected to involve a variety of dissociations this TH model is expected to be appropriate for study of strained conditions as is done here. Thus the parameterized TH model is expected to provide accurate results for such processes, while requiring only the computational intensity of semi-empirical methods which are orders of magnitude less demanding than coupled cluster calculations on the same systems. In the next few subsections, the charge densities and forces that arise from the TH are the reference data by which the proposed method for multiscale modeling is judged.
### VI.2 Pair Potentials for the Classical Domain
The method that we propose for constructing a classical pair potential has been previously presented in detail Mallik20042 , so only a brief summary is given here. The functional form for the potential is chosen to be one that is commonly used in the literature. In the case of silica, this is the TTAM Tsuneyuki1998 and BKS Beest1990 forms which are pair potentials comprised of a Coulombic interaction, a short-range exponential repulsion, and a van der Waals attraction (the last two terms are collectively known as the ’Buckingham’ potential)
$$V_{ij}(r)=\frac{q_iq_j}{r}+a_{ij}e^{b_{ij}r}\frac{c_{ij}}{r^6}$$
(37)
There are 10 free parameters as the charge neutrality of the $`SiO_2`$ unit requires that $`q_{Si}=2q_O`$, so that $`a_{ij}`$, $`b_{ij}`$, and $`c_{ij}`$ for each atom pair and one charge define the potential. For the purposes here, this potential is used to generate the dynamics (MD simulations) so that forces have been used as the property ”training” the potential (determining the free parameters). More specifically, force data from the silica nanorod calculated with the TH at equilibrium and a few small strain configurations have been used.
Fitting a set of ten parameters from hundreds of data points requires an optimization approach that is capable of exploring a large parameter space efficiently without being trapped in local minima. For this task, a genetic algorithm (GA) has been chosen as detailed in reference Mallik20042 . The GA generates a number of sets of parameters of a given fitness which are then tested using a standard geometry optimization technique (BFGS Zhu1994 ) to ensure that the parameterization gives a stable equilibrium configuration. The parameter set that gives a geometry closest to the TH nanorod is then chosen and rescaled to reproduce the equilibrium structure to within a few percent. (See the appendix of reference Mallik20042 for details of the rescaling used.) Table 1 presents the final parameters obtained compared to the TTAM and BKS parameters. Table 2 shows that the new classical potential reproduces the nanorod structure to within a few percent compared to larger discrepancies with the two standard potentials. The shapes of these three potentials are shown in Figure 2 for the Si-O pair interaction.
Finally, the small strain elastic behavior of this new potential is studied. This is accomplished using MD simulation of uniaxial strain along the long axis (z) of the silica nanorod at a strain rate of 25 m/s. The integration is done using 2 fs time steps and a temperature of 10 K enforced by velocity rescaling. Figure 3 shows the resulting stress-strain curves for small strains up to $`5\%`$. It is noted that the new potential indeed reproduces the linear elastic response (Young’s modulus) behavior of the underlying TH while the TTAM and BKS potentials are somewhat stiffer.
A key to successful multiscale modeling lies in the consistent embedding of a QM domain in its classical MD region. This requires that the CM region have the same structure and elastic properties as the QM domain. The potential obtained in this section meets these criteria for the silica nanorod. Essentially, a classical representation of the TH silica nanorod has been found for both equilibrium and near-equilibrium configurations.
### VI.3 The Quantum Domain
A proper description of the quantum domain requires incorporating two kinds of environmental effects: a short-range electronic exchange interaction and long-range Coulombic interactions. Most previous studies of QM/CM simulations have focused on the former and have neglected the long-range interactions. It will be shown here that both kinds of interactions must be taken into account for an accurate description of forces and charge densities in the QM domain. This subsection first reviews briefly the various termination schemes used in the literature to treat the bond cutting region. Then, the method proposed here for construction of a pseudo-atom to saturate dangling bonds is described, together with modeling the rest of the CM environment by lowest order multipoles. The charge densities and forces in the QM domain obtained from our scheme of pseudo-atoms and dipoles with those obtained from the conventionally used link atoms (bond saturation with hydrogen atoms) and dipoles are compared to the ”exact” TH reference data for the entire nanorod.
Over the past decade there have been numerous proposals for different types of termination schemes to accommodate dangling bonds. Only a few of the more relevant ones are noted for context. (i) Link Atom method: This commonly used method is the Link Atom method, presented by Singh and Kollman Singh1986 . In this method, hydrogen atoms are added to the CM side of broken covalent bond to satisfy the valency of the QM system. There are many variations within the implementation of the LA method, for example, the double Link Atom method Das2002 , the Add-Remove Link Atom method Swart2003 or the Scaled-Position-Link-Atom Method (SPLAM)Eichinger1999 . In the present work, the hydrogen atoms are placed at a fixed distance of 0.97 Å from the cut Si bond. (ii) Connection Atom method Antes1999 : A connection atom is developed to saturate a C-C bond, such that the connection atom mimics the effect of a methyl group. This connection atom interacts with the QM atom quantum mechanically and the interactions with the CM atoms are handled classically using a carbon force field. The parameters of the connection atom are determined using semi-empirical methods Stewart1990 such as AM1, MNDO or PM3 designed to reproduce theoretical QM data for energies, geometry and net charges. About 30 different methyl hydrocarbons were used as reference molecules. The parameters adjusted are the orbital exponent , one-center one-electron energy , one center two-electron integral , resonance parameter , and repulsion term . The mathematical functional forms of these parameters can be found in any book on semi-empirical theory Stewart1990 or papers by Dewar and Theil Dewar1976 .
There are numerous other termination schemes for the QM/CM boundary like the ’pseudobond’ scheme Zhang1999 , ‘IMOMM’ Maseras1995 , Humbel1996 and ’ONIOM’Svensson1998 , Maseras1995 , Dapprish1999 procedures , ’effective group potential (EGP)’ Poteau2001 . However these methods are not discussed here.
The method proposed here will be referred to as the ”pseudo-atom” method. In the present case it is based on the TH approach for saturating a bond terminating in Si (in silica systems). The training cluster is a pyrosilicic acid molecule (Figure 4). The part of the molecule within the dotted lines is replaced by a fluorine (F) atom whose NDDO parameters are then adjusted to give the correct QM forces (which implies correct geometry as well) and charge density in other parts of the molecule (outside the dotted lines). The pseudo atom is placed at the same position as the neighboring CM atom (O in this case) in the bond being cut and hence geometry optimization as in the LA method is not required. The specific NDDO parameters modified are: one-center-one-electron integrals , Coulomb integrals, exchange integral, and two-center-one-electron resonance integrals Dewar1976 . This method is similar to the previously described connection atom developed in reference Antes1999 . However the advantage of using pseudo-atoms is that they are trained to give the correct QM forces and charge densities for both the equilibrium Si-O bond being cut and for small distortions (up to $`34\%`$ from the equilibrium) in the Si-O bond length as well. This allows use of the pseudo-atoms even while studying dynamics in the system.
To test this method, the actual charge density from the TH method is calculated instead of the more commonly used Mulliken populations Mulliken1955 , which are known to have several common problems ( e.g., equal apportioning of electrons between pairs of atoms, even if their electronegetavities are very different).
The remainder of the environment is represented as two dipoles for the top and bottom portions of the rod (Figure 5). The values of the dipole have been calculated using the TH-NDDO charge density for these two portions of the rod. These domains are taken to be charge neutral, but are polarized by the presence of the QM domain. The validity of the approximation by dipoles has been checked by comparing the force on an Si nuclei of the ring due to all charges and that due to the dipole, with excellent agreement. Constructing the full system as a composite of a QM domain terminated by pseudo atoms and embedded in dipolar fields leads to the forces displayed in Figure 6. It is found that the normalized difference of charge density is reproduced to within $`0.1\%`$ in the plane of the ring Mallik2004 .
This scheme was found to be applicable to both equilibrium and strained configurations, as illustrated in Figure 6 for the following cases: a) equilibrium, b) the ring of the QM domain radially expanded by $`5\%`$, and c) a distorted ring in which one Si atom is radially pushed out and one Si pushed in. Also shown are the corresponding results for a longer 10 ring rod (cases d),e), and f)).
The pseudo-atoms and the dipoles reproduce the force in the QM domain within $`1\%`$ error. In contrast, the LA placed at 0.97 Å from the terminating Si atom and along the Si-O bond lead to a large force on the Si atom at equilibrium, and generally poor results for the strained cases. One might argue that the LA method could be improved if placed at ”optimal” positions. To test this, both the Si-H bond length as well as the alignment of the LA was varied to give minimum force on the Si atom. This was obtained when the LA is placed at a distance of 1.45 Å from the silicon and the bond angle is decreased by about 5 degrees. Although the LA at this position gives forces comparable to those for pseudo-atoms plus the dipoles, it fails to reproduce the correct charge densities.
### VI.4 The Composite Rod
The composite rod is built by embedding the QM region in its CM environment as described in Section III. The forces on the atoms in the CM region are calculated from the pair potential developed in IV B above, while the forces on the atoms in the QM domain are obtained from the charge density calculated by the TH method as described in VI C . An important test of the composite rod is its indistinguishability from the TH rod for near equilibrium states, i.e. its structure and elastic properties.
The structural properties (bond lengths, bond angles) have an accuracy comparable to that of the purely classical rod described using this pair potential (Table 2). Consequently, attention here will be focused on the elastic properties. Figure 7 shows the stress-strain behavior of the nanorod obtained from three different methods: (i) quantum mechanics TH method for the entire rod, (ii) pair potentials for the entire rod, and (iii) the composite rod constructed as described above. The three overlaying curves (measured at 0.01K) indicate that the composite rod is identical to the rod obtained from TH and the pair potential nanorod in terms of small strain elastic properties and structure. The stress-strain results shows the success of our multiscale method indicating that the composite rod is indistinguishable from the underlying quantum mechanics for states near equilibrium.
### VI.5 Notched Nanorod
The elastic properties of the composite rod obtained from both of these potentials do not agree beyond $`10\%`$ strain. This is because the pseudo-atoms, trained at regions only close to the equilibrium configuration, fail to give the correct charge densities at such high strains. This diagnosis was checked by comparing the charge density in the QM domain of the $`12\%`$ strained rod with that of equilibrium configuration and a difference of $`6\%`$ was found. Retraining of the pseudo-atoms improve the stress-strain performance of the composite rod beyond a strain of $`10\%`$ is possible, but this was not done because of the large strains involved. Real systems like glasses have many inherent defects which act as stress concentrators that cause the material to break at much lower strains than observed for the nanorod (with a yield stress of about 190 GPa).
To illustrate this, a defect notch was placed in the 108 atom nanorod by removal of an oxygen atom as shown in Figure 8. The MD stress-strain curves for this notched rod were found using TH quantum mechanics, the trained classical potential, and the composite. We note that the presence of only a small defect can significantly reduce the yield stress of the material and make it more prone to fracture. The TH curve for the defect-free rod is plotted in the same figure to contrast the value of the yield stress. As can be seen there is a reduction of ~60 GPa in the yield stress.
For the composite rod in this case, the QM domain was chosen to consist of 2 silica planes and the intermediate 5 oxygen atoms (see Figure 8), so that the defect could be located in the QM region. The stress - strain curve for this composite notched rod agrees well with that for the TH quantum calculation up to 10% strain. Above 4% strain the curve for the composite notched nanorod follows that of the TH instead that of the trained classical potential nanorod, showing that the composite rod is representing the ”real” material. This is exactly what is required of multiscale modeling.
### VI.6 DFT Quantum Mechanics
A proper multiscale modeling procedure should be independent of the choice of underlying quantum mechanical method. To test the method proposed here, the analysis of subsections (B)-(D) is repeated using a density functional theory (DFT) instead of the Transfer Hamiltonian as the quantum mechanical approximation. The primary results of this section are a confirmation that the isolation of the QM domain with pseudo-atoms and dipoles, plus the construction of a classical potential based on the DFT forces leads to accuracies of the same quality as those described already using the TH quantum mechanics. Hence, the multiscale modeling scheme is faithful to the chosen form for the underlying quantum mechanics in both cases. The DFT code used is a parallel multiscale program package, known as Born-Oppenheimer molecular dynamics “BOMD” Barnett1993 . It is a generalized gradient approximation (GGA) within local spin DFT (LSDFT). A Troullier-Martin Pseudo-Potential Troullier1991a ; Troullier1991b is used for the effect of the chemically inert core states on the valence states. The code employs the dual space formalism for calculation of the DFT energy. A plane wave basis set and cut-off energy of 30.84 Rydbergs is chosen.
A pseudo-atom for the quantum domain is constructed based on parameterization of the Troullier-Martin (TM) pseudopotential, using a cut-off radius of $`1.5\AA `$. Once again the pyrosilicic acid molecule is chosen for parameterization of the fluorine-like pseudo atom and its position is constrained to be at the same place as the O atom (Figure 4). Unlike the TH-NDDO method in which both the electron-ion and electron-electron interaction parameters could be changed , in DFT only the electron-ion interactions can be modified. The three options to alter these interactions for the F atom using the TM pseudopotential are: (i) the core charge on F, (ii) omission of the non-local part in the potential, and (iii) switching the local and non-local part between s and p orbitals. All three possible choices and their combinations were explored to find the optimal reproduction of forces on the terminated Si atom in the pyrosilicic acid. The best results were obtained when the core charge is 7.0 and the non-local part is omitted. This pseudo-atom was then applied to the composite nanorod constructed as above (Figure 5). The values of the dipoles were recalculated from the charge density obtained using the DFT results for the CM portions of the rod. The force on a Si atom in the QM domain was calculated for the rod in equilibrium and all the strained cases considered in subsection C. The results obtained from a DFT calculation on the whole rod are now taken to be the ”exact” reference forces. It is found that the forces and the charge densities in the QM domain can be generated using DFT pseudo-atoms and dipoles to $`1\%`$ accuracy.
Next, a new potential having the same form as TTAM is constructed to predict the same structure and elastic properties as for the QM rod. A GA with DFT force data up to $`4\%`$ expansion followed by the scaling procedure is used as described in section VI B above to find the parameters for the potential. The charge on the ions is lower (as was found for TH quantum mechanics) than that given in the TTAM and BKS potentials. Also, the van der Waals interaction is much weaker than in those potentials. This is expected since the DFT forces fail to represent this effect, falling exponentially rather than algebraically with separation. The parameters for the DFT potential and a comparison of the resulting nanorod structure with that for the BKS and TTAM nanorods are given in reference Mallik20042 . The agreement between the results from the DFT potential and DFT quantum mechanics is similar to that found in Table 2 for the TH quantum mechanics.
A stress strain curve for the entire rod using MD with DFT quantum mechanics is computationally too intensive. Instead, only selected equilibrium and adiabatic strain configurations were calculated with DFT. The equilibrium structure was determined by sequential DFT calculations and nuclear relaxation to find the minimum energy configuration. The strained configurations were obtained from an affine transformation of the minimum energy configuration by 1, 2, 3, and 4 %, with a single DFT calculation of forces at each of the expanded configurations. Then the average force on the Si atom in each ring of the DFT nanorod was computed for these four cases. The stresses for adiabatic configurations using DFT for the entire rod, the constructed DFT potential, and the TTAM potential were then compared (the values of stress obtained using BKS potential are similar to those of TTAM potential). The results are shown in Figure 9.
These results suggest that the multiscale modeling method proposed here is accurate for a wide range of possible choices for the underlying quantum mechanical method employed.
## VII Summary and Discussion
The objectives here have been three-fold. The first was to describe the formal quantum structure for a real solid from which a practical model should be obtained through a sequence of well-identified (if not fully controlled) approximations. The second was to propose a method for partitioning this structure into classical and quantum domains while preserving the properties of interest for the quantum structure in the replica composite solid. The final objective was a quantitative test of the proposed method by its application to a non trivial mesoscopic ”solid”, the $`SiO_2`$ nanorod.
For the first objective, coupled Liouville - von Neumann equations for the reduced density operators for the ions and electrons of the system were considered. Each is coupled to the other through the mean charge density of the complementary subsystem. Determination of these charge densities then becomes the central problem for further analysis. For the heavy ion component analysis by classical MD simulation provides a practical and accurate approach. The necessity for a detailed quantum treatment of the electronic charge density is the primary bottleneck for progress. For bulk samples of interest, direct application accurate quantum chemical methods are precluded, so composite constructs with smaller quantum subdomains are a potentially fruitful compromise. Here, the isolation of a quantum subdomain was accomplished by defining the reduced density matrix for electrons in the vicinity of a selected small set of ions. The corresponding Liouville - von Neumann equation describes the dynamics of those electrons and ions coupled to the remainder of the solid (the classical domain). Its environment is entirely characterized by its mean charge density. This formulation and emphasis on the charge density is a guiding feature of the subsequent approximations and modeling.
The specific method for constructing the composite solid consists of modeling the charge densities of the environment for the two subsystems. For the environment of the quantum domain the charge density is separated into that border component responsible for electronic bonding across the border, and a longer range Coulomb interaction. The former is treated in detail by pseudo atoms trained on smaller clusters to provide highly accurate local forces and charge densities. The latter is represented by leading order terms in a multi-pole expansion. The interactions within the classical domain are described phenomenologically by pair potentials. These pair potentials are designed specifically to represent selected properties determined by the quantum chemical method assumed to represent the solid. This presumes conditions such that the chosen properties have such a classical representation (e.g., near equilibrium states), and the method constitutes a ”tuning” of the potential to match the quantum description on the other side of the border.
There are several critical tests of this approach, all requiring a benchmark system for which the ”exact” quantum data for the entire solid is required. This system is provided here by the nanorod for which the global quantum calculations are possible. Then, choosing a central ring as the quantum subdoman and the remainder as the classical domain many quantiative tests are possible. Some of these provided here are:
1. The charge density and forces in the central ring computed from the model are accurate to within one percent, both at equilibrium and strains up to 5%.
2. The pair potential fit to quantum data was used to construct the entire nanorod. The resulting structure and elastic properties up to 5% strain were indistinguishable from those of the quantum nanorod.
3. The proposed method was used to construct a composite classical / quantum nanorod. Again the structure and elastic properties were indistinguishable from those of the quantum nanorod.
These and other results presented above constitute a demonstration that the multiscale modeling is not creating a new solid, but rather is faithful to the real system of interest. Similar (and perhaps more physical) results were obtained for the nanorod with a defect (missing Oxygen). Finally, the entire multiscale analysis was repeated using a quite different choice for the underlying quantum mechanics with the same degree of accuracy.
The predictions of multiscale modeling for conditions of interest, states far from equilibrium, rely on the assumption that the properties calculated in the quantum subdomain are indeed those of the given quantum structure. This means that the quantum imbedding and representation of the classical domain are ”passive”, i.e. no new physical effects have been introduced. The method described here manifestly satisfies this constraint for the benchmark nanorod, and provides some concrete support for its use in realistic applications.
## VIII Acknowledgements
This research was supported by NSF-ITR under Grant No. DMR-0325553. The authors are indebted to S. B. Trickey, K. Muralidharan, and D. E. Taylor for helpful discussions and assistance. |
warning/0507/math0507215.html | ar5iv | text | # Invariant Metrics and Laplacians on Siegel-Jacobi Space
## 1. Introduction
For a given fixed positive integer $`n`$, we let
$$_n=\{Z^{(n,n)}|Z=^tZ,\text{Im}Z>0\}$$
be the Siegel upper half plane of degree $`n`$ and let
$$Sp(n,)=\{M^{(2n,2n)}|^tMJ_nM=J_n\}$$
be the symplectic group of degree $`n`$, where
$$J_n=\left(\begin{array}{cc}0& E_n\\ E_n& 0\end{array}\right).$$
We see that $`Sp(n,)`$ acts on $`_n`$ transitively by
(1.1)
$$MZ=(AZ+B)(CZ+D)^1,$$
where $`M=\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)Sp(n,)`$ and $`Z_n.`$
For two positive integers $`n`$ and $`m`$, we consider the Heisenberg group
$$H_{}^{(n,m)}=\{(\lambda ,\mu ;\kappa )|\lambda ,\mu ^{(m,n)},\kappa ^{(m,m)},\kappa +\mu ^t\lambda \text{symmetric}\}$$
endowed with the following multiplication law
$$(\lambda ,\mu ;\kappa )(\lambda ^{},\mu ^{};\kappa ^{})=(\lambda +\lambda ^{},\mu +\mu ^{};\kappa +\kappa ^{}+\lambda ^t\mu ^{}\mu ^t\lambda ^{}).$$
We define the semidirect product of $`Sp(n,)`$ and $`H_{}^{(n,m)}`$
$$G^J:=Sp(n,)H_{}^{(n,m)}$$
endowed with the following multiplication law
$$(M,(\lambda ,\mu ;\kappa ))(M^{},(\lambda ^{},\mu ^{};\kappa ^{}))=(MM^{},(\stackrel{~}{\lambda }+\lambda ^{},\stackrel{~}{\mu }+\mu ^{};\kappa +\kappa ^{}+\stackrel{~}{\lambda }^t\mu ^{}\stackrel{~}{\mu }^t\lambda ^{}))$$
with $`M,M^{}Sp(n,),(\lambda ,\mu ;\kappa ),(\lambda ^{},\mu ^{};\kappa ^{})H_{}^{(n,m)}`$ and $`(\stackrel{~}{\lambda },\stackrel{~}{\mu })=(\lambda ,\mu )M^{}`$. We call this group $`G^J`$ the Jacobi group of degree $`n`$ and index $`m`$. We have the natural action of $`G^J`$ on $`_n\times ^{(m,n)}`$ defined by
(1.2)
$$(M,(\lambda ,\mu ;\kappa ))(Z,W)=(MZ,(W+\lambda Z+\mu )(CZ+D)^1),$$
where $`M=\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)Sp(n,),(\lambda ,\mu ;\kappa )H_{}^{(n,m)}`$ and $`(Z,W)_n\times ^{(m,n)}.`$ The homogeneous space $`_n\times ^{(m,n)}`$ is called the Siegel-Jacobi space of degree $`n`$ and index $`m.`$ We refer to \[2-3\], \[6-7\], , \[14-21\] for more details on materials related to the Siegel-Jacobi space.
For brevity, we write $`_{n,m}:=_n\times ^{(m,n)}.`$ For a coordinate $`(Z,W)_{n,m}`$ with $`Z=(z_{\mu \nu })_n`$ and $`W=(w_{kl})^{(m,n)},`$ we put
$`Z=`$ $`X+iY,X=(x_{\mu \nu }),Y=(y_{\mu \nu })\text{real},`$
$`W=`$ $`U+iV,U=(u_{kl}),V=(v_{kl})\text{real},`$
$`dZ=`$ $`(dz_{\mu \nu }),d\overline{Z}=(d\overline{z}_{\mu \nu }),dY=(dy_{\mu \nu }),`$
$`dW=`$ $`(dw_{kl}),d\overline{W}=(d\overline{w}_{kl}),dV=(dv_{kl}),`$
$`{\displaystyle \frac{}{Z}}=`$ $`\left({\displaystyle \frac{1+\delta _{\mu \nu }}{2}}{\displaystyle \frac{}{z_{\mu \nu }}}\right),{\displaystyle \frac{}{\overline{Z}}}=\left({\displaystyle \frac{1+\delta _{\mu \nu }}{2}}{\displaystyle \frac{}{\overline{z}_{\mu \nu }}}\right),`$
$`{\displaystyle \frac{}{X}}=`$ $`\left({\displaystyle \frac{1+\delta _{\mu \nu }}{2}}{\displaystyle \frac{}{x_{\mu \nu }}}\right),{\displaystyle \frac{}{Y}}=\left({\displaystyle \frac{1+\delta _{\mu \nu }}{2}}{\displaystyle \frac{}{y_{\mu \nu }}}\right),`$
$$\frac{}{W}=\left(\begin{array}{ccc}\frac{}{w_{11}}\hfill & \mathrm{}& \frac{}{w_{m1}}\hfill \\ \mathrm{}& \mathrm{}& \mathrm{}\\ \frac{}{w_{1n}}\hfill & \mathrm{}& \frac{}{w_{mn}}\hfill \end{array}\right),\frac{}{\overline{W}}=\left(\begin{array}{ccc}\frac{}{\overline{w}_{11}}\hfill & \mathrm{}& \frac{}{\overline{w}_{m1}}\\ \mathrm{}& \mathrm{}& \mathrm{}\\ \frac{}{\overline{w}_{1n}}& \mathrm{}& \frac{}{\overline{w}_{mn}}\hfill \end{array}\right),$$
$$\frac{}{U}=\left(\begin{array}{ccc}\frac{}{u_{11}}\hfill & \mathrm{}& \frac{}{u_{m1}}\hfill \\ \mathrm{}& \mathrm{}& \mathrm{}\\ \frac{}{u_{1n}}\hfill & \mathrm{}& \frac{}{u_{mn}}\hfill \end{array}\right),\frac{}{V}=\left(\begin{array}{ccc}\frac{}{v_{11}}& \mathrm{}& \frac{}{v_{m1}}\\ \mathrm{}& \mathrm{}& \mathrm{}\\ \frac{}{v_{1n}}& \mathrm{}& \frac{}{v_{mn}}\hfill \end{array}\right),$$
where $`\delta _{ij}`$ denotes the Kronecker delta symbol.
C. L. Siegel introduced the symplectic metric $`ds_n^2`$ on $`_n`$ invariant under the action (1.1) of $`Sp(n,)`$ given by
(1.3) $`ds_n^2=\sigma \left(Y^1dZY^1d\overline{Z}\right)`$
and H. Maass proved that the differential operator
(1.4) $`\mathrm{\Delta }_n=4\sigma \left(Y^t\left(Y{\displaystyle \frac{}{\overline{Z}}}\right){\displaystyle \frac{}{Z}}\right)`$
is the Laplacian of $`_n`$ for the symplectic metric $`ds_n^2.`$ Here $`\sigma (A)`$ denotes the trace of a square matrix $`A`$.
In this paper, for arbitrary positive integers $`n`$ and $`m`$, we express the $`G^J`$-invariant metrics on $`_n\times ^{(m,n)}`$ and their Laplacians explicitly.
In fact, we prove the following theorems.
###### Theorem 1.1.
For any two positive real numbers $`A`$ and $`B`$, the following metric
$`ds_{n,m;A,B}^2`$ $`=`$ $`A\sigma \left(Y^1dZY^1d\overline{Z}\right)`$
$`+B\{\sigma \left(Y^1{}_{}{}^{t}VVY^1dZY^1d\overline{Z}\right)+\sigma \left(Y^1{}_{}{}^{t}(dW)d\overline{W}\right)`$
$`\sigma \left(VY^1dZY^1{}_{}{}^{t}(d\overline{W})\right)\sigma \left(VY^1d\overline{Z}Y^1{}_{}{}^{t}(dW)\right)\}`$
is a Riemannian metric on $`_{n,m}`$ which is invariant under the action (1.2) of the Jacobi group $`G^J`$.
###### Theorem 1.2.
For any two positive real numbers $`A`$ and $`B`$, the Laplacian $`\mathrm{\Delta }_{n,m;A,B}`$ of $`(_{n,m},ds_{n,m;A,B}^2)`$ is given by
$`\mathrm{\Delta }_{n,m;A,B}`$ $`=`$ $`{\displaystyle \frac{4}{A}}\{\sigma \left(Y^t\left(Y{\displaystyle \frac{}{\overline{Z}}}\right){\displaystyle \frac{}{Z}}\right)+\sigma \left(VY^1{}_{}{}^{t}V_{}^{t}\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{W}}\right)`$
$`+\sigma \left(V^t\left(Y{\displaystyle \frac{}{\overline{Z}}}\right){\displaystyle \frac{}{W}}\right)+\sigma \left({}_{}{}^{t}V_{}^{t}\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{Z}}\right)\}`$
$`+{\displaystyle \frac{4}{B}}\sigma \left(Y{\displaystyle \frac{}{W}}^t\left({\displaystyle \frac{}{\overline{W}}}\right)\right).`$
The following differential form
$`dv=\left(\text{det}Y\right)^{(n+m+1)}[dX][dY][dU][dV]`$
is a $`G^J`$-invariant volume element on $`_{n,m}`$, where
$$[dX]=_{\mu \nu }dx_{\mu \nu },[dY]=_{\mu \nu }dy_{\mu \nu },[dU]=_{k,l}du_{kl}\text{and}[dV]=_{k,l}dv_{kl}.$$
The point is that the invariant metric $`ds_{n,m;A,B}^2`$ and its Laplacian $`\mathrm{\Delta }_{n,m;A,B}`$ are expressed in terms of the trace form.
For the case $`n=m=1`$ and $`A=B=1`$, Berndt proved in (cf. ) that the metric $`ds_{1,1}^2`$ on $`\times `$ defined by
$`ds_{1,1}^2:=ds_{1,1;1,1}=`$ $`{\displaystyle \frac{y+v^2}{y^3}}(dx^2+dy^2)+{\displaystyle \frac{1}{y}}(du^2+dv^2)`$
$`{\displaystyle \frac{2v}{y^2}}(dxdu+dydv)`$
is a Riemannian metric on $`\times `$ invariant under the action (1.2) of the Jacobi group and its Laplacian $`\mathrm{\Delta }_{1,1}`$ is given by
$`\mathrm{\Delta }_{1,1}:=\mathrm{\Delta }_{1,1;1,1}=`$ $`y^2\left({\displaystyle \frac{^2}{x^2}}+{\displaystyle \frac{^2}{y^2}}\right)+(y+v^2)\left({\displaystyle \frac{^2}{u^2}}+{\displaystyle \frac{^2}{v^2}}\right)`$
$`+\mathrm{\hspace{0.17em}2}yv\left({\displaystyle \frac{^2}{xu}}+{\displaystyle \frac{^2}{yv}}\right).`$
It is a pleasure to thank Eberhard Freitag for his helpful advice and letting me know the paper of H. Maass. Notations: We denote by $``$ and $``$ the field of real numbers, and the field of complex numbers respectively. The symbol “:=” means that the expression on the right is the definition of that on the left. For two positive integers $`k`$ and $`l`$, $`F^{(k,l)}`$ denotes the set of all $`k\times l`$ matrices with entries in a commutative ring $`F`$. For a square matrix $`AF^{(k,k)}`$ of degree $`k`$, $`\sigma (A)`$ denotes the trace of $`A`$. For any $`MF^{(k,l)},^tM`$ denotes the transpose matrix of $`M`$. $`E_n`$ denotes the identity matrix of degree $`n`$. For $`AF^{(k,l)}`$ and $`BF^{(k,k)}`$, we set $`B[A]=^tABA.`$ For a complex matrix $`A`$, $`\overline{A}`$ denotes the complex conjugate of $`A`$. For $`A^{(k,l)}`$ and $`B^{(k,k)}`$, we use the abbreviation $`B\{A\}=^t\overline{A}BA.`$
## 2. Proof of Theorem 1.1
Let $`g=(M,(\lambda ,\mu ;\kappa ))`$ be an element of $`G^J`$ with $`M=\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)Sp(n,)`$ and $`(Z,W)_{n,m}`$ with $`Z_n`$ and $`W^{(m,n)}.`$ If we put $`(Z_{},W_{}):=g(Z,W),`$ then we have
$`Z_{}=MZ=(AZ+B)(CZ+D)^1,`$
$`W_{}=(W+\lambda Z+\mu )(CZ+D)^1.`$
Thus we obtain
(2.1) $`dZ_{}=dZ[(CZ+D)^1]={}_{}{}^{t}(CZ+D)_{}^{1}dZ(CZ+D)^1`$
and
(2.2) $`dW_{}=dW(CZ+D)^1+\{\lambda (W+\lambda Z+\mu )(CZ+D)^1C\}dZ(CZ+D)^1.`$
Here we used the following facts that
$$d(CZ+D)^1=(CZ+D)^1CdZ(CZ+D)^1$$
and that $`(CZ+D)^1C`$ is symmetric.
We put
$$Z_{}=X_{}+iY_{},W_{}=U_{}+iV_{},X_{},Y_{},U_{},V_{}\text{real}.$$
From , p.33 or , p.128, we know that
(2.3) $`Y_{}=Y\{(CZ+D)^1\}={}_{}{}^{t}(C\overline{Z}+D)_{}^{1}Y(CZ+D)^1.`$
First of all, we recall that the following matrices
$`t(b)=\left(\begin{array}{cc}E_n& b\\ 0& E_n\end{array}\right),b={}_{}{}^{t}b\text{real},`$
$`g_0(h)=\left(\begin{array}{cc}{}_{}{}^{t}h& 0\\ 0& h^1\end{array}\right),hGL(n,),`$
$`J_n=\left(\begin{array}{cc}0& E_n\\ E_n& 0\end{array}\right)`$
generate the symplectic group $`Sp(n,)`$ (cf. , ). Therefore the following elements
$`t(b;\lambda ,\mu ,\kappa ),g(h)`$ and $`\sigma _n`$ of $`G^J`$ defined by
$`t(b;\lambda ,\mu ,\kappa )=(\left(\begin{array}{cc}E_n& b\\ 0& E_n\end{array}\right),(\lambda ,\mu ;\kappa )),b={}_{}{}^{t}b\text{real},(\lambda ,\mu ;\kappa )H_{}^{(n,m)},`$
$`g(h)=(\left(\begin{array}{cc}{}_{}{}^{t}h& 0\\ 0& h^1\end{array}\right),(0,0;0)),hGL(n,),`$
$`\sigma _n=(\left(\begin{array}{cc}0& E_n\\ E_n& 0\end{array}\right),(0,0;0))`$
generate the Jacobi group $`G^J.`$ So it suffices to prove the invariance of the metric $`ds_{n,m;A,B}^2`$ under the action of the generators $`t(b;\lambda ,\mu ,\kappa ),g(h)`$ and $`\sigma _n.`$ For brevity, we write
$`(a)=\sigma \left(Y^1dZY^1d\overline{Z}\right),`$
$`(b)=\sigma \left(Y^1{}_{}{}^{t}VVY^1dZY^1d\overline{Z}\right),`$
$`(c)=\sigma \left(Y^1{}_{}{}^{t}(dW)d\overline{W}\right),`$
$`(d)=\sigma (VY^1dZY^1{}_{}{}^{t}(d\overline{W})+VY^1d\overline{Z}Y^1{}_{}{}^{t}(dW))`$
and
$`(a)_{}=\sigma \left(Y_{}^1dZ_{}Y_{}^1d\overline{Z}_{}\right),`$
$`(b)_{}=\sigma \left(Y_{}^1{}_{}{}^{t}V_{}^{}V_{}Y_{}^1dZ_{}Y_{}^1d\overline{Z}_{}\right),`$
$`(c)_{}=\sigma \left(Y_{}^1{}_{}{}^{t}(dW_{})d\overline{W}_{}\right),`$
$`(d)_{}=\sigma (V_{}Y_{}^1dZ_{}Y_{}^1{}_{}{}^{t}(d\overline{W}_{})+V_{}Y_{}^1d\overline{Z}_{}Y_{}^1{}_{}{}^{t}(dW_{}))`$
Case I. $`g=t(b;\lambda ,\mu ,\kappa )`$ with $`b={}_{}{}^{t}b`$ real and $`(\lambda ,\mu ;\kappa )H_{}^{(n,m)}.`$ In this case, we have
$$Z_{}=Z+b,Y_{}=Y,W_{}=W+\lambda Z+\mu ,V_{}=V+\lambda Y$$
and
$$dZ_{}=dZ,dW_{}=dW+\lambda dZ.$$
Therefore
$`(a)_{}`$ $`=`$ $`\sigma \left(Y_{}^1dZ_{}Y_{}^1d\overline{Z_{}}\right)=\sigma \left(Y^1dZY^1d\overline{Z}\right)=(a),`$
$`(b)_{}`$ $`=`$ $`\sigma \left(Y^1{}_{}{}^{t}VVY^1dZY^1d\overline{Z}\right)+\sigma \left(Y^1{}_{}{}^{t}V\lambda dZY^1d\overline{Z}\right)`$
$`+\sigma \left({}_{}{}^{t}\lambda VY^1dZY^1d\overline{Z}\right)+\sigma \left({}_{}{}^{t}\lambda \lambda dZY^1d\overline{Z}\right),`$
$`(c)_{}`$ $`=`$ $`\sigma \left(Y^1{}_{}{}^{t}(dW)d\overline{W}\right)+\sigma \left(Y^1{}_{}{}^{t}(dW)\lambda d\overline{Z}\right)`$
$`+\sigma \left(Y^1dZ{}_{}{}^{t}\lambda d\overline{W}\right)+\sigma \left(Y^1dZ{}_{}{}^{t}\lambda \lambda d\overline{Z}\right)`$
and
$`(d)_{}`$ $`=`$ $`\sigma \left(VY^1dZY^1{}_{}{}^{t}(d\overline{W})\right)\sigma \left(\lambda dZY^1{}_{}{}^{t}(d\overline{W})\right)`$
$`\sigma \left(VY^1dZY^1d\overline{Z}{}_{}{}^{t}\lambda \right)\sigma \left(\lambda dZY^1d\overline{Z}{}_{}{}^{t}\lambda \right)`$
$`\sigma \left(VY^1d\overline{Z}Y^1{}_{}{}^{t}(dW)\right)\sigma \left(\lambda d\overline{Z}Y^1{}_{}{}^{t}(dW)\right)`$
$`\sigma \left(VY^1d\overline{Z}Y^1dZ{}_{}{}^{t}\lambda \right)\sigma \left(\lambda d\overline{Z}Y^1dZ{}_{}{}^{t}\lambda \right).`$
Thus we see that
$`(a)=(a)_{}\text{and}(b)+(c)+(d)=(b)_{}+(c)_{}+(d)_{}.`$
Hence
$`ds_{n,m;A,B}^2=A(a)+B\left\{(b)+(c)+(d)\right\}`$
is invariant under the action of $`t(B;\lambda ,\mu ,\kappa ).`$
Case II. $`g=g(h)`$ with $`hGL(n,).`$ In this case, we have
$$Z_{}={}_{}{}^{t}hZh,Y_{}={}_{}{}^{t}hYh,W_{}=Wh,V_{}=Vh$$
and
$$dZ_{}={}_{}{}^{t}hdZh,dW_{}=dWh.$$
Therefore by an easy computation, we see that each of $`(a),(b),(c)`$ and $`(d)`$ is invariant under the action of all $`g(h)`$ with $`hGL(n,).`$ Hence the metric $`ds_{n,m;A,B}^2`$ is invariant under the action of all $`g(h)`$ with $`hGL(n,).`$
Case III. $`g=\sigma _n=(\left(\begin{array}{cc}0& E_n\\ E_n& 0\end{array}\right),(0,0;0)).`$
In this case, we have
(2.4)
$$Z_{}=Z^1\text{and}W_{}=WZ^1.$$
We set
$$\theta _1:=\text{Re}Z^1\text{and}\theta _2:=\text{Im}Z^1.$$
Then $`\theta _1`$ and $`\theta _2`$ are symmetric matrices and we have
(2.5)
$$Y_{}=\theta _2\text{and}V_{}:=\text{Im}W_{}=V\theta _1+U\theta _2.$$
It is easy to see that
(2.6)
$$Y=Z\theta _2\overline{Z}=\overline{Z}\theta _2Z,$$
(2.7)
$$\theta _1Y+\theta _2X=0$$
and
(2.8)
$$\theta _1X\theta _2Y=E_n.$$
According to (2.6) and (2.7), we obtain
(2.9)
$$X=(\theta _2)^1\theta _1Y\text{and}Y^1=\theta _1(\theta _2)^1\theta _1\theta _2.$$
From (2.1) and (2.2), we have
(2.10)
$$dZ_{}=Z^1dZZ^1$$
and
(2.11)
$$dW_{}=dWZ^1WZ^1dZZ^1=\left(dWWZ^1dZ\right)Z^1.$$
Therefore we have, according to (2.6) and (2.10),
$`(a)_{}`$ $`=`$ $`\sigma \left((\theta _2)^1Z^1dZZ^1(\theta _2)^1\overline{Z}^1d\overline{Z}\overline{Z}^1\right)`$
$`=`$ $`\sigma \left(Y^1dZY^1d\overline{Z}\right)=(a).`$
According to (2.5)-(2.10), we have
$`(b)_{}`$ $`=`$ $`\sigma \left((\theta _2)^1(\theta _1{}_{}{}^{t}V+\theta _2{}_{}{}^{t}U)(V\theta _1+U\theta _2)(\theta _2)^1Z^1dZY^1d\overline{Z}\overline{Z}^1\right)`$
$`=`$ $`\sigma \left(\{{}_{}{}^{t}U(\theta _2)^1\theta _1{}_{}{}^{t}V\}\{UV\theta _1(\theta _2)^1\}Z^1dZY^1d\overline{Z}\overline{Z}^1\right)`$
$`=`$ $`\sigma (\{{}_{}{}^{t}\overline{W}+(iE_n(\theta _2)^1\theta _1){}_{}{}^{t}V\}\{WV(iE_n+\theta _1(\theta _2)^1)\}`$
$`Z^1dZY^1d\overline{Z}\overline{Z}^1),`$
$`(c)_{}`$ $`=`$ $`\sigma \left((\theta _2)^1(Z^1{}_{}{}^{t}(dW)Z^1dZZ^1{}_{}{}^{t}W)(d\overline{W}\overline{Z}^1\overline{W}\overline{Z}^1d\overline{Z}\overline{Z}^1)\right)`$
$`=`$ $`\sigma ((\theta _2)^1Z^1{}_{}{}^{t}(dW)d\overline{W}\overline{Z}^1(\theta _2)^1Z^1{}_{}{}^{t}(dW)\overline{W}\overline{Z}^1d\overline{Z}\overline{Z}^1`$
$`(\theta _2)^1Z^1dZZ^1{}_{}{}^{t}Wd\overline{W}\overline{Z}^1`$
$`+(\theta _2)^1Z^1dZZ^1{}_{}{}^{t}W\overline{W}\overline{Z}^1d\overline{Z}\overline{Z}^1)`$
$`=`$ $`\sigma \left(Y^1{}_{}{}^{t}(dW)d\overline{W}\right)\sigma \left(Y^1{}_{}{}^{t}(dW)\overline{W}\overline{Z}^1d\overline{Z}\right)`$
$`\sigma \left(Y^1dZZ^1{}_{}{}^{t}Wd\overline{W}\right)+\sigma \left(Y^1dZZ^1{}_{}{}^{t}W\overline{W}\overline{Z}^1d\overline{Z}\right)`$
and
$`(d)_{}`$ $`=`$ $`\sigma \left((V\theta _1+U\theta _2)(\theta _2)^1Z^1dZZ^1(\theta _2)^1\{\overline{Z}^1{}_{}{}^{t}(d\overline{W})\overline{Z}^1d\overline{Z}\overline{Z}^1{}_{}{}^{t}\overline{W}\}\right)`$
$`\sigma \left((V\theta _1+U\theta _2)(\theta _2)^1\overline{Z}^1d\overline{Z}\overline{Z}^1(\theta _2)^1\{Z^1{}_{}{}^{t}(dW)Z^1dZZ^1{}_{}{}^{t}W\}\right)`$
$`=`$ $`\sigma \left((V\theta _1+U\theta _2)(\theta _2)^1Z^1dZY^1{}_{}{}^{t}(d\overline{W})\right)`$
$`+\sigma \left((V\theta _1+U\theta _2)(\theta _2)^1Z^1dZY^1d\overline{Z}\overline{Z}^1{}_{}{}^{t}\overline{W}\right)`$
$`\sigma \left((V\theta _1+U\theta _2)(\theta _2)^1\overline{Z}^1d\overline{Z}Y^1{}_{}{}^{t}(dW)\right)`$
$`+\sigma \left((V\theta _1+U\theta _2)(\theta _2)^1\overline{Z}^1d\overline{Z}Y^1dZZ^1{}_{}{}^{t}W\right).`$
Taking the $`(dZ,d\overline{W})`$-part $`\mathrm{}(Z,\overline{W})`$ in $`(b)_{}+(c)_{}+(d)_{},`$ we have
$`\mathrm{}(Z,\overline{W})`$ $`=`$ $`\sigma \left(VY^1dZY^1{}_{}{}^{t}(d\overline{W})\right)+\sigma \left(Y^1dZ({}_{}{}^{t}W_{}^{}Z^1{}_{}{}^{t}W)d\overline{W}\right)`$
$`=`$ $`\sigma \left(VY^1dZY^1{}_{}{}^{t}(d\overline{W})\right)\text{because}W_{}=WZ^1(\text{cf.}(2.4)).`$
Similiarly, if we take the $`(d\overline{Z},dW)`$-part $`\mathrm{}(\overline{Z},W)`$ in $`(b)_{}+(c)_{}+(d)_{},`$ we have
$`\mathrm{}(\overline{Z},W)`$ $`=`$ $`\sigma \left(VY^1d\overline{Z}Y^1{}_{}{}^{t}(dW)\right)+\sigma \left(d\overline{Z}Y^1{}_{}{}^{t}(dW)(\overline{W_{}}\overline{W}\overline{Z}^1)\right)`$
$`=`$ $`\sigma \left(VY^1d\overline{Z}Y^1{}_{}{}^{t}(dW)\right)\text{because}W_{}=WZ^1.`$
If we take the $`(dW,d\overline{W})`$-part $`\mathrm{}(W,\overline{W})`$ in $`(b)_{}+(c)_{}+(d)_{},`$ we have
$$\mathrm{}(W,\overline{W})=\sigma \left(Y^1{}_{}{}^{t}(dW)d\overline{W}\right).$$
Finally, if we take the $`(dZ,d\overline{Z})`$-part $`\mathrm{}(Z,\overline{Z})`$ in $`(b)_{}+(c)_{}+(d)_{},`$ we have
$`\mathrm{}(Z,\overline{Z})`$ $`=`$ $`\sigma (\{{}_{}{}^{t}\overline{W}+(iE_n(\theta _2)^1\theta _1){}_{}{}^{t}V\}\{WV(iE_n+\theta _1(\theta _2)^1)\}`$
$`Z^1dZY^1d\overline{Z}\overline{Z}^1)`$
$`+\sigma \left(Z^1{}_{}{}^{t}W\overline{W}\overline{Z}^1d\overline{Z}Y^1dZ\right)`$
$`+\sigma \left({}_{}{}^{t}\overline{W}(V\theta _1+U\theta _2)(\theta _2)^1Z^1dZY^1d\overline{Z}\overline{Z}^1\right)`$
$`+\sigma \left({}_{}{}^{t}W(V\theta _1+U\theta _2)(\theta _2)^1\overline{Z}^1d\overline{Z}Y^1dZZ^1\right).`$
Since
$`(V\theta _1+U\theta _2)(\theta _2)^1`$ $`=`$ $`U+V\theta _1(\theta _2)^1`$
$`=`$ $`W+V\{iE_n+\theta _1(\theta _2)^1\}`$
$`=`$ $`\overline{W}V\{iE_n\theta _1(\theta _2)^1\},`$
we have
$`\mathrm{}(Z,\overline{Z})`$ $`=`$ $`\sigma (\overline{Z}^1\{iE_n(\theta _2)^1\theta _1\}{}_{}{}^{t}V\{WV(iE_n+\theta _1(\theta _2)^1)\}`$
$`Z^1dZY^1d\overline{Z})`$
$`\sigma \left(\overline{Z}^1\left\{iE_n(\theta _2)^1\theta _1\right\}{}_{}{}^{t}VWZ^1dZY^1d\overline{Z}\right)`$
$`=`$ $`\sigma (\overline{Z}^1\{iE_n(\theta _2)^1\theta _1\}{}_{}{}^{t}VV\{iE_n+\theta _1(\theta _2)^1\}`$
$`Z^1dZY^1d\overline{Z}).`$
By the way, according to (2.9), we obtain
$`\overline{Z}^1\left\{iE_n(\theta _2)^1\theta _1\right\}`$ $`=`$ $`(\theta _1i\theta _2)\left\{iE_n(\theta _2)^1\theta _1\right\}`$
$`=`$ $`\theta _2\theta _1(\theta _2)^1\theta _1=Y^1`$
and
$`\left\{iE_n+\theta _1(\theta _2)^1\right\}Z^1=\theta _1(\theta _2)^1\theta _1\theta _2=Y^1.`$
Therefore
$$\mathrm{}(Z,\overline{Z})=\sigma \left(Y^1{}_{}{}^{t}VVY^1dZY^1d\overline{Z}\right).$$
Hence $`(a)=(a)_{}`$ and
$`(b)_{}+(c)_{}+(d)_{}`$ $`=`$ $`\mathrm{}(Z,\overline{W})+\mathrm{}(\overline{Z},W)+\mathrm{}(W,\overline{W})+\mathrm{}(Z,\overline{Z})`$
$`=`$ $`(b)+(c)+(d).`$
This implies that the metric
$`ds_{n,m;A,B}^2=A(a)+B\left\{(b)+(c)+(d)\right\}`$
is invariant under the action (1.2) of $`\sigma _n.`$
Consequently $`ds_{n,m;A,B}^2`$ is invariant under the action (1.2) of the Jacobi group $`G^J.`$ In particular, for $`(Z,W)=(iE_n,0),`$ we have
$`ds_{n,m;A,B}^2`$ $`=`$ $`A\sigma \left(dZd\overline{Z}\right)+B\sigma \left({}_{}{}^{t}(dW)d\overline{W}\right)`$
$`=`$ $`A\left\{{\displaystyle \underset{\mu =1}{\overset{n}{}}}(dx_{\mu \mu }^2+dy_{\mu \mu }^2)+2{\displaystyle \underset{1\mu <\nu n}{}}(dx_{\mu \nu }^2+dy_{\mu \nu }^2)\right\}`$
$`+B\left\{{\displaystyle \underset{1km,\mathrm{\hspace{0.17em}1}ln}{}}(du_{kl}^2+dv_{kl}^2)\right\},`$
which is clearly positive definite. Since $`G^J`$ acts on $`_{n,m}`$ transitively, $`ds_{n,m;A,B}^2`$ is positive definite everywhere in $`_{n,m}.`$ This completes the proof of Theorem 1.1. $`\mathrm{}`$
Remark 2.1. The scalar curvature of the Siegel-Jacobi space $`(_{n,m},ds_{n,m;A,B}^2)`$ is constant because of the transitive group action of $`G^J`$ on $`_{n,m}.`$ In the special case $`n=m=1`$ and $`A=B=1`$, by a direct computation, we see that the scalar curvature of $`(_{1,1},ds_{1,1;1,1}^2)`$ is $`3.`$
## 3. Proof of Theorem 1.2
If $`(Z_{},W_{})=g(Z,W)`$ with $`g=(\left(\begin{array}{cc}A& B\\ C& D\end{array}\right),(\lambda ,\mu ;\kappa ))G^J,`$ we can see easily that
$`{\displaystyle \frac{}{Z_{}}}`$ $`=`$ $`(CZ+D){}_{}{}^{t}\left\{(CZ+D){\displaystyle \frac{}{Z}}\right\}`$
$`+(CZ+D){}_{}{}^{t}\left\{(C{}_{}{}^{t}W+C{}_{}{}^{t}\mu D{}_{}{}^{t}\lambda ){}_{}{}^{t}\left({\displaystyle \frac{}{W}}\right)\right\}`$
and
(3.2)
$$\frac{}{W_{}}=(CZ+D)\frac{}{W}.$$
For brevity, we put
$`(\alpha ):=4\sigma \left(Y^t\left(Y{\displaystyle \frac{}{\overline{Z}}}\right){\displaystyle \frac{}{Z}}\right),`$
$`(\beta ):=\mathrm{\hspace{0.33em}4}\sigma \left(Y{\displaystyle \frac{}{W}}{}_{}{}^{t}\left({\displaystyle \frac{}{\overline{W}}}\right)\right),`$
$`(\gamma ):=\mathrm{\hspace{0.33em}4}\sigma \left(VY^1{}_{}{}^{t}V_{}^{t}\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{W}}\right),`$
$`(\delta ):=\mathrm{\hspace{0.33em}4}\sigma \left(V^t\left(Y{\displaystyle \frac{}{\overline{Z}}}\right){\displaystyle \frac{}{W}}\right)`$
and
$`(ϵ):=\mathrm{\hspace{0.33em}4}\sigma \left({}_{}{}^{t}V_{}^{t}\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{Z}}\right).`$
We also set
$`(\alpha )_{}:=4\sigma \left(Y_{}^t\left(Y_{}{\displaystyle \frac{}{\overline{Z}}}_{}\right){\displaystyle \frac{}{Z}}_{}\right),`$
$`(\beta )_{}:=\mathrm{\hspace{0.33em}4}\sigma \left(Y_{}{\displaystyle \frac{}{W}}_{}{}_{}{}^{t}\left({\displaystyle \frac{}{\overline{W}}}_{}\right)\right),`$
$`(\gamma )_{}:=\mathrm{\hspace{0.33em}4}\sigma \left(V_{}Y_{}^1{}_{}{}^{t}V_{}^{t}\left(Y_{}{\displaystyle \frac{}{\overline{W}}}_{}\right){\displaystyle \frac{}{W}}_{}\right),`$
$`(\delta )_{}:=\mathrm{\hspace{0.33em}4}\sigma \left(V_{}^t\left(Y_{}{\displaystyle \frac{}{\overline{Z}}}_{}\right){\displaystyle \frac{}{W}}_{}\right)`$
and
$`(ϵ)_{}:=\mathrm{\hspace{0.33em}4}\sigma \left({}_{}{}^{t}V_{}^{t}\left(Y_{}{\displaystyle \frac{}{\overline{W}}}_{}\right){\displaystyle \frac{}{Z}}_{}\right).`$
We need the following lemma for the proof of Theorem 1.2. H. Maass observed the following useful fact.
###### Lemma 3.1.
(a) Let $`A`$ be an $`n\times k`$ matrix and let $`B`$ be a $`k\times n`$ matrix. Assume that the entries of $`A`$ commute with the entries of $`B`$. Then $`\sigma (AB)=\sigma (BA).`$ (b) Let $`A`$ be an $`m\times n`$ matrix and $`B`$ an $`n\times l`$ matrix. Assume that the entries of $`A`$ commute with the entries of $`B`$. Then $`{}_{}{}^{t}(AB)={}_{}{}^{t}B{}_{}{}^{t}A.`$ (c) Let $`A,B`$ and $`C`$ be a $`k\times l`$, an $`n\times m`$ and an $`m\times l`$ matrix respectively. Assume that the entries of $`A`$ commute with the entries of $`B`$. Then
$`{}_{}{}^{t}(A{}_{}{}^{t}(BC))=B{}_{}{}^{t}(A^tC).`$
Proof. The proof follows immediately from the direct computation. $`\mathrm{}`$ Now we are ready to prove Theorem 1.2. First of all, we shall prove that $`\mathrm{\Delta }_{n,m;A,B}`$ is invariant under the action of the generators $`t(b;\lambda ,\mu ,\kappa ),g(h)`$ and $`\sigma _n.`$ Case I. $`g=t(b;\lambda ,\mu ,\kappa )=(\left(\begin{array}{cc}E_n& b\\ 0& E_n\end{array}\right),(\lambda ,\mu ;\kappa ))`$ with $`b=^tb`$ real.
In this case, we have
$$Y_{}=Y,V_{}=V+\lambda Y$$
and
$$\frac{}{Z_{}}=\frac{}{Z}^t\left({}_{}{}^{t}\lambda _{}^{t}\left(\frac{}{W}\right)\right)\text{and}\frac{}{W_{}}=\frac{}{W}.$$
Using Lemma 3.1, we obtain
$`(\alpha )_{}`$ $`=`$ $`(\alpha )\sigma \left(\lambda Y^t\left(Y{\displaystyle \frac{}{\overline{Z}}}\right){\displaystyle \frac{}{W}}\right)`$
$`\sigma \left(Y^t\lambda ^t\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{Z}}\right)+\sigma \left(\lambda Y^t\lambda ^t\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{W}}\right)`$
$`(\beta )_{}`$ $`=`$ $`\sigma \left(Y_{}{\displaystyle \frac{}{W_{}}}^t\left({\displaystyle \frac{}{\overline{W_{}}}}\right)\right)=(\beta ),`$
$`(\gamma )_{}`$ $`=`$ $`(\gamma )+\sigma \left(\lambda ^tV^t\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{W}}\right)`$
$`+\sigma \left(V^t\lambda ^t\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{W}}\right)+\sigma \left(\lambda Y^t\lambda ^t\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{W}}\right),`$
$`(\delta )_{}`$ $`=`$ $`(\delta )+\sigma \left(\lambda Y^t\left(Y{\displaystyle \frac{}{\overline{Z}}}\right){\displaystyle \frac{}{W}}\right)`$
$`\sigma \left(V^t\lambda ^t\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{W}}\right)+\sigma \left(\lambda Y^t\lambda ^t\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{W}}\right)`$
and
$`(ϵ)_{}`$ $`=`$ $`(ϵ)+\sigma \left(Y^t\lambda ^t\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{Z}}\right)`$
$`\sigma \left(\lambda ^tV^t\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{W}}\right)\sigma \left(\lambda Y^t\lambda ^t\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{W}}\right).`$
Thus $`(\beta )=(\beta )_{}`$ and
$`(\alpha )+(\gamma )+(\delta )+(ϵ)=(\alpha )_{}+(\gamma )_{}+(\delta )_{}+(ϵ)_{}.`$
Hence
$`\mathrm{\Delta }_{n,m;A,B}={\displaystyle \frac{4}{B}}(\beta )+{\displaystyle \frac{4}{A}}\left\{(\alpha )+(\gamma )+(\delta )+(ϵ)\right\}`$
is invariant under the action of all $`t(b;\lambda ,\mu ,\kappa ).`$
Case II. $`g=g(h)=(\left(\begin{array}{cc}{}_{}{}^{t}h& 0\\ 0& h^1\end{array}\right),(0,0;0))`$ with $`hGL(n,)`$.
In this case, we have
$$Y_{}=^thYhV_{}=Vh$$
and
$$\frac{}{Z_{}}=h^1{}_{}{}^{t}\left(h^1\frac{}{Z}\right),\frac{}{W_{}}=h^1\frac{}{W}.$$
According to Lemma 3.1, we see that each of $`(\alpha ),(\beta ),(\gamma ),(\delta )`$ and $`(ϵ)`$ is invariant under the action of all $`g(h)`$ with $`hGL(n,).`$ Therefore $`\mathrm{\Delta }_{n,m;A,B}`$ is invariant under the action of all $`g(h)`$ with $`hGL(n,).`$
Case III. $`g=\sigma _n=(\left(\begin{array}{cc}0& E_n\\ E_n& 0\end{array}\right),(0,0;0)).`$
In this case, we have
$$Z_{}=Z^1\text{and}W_{}=WZ^1.$$
We set
$$\theta _1:=\text{Re}Z^1\text{and}\theta _2:=\text{Im}Z^1.$$
Then we obtain the relations (2.5)-(2.9). From (2.6), we have the relation
(3.3)
$$\theta _2\overline{Z}=Z^1Y.$$
It follows from the relation (2.3) that
(3.4)
$$Y_{}=\overline{Z}^1YZ^1=Z^1Y\overline{Z}^1=\theta _2.$$
From (2.9), we obtain
(3.5)
$$\theta _1\theta _2^1\theta _1=Y^1\theta _2.$$
According to (3.1) and (3.2), we have
(3.6)
$$\frac{}{Z_{}}=Z^t\left(Z\frac{}{Z}\right)+Z^t\left({}_{}{}^{t}W_{}^{t}\left(\frac{}{W}\right)\right)$$
and
(3.7)
$$\frac{}{W_{}}=Z\frac{}{W}.$$
From (2.6), (3.3) and Lemma 3.1, we obtain
$`(\alpha )_{}`$ $`=`$ $`(\alpha )\sigma \left(U\theta _2\overline{Z}^t\left(Y{\displaystyle \frac{}{\overline{Z}}}\right){\displaystyle \frac{}{W}}\right)i\sigma \left(V\theta _2\overline{Z}^t\left(Y{\displaystyle \frac{}{\overline{Z}}}\right){\displaystyle \frac{}{W}}\right)`$
$`\sigma \left(Z\theta _2^tU^t\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{Z}}\right)+i\sigma \left(Z\theta _2^tV^t\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{Z}}\right)`$
$`\sigma \left(W\theta _2^tU^t\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{W}}\right)+\sigma \left(W\theta _2^tV^t\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{W}}\right).`$
From the relation (3.4), we see $`(\beta )_{}=(\beta )`$. According to (3.3), (3.5) and Lemma 3.1, we otain
$`(\gamma )_{}=(\gamma )+\sigma \left((V\theta _2^tVV\theta _1^tUU\theta _1^tVU\theta _2^tU)^t\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{W}}\right).`$
Using the relation (3.3) and Lemma 3.1, we finally obtain
$`(\delta )_{}`$ $`=`$ $`\sigma \left(V\theta _1\overline{Z}^t\left(Y{\displaystyle \frac{}{\overline{Z}}}\right){\displaystyle \frac{}{W}}\right)+\sigma \left(U\theta _2\overline{Z}^t\left(Y{\displaystyle \frac{}{\overline{Z}}}\right){\displaystyle \frac{}{W}}\right)`$
$`+\sigma \left(V\theta _1^t\overline{W}^t\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{W}}\right)+\sigma \left(U\theta _2^t\overline{W}^t\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{W}}\right)`$
and
$`(ϵ)_{}`$ $`=`$ $`\sigma \left(Z\theta _1^tV^t\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{Z}}\right)+\sigma \left(Z\theta _2^tU^t\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{Z}}\right)`$
$`+\sigma \left(W\theta _1^tV^t\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{W}}\right)+\sigma \left(W\theta _2^tU^t\left(Y{\displaystyle \frac{}{\overline{W}}}\right){\displaystyle \frac{}{W}}\right).`$
Using the fact $`Z^1=\theta _1+i\theta _2`$, we can show that
$$(\alpha )+(\gamma )+(\delta )+(ϵ)=(\alpha )_{}+(\gamma )_{}+(\delta )_{}+(ϵ)_{}.$$
Hence
$`\mathrm{\Delta }_{n,m;A,B}={\displaystyle \frac{4}{B}}(\beta )+{\displaystyle \frac{4}{A}}\left\{(\alpha )+(\gamma )+(\delta )+(ϵ)\right\}`$
is invariant under the action of $`\sigma _n`$.
Consequently $`\mathrm{\Delta }_{n,m;A,B}`$ is invariant under the action (1.2) of $`G^J`$. In particular, for $`(Z,W)=(iE_n,0),`$ the differential operator $`\mathrm{\Delta }_{n,m;A,B}`$ coincides with the Laplacian for the metric $`ds_{n,m;A,B}^2`$. It follows from the invariance of $`\mathrm{\Delta }_{n,m;A,B}`$ under the action (1.2) and the transitivity of the action of $`G^J`$ on $`_{n,m}`$ that $`\mathrm{\Delta }_{n,m;A,B}`$ is the Laplacian of $`(_{n,m},ds_{n,m;A,B}^2).`$ The invariance of the differential form $`dv`$ follows from the fact that the following differential form
$$(detY)^{(n+1)}[dX][dY]$$
is invariant under the action (1.1) of $`Sp(n,)`$ (cf. , p. 130). $`\mathrm{}`$
## 4. Remark on Spectral Theory of $`\mathrm{\Delta }_{n,m;A,B}`$ on Siegel-Jacobi Space
Before we describe a fundamental domain for the Siegel-Jacobi space, we review the Siegel’s fundamental domain for the Siegel upper half plane. We let
$$𝒫_n=\left\{Y^{(n,n)}\right|Y=^tY>0\}$$
be an open cone in $`^{n(n+1)/2}`$. The general linear group $`GL(n,)`$ acts on $`𝒫_n`$ transitively by
$$hY:=hY{}_{}{}^{t}h,hGL(n,),Y𝒫_n.$$
Thus $`𝒫_n`$ is a symmetric space diffeomorphic to $`GL(n,)/O(n).`$ We let
$`GL(n,)=\left\{hGL(n,)\right|h\text{is integral}\}`$
be the discrete subgroup of $`GL(n,)`$.
The fundamental domain $`_n`$ for $`GL(n,)\backslash 𝒫_n`$ which was found by H. Minkowski is defined as a subset of $`𝒫_n`$ consisting of $`Y=(y_{ij})𝒫_n`$ satisfying the following conditions (M.1)-(M.2) (cf. p. 123): (M.1) $`aY^tay_{kk}`$ for every $`a=(a_i)^n`$ in which $`a_k,\mathrm{},a_n`$ are relatively prime for
$`k=1,2,\mathrm{},n`$. (M.2) $`y_{k,k+1}0`$ for $`k=1,\mathrm{},n1.`$ We say that a point of $`_n`$ is Minkowski reduced or simply M-reduced.
Siegel determined a fundamental domain $`_n`$ for $`\mathrm{\Gamma }_n\backslash _n,`$ where $`\mathrm{\Gamma }_n=Sp(n,)`$ is the Siegel modular group of degree $`n`$. We say that $`\mathrm{\Omega }=X+iY_n`$ with $`X,Y`$ real is Siegel reduced or S-reduced if it has the following three properties: (S.1) $`det(\text{Im}(\gamma \mathrm{\Omega }))det(\text{Im}(\mathrm{\Omega }))\text{for all}\gamma \mathrm{\Gamma }_n`$; (S.2) $`Y=\text{Im}\mathrm{\Omega }`$ is M-reduced, that is, $`Y_n;`$ (S.3) $`|x_{ij}|\frac{1}{2}\text{for}1i,jn,\text{where}X=(x_{ij}).`$ $`_n`$ is defined as the set of all Siegel reduced points in $`_n.`$ Using the highest point method, Siegel proved the following (F1)-(F3) (cf. , p. 169): (F1) $`\mathrm{\Gamma }_n_n=_n,`$ i.e., $`_n=_{\gamma \mathrm{\Gamma }_n}\gamma _n.`$ (F2) $`_n`$ is closed in $`_n.`$ (F3) $`_n`$ is connected and the boundary of $`_n`$ consists of a finite number of hyperplanes.
The metric $`ds_n^2`$ given by (1.3) induces a metric $`ds__n^2`$ on $`_n`$. Siegel computed the volume of $`_n`$
$$\text{vol}(_n)=2\underset{k=1}{\overset{n}{}}\pi ^k\mathrm{\Gamma }(k)\zeta (2k),$$
where $`\mathrm{\Gamma }(s)`$ denotes the Gamma function and $`\zeta (s)`$ denotes the Riemann zeta function. For instance,
$$\text{vol}(_1)=\frac{\pi }{3},\text{vol}(_2)=\frac{\pi ^3}{270},\text{vol}(_3)=\frac{\pi ^6}{127575},\text{vol}(_4)=\frac{\pi ^{10}}{200930625}.$$
Let $`f_{kl}(1km,1ln)`$ be the $`m\times n`$ matrix with entry $`1`$ where the $`k`$-th row and the $`l`$-th column meet, and all other entries $`0`$. For an element $`\mathrm{\Omega }_n`$, we set for brevity
$$h_{kl}(\mathrm{\Omega }):=f_{kl}\mathrm{\Omega },1km,1ln.$$
For each $`\mathrm{\Omega }_n,`$ we define a subset $`P_\mathrm{\Omega }`$ of $`^{(m,n)}`$ by
$$P_\mathrm{\Omega }=\left\{\underset{k=1}{\overset{m}{}}\underset{j=1}{\overset{n}{}}\lambda _{kl}f_{kl}+\underset{k=1}{\overset{m}{}}\underset{j=1}{\overset{n}{}}\mu _{kl}h_{kl}(\mathrm{\Omega })\right|0\lambda _{kl},\mu _{kl}1\}.$$
For each $`\mathrm{\Omega }_n,`$ we define the subset $`D_\mathrm{\Omega }`$ of $`_n\times ^{(m,n)}`$ by
$$D_\mathrm{\Omega }:=\left\{(\mathrm{\Omega },Z)_n\times ^{(m,n)}\right|ZP_\mathrm{\Omega }\}.$$
We define
$$_{n,m}:=_{\mathrm{\Omega }_n}D_\mathrm{\Omega }.$$
###### Theorem 4.1.
Let
$$\mathrm{\Gamma }_{n,m}:=Sp(n,)H_{}^{(n,m)}$$
be the discrete subgroup of $`G^J`$, where
$$H_{}^{(n,m)}=\left\{(\lambda ,\mu ;\kappa )H_{}^{(n,m)}\right|\lambda ,\mu ,\kappa \text{are integral}\}.$$
Then $`_{n,m}`$ is a fundamental domain for $`\mathrm{\Gamma }_{n,m}\backslash _{n,m}.`$
Proof. The proof can be found in . $`\mathrm{}`$
In the case $`n=m=1`$, R. Berndt introduced the notion of Maass-Jacobi forms. Now we generalize this notion to the general case.
###### Definition 4.1.
For brevity, we set $`\mathrm{\Delta }_{n,m}:=\mathrm{\Delta }_{n,m;1,1}`$ (cf. Theorem 1.2). Let
$$\mathrm{\Gamma }_{n,m}:=Sp(n,)H_{}^{(n,m)}$$
be the discrete subgroup of $`G^J`$, where
$$H_{}^{(n,m)}=\left\{(\lambda ,\mu ;\kappa )H_{}^{(n,m)}\right|\lambda ,\mu ,\kappa \text{are integral}\}.$$
A smooth function $`f:_{n,m}`$ is called a Maass-Jacobi form on $`_{n,m}`$ if $`f`$ satisfies the following conditions (MJ1)-(MJ3) :(MJ1) $`f`$ is invariant under $`\mathrm{\Gamma }_{n,m}.`$
(MJ2) $`f`$ is an eigenfunction of the Laplacian $`\mathrm{\Delta }_{n,m}`$.
(MJ3) $`f`$ has a polynomial growth, that is, there exist a constant $`C>0`$ and a positive
integer $`N`$ such that
$$|f(X+iY,Z)|C|p(Y)|^N\text{as}detY\mathrm{},$$
where $`p(Y)`$ is a polynomial in $`Y=(y_{ij}).`$
It is natural to propose the following problems.
Problem A : Construct Maass-Jacobi forms.
Problem B : Find all the eigenfunctions of $`\mathrm{\Delta }_{n,m}.`$
We consider the simple case $`n=m=1.`$ A metric $`ds_{1,1}^2`$ on $`𝐇_1\times `$ given by
$`ds_{1,1}^2=`$ $`{\displaystyle \frac{y+v^2}{y^3}}(dx^2+dy^2)+{\displaystyle \frac{1}{y}}(du^2+dv^2)`$
$`{\displaystyle \frac{2v}{y^2}}(dxdu+dydv)`$
is a $`G^J`$-invariant Kähler metric on $`𝐇_1\times `$. Its Laplacian $`\mathrm{\Delta }_{1,1}`$ is given by
$`\mathrm{\Delta }_{1,1}=`$ $`y^2\left({\displaystyle \frac{^2}{x^2}}+{\displaystyle \frac{^2}{y^2}}\right)`$
$`+(y+v^2)\left({\displaystyle \frac{^2}{u^2}}+{\displaystyle \frac{^2}{v^2}}\right)`$
$`+\mathrm{\hspace{0.17em}2}yv\left({\displaystyle \frac{^2}{xu}}+{\displaystyle \frac{^2}{yv}}\right).`$
We provide some examples of eigenfunctions of $`\mathrm{\Delta }_{1,1}`$. (1) $`h(x,y)=y^{\frac{1}{2}}K_{s\frac{1}{2}}(2\pi |a|y)e^{2\pi iax}(s,`$ $`a0)`$ with eigenvalue $`s(s1).`$ Here
$$K_s(z):=\frac{1}{2}_0^{\mathrm{}}\mathrm{exp}\left\{\frac{z}{2}(t+t^1)\right\}t^{s1}𝑑t,$$
where $`\mathrm{Re}z>0.`$
(2) $`y^s,y^sx,y^su(s)`$ with eigenvalue $`s(s1).`$
(3) $`y^sv,y^suv,y^sxv`$ with eigenvalue $`s(s+1).`$
(4) $`x,y,u,v,xv,uv`$ with eigenvalue $`0`$.
(5) All Maass wave forms.
We fix two positive integers $`m`$ and $`n`$ throughout this section. For an element $`\mathrm{\Omega }_n,`$ we set
$$L_\mathrm{\Omega }:=^{(m,n)}+^{(m,n)}\mathrm{\Omega }$$
It follows from the positivity of $`\text{Im}\mathrm{\Omega }`$ that the elements $`f_{kl},h_{kl}(\mathrm{\Omega })(1km,1ln)`$ of $`L_\mathrm{\Omega }`$ are linearly independent over $``$. Therefore $`L_\mathrm{\Omega }`$ is a lattice in $`^{(m,n)}`$ and the set $`\left\{f_{kl},h_{kl}(\mathrm{\Omega })\right|1km,1ln\}`$ forms an integral basis of $`L_\mathrm{\Omega }`$. We see easily that if $`\mathrm{\Omega }`$ is an element of $`_n`$, the period matrix $`\mathrm{\Omega }_{}:=(I_n,\mathrm{\Omega })`$ satisfies the Riemann conditions (RC.1) and (RC.2) : (RC.1) $`\mathrm{\Omega }_{}J_n^t\mathrm{\Omega }_{}=0`$; (RC.2) $`\frac{1}{i}\mathrm{\Omega }_{}J_n^t\overline{\mathrm{\Omega }}_{}>0`$.
Thus the complex torus $`A_\mathrm{\Omega }:=^{(m,n)}/L_\mathrm{\Omega }`$ is an abelian variety. It might be interesting to investigate the spectral theory of the Laplacian $`\mathrm{\Delta }_{n,m}`$ on a fundamental domain $`_{n,m}`$. But this work is very complicated and difficult at this moment. It may be that the first step is to develop the spectral theory of the Laplacian $`\mathrm{\Delta }_\mathrm{\Omega }`$ on the abelian variety $`A_\mathrm{\Omega }.`$ The second step will be to study the spectral theory of the Laplacian $`\mathrm{\Delta }_n`$ (see (1.4)) on the moduli space $`\mathrm{\Gamma }_n\backslash _n`$ of principally polarized abelian varieties of dimension $`n`$. The final step would be to combine the above steps and more works to develop the spectral theory of the Lapalcian $`\mathrm{\Delta }_{n,m}`$ on $`_{n,m}.`$ Maass-Jacobi forms play an important role in the spectral theory of $`\mathrm{\Delta }_{n,m}`$ on $`_{n,m}.`$ Here we deal only with the spectral theory $`\mathrm{\Delta }_\mathrm{\Omega }`$ on $`L^2(A_\mathrm{\Omega }).`$
We fix an element $`\mathrm{\Omega }=X+iY`$ of $`_n`$ with $`X=\text{Re}\mathrm{\Omega }`$ and $`Y=\text{Im}\mathrm{\Omega }.`$ For a pair $`(A,B)`$ with $`A,B^{(m,n)},`$ we define the function $`E_{\mathrm{\Omega };A,B}:^{(m,n)}`$ by
$$E_{\mathrm{\Omega };A,B}(Z)=e^{2\pi i(\sigma (^tAU)+\sigma ((BAX)Y^1{}_{}{}^{t}V))},$$
where $`Z=U+iV`$ is a variable in $`^{(m,n)}`$ with real $`U,V`$.
###### Lemma 4.1.
For any $`A,B^{(m,n)},`$ the function $`E_{\mathrm{\Omega };A,B}`$ satisfies the following functional equation
$$E_{\mathrm{\Omega };A,B}(Z+\lambda \mathrm{\Omega }+\mu )=E_{\mathrm{\Omega };A,B}(Z),Z^{(m,n)}$$
for all $`\lambda ,\mu ^{(m,n)}.`$ Thus $`E_{\mathrm{\Omega };A,B}`$ can be regarded as a function on $`A_\mathrm{\Omega }.`$
###### Proof.
We write $`\mathrm{\Omega }=X+iY`$ with real $`X,Y.`$ For any $`\lambda ,\mu ^{(m,n)},`$ we have
$`E_{\mathrm{\Omega };A,B}(Z+\lambda \mathrm{\Omega }+\mu )`$ $`=E_{\mathrm{\Omega };A,B}((U+\lambda X+\mu )+i(V+\lambda Y))`$
$`=e^{2\pi i\{\sigma (^tA(U+\lambda X+\mu ))+\sigma ((BAX)Y^1{}_{}{}^{t}(V+\lambda Y))\}}`$
$`=e^{2\pi i\{\sigma (^tAU+^tA\lambda X+^tA\mu )+\sigma ((BAX)Y^1{}_{}{}^{t}V+B^t\lambda AX^t\lambda )\}}`$
$`=e^{2\pi i\{\sigma (^tAU)+\sigma ((BAX)Y^1{}_{}{}^{t}V)\}}`$
$`=E_{\mathrm{\Omega };A,B}(Z).`$
Here we used the fact that $`{}_{}{}^{t}A\mu `$ and $`B^t\lambda `$ are integral. ∎
###### Lemma 4.2.
The metric
$$ds_\mathrm{\Omega }^2=\sigma ((\text{Im}\mathrm{\Omega })^1{}_{}{}^{t}(dZ)d\overline{Z}))$$
is a Kähler metric on $`A_\mathrm{\Omega }`$ invariant under the action (5.15) of $`\mathrm{\Gamma }^J=Sp(n,)H_{}^{(m,n)}`$ on $`(\mathrm{\Omega },Z)`$ with $`\mathrm{\Omega }`$ fixed. Its Laplacian $`\mathrm{\Delta }_\mathrm{\Omega }`$ of $`ds_\mathrm{\Omega }^2`$ is given by
$$\mathrm{\Delta }_\mathrm{\Omega }=\sigma \left((\text{Im}\mathrm{\Omega })\frac{}{Z}^t\left(\frac{}{\overline{Z}}\right)\right).$$
###### Proof.
The proof can be found .
We let $`L^2(A_\mathrm{\Omega })`$ be the space of all functions $`f:A_\mathrm{\Omega }`$ such that
$$f_\mathrm{\Omega }:=_{A_\mathrm{\Omega }}|f(Z)|^2𝑑v_\mathrm{\Omega },$$
where $`dv_\mathrm{\Omega }`$ is the volume element on $`A_\mathrm{\Omega }`$ normalized so that $`_{A_\mathrm{\Omega }}𝑑v_\mathrm{\Omega }=1.`$ The inner product $`(,)_\mathrm{\Omega }`$ on the Hilbert space $`L^2(A_\mathrm{\Omega })`$ is given by
$$(f,g)_\mathrm{\Omega }:=_{A_\mathrm{\Omega }}f(Z)\overline{g(Z)}𝑑v_\mathrm{\Omega },f,gL^2(A_\mathrm{\Omega }).$$
###### Theorem 4.2.
The set $`\left\{E_{\mathrm{\Omega };A,B}\right|A,B^{(m,n)}\}`$ is a complete orthonormal basis for $`L^2(A_\mathrm{\Omega })`$. Moreover we have the following spectral decomposition of $`\mathrm{\Delta }_\mathrm{\Omega }`$:
$$L^2(A_\mathrm{\Omega })=_{A,B^{(m,n)}}E_{\mathrm{\Omega };A,B}.$$
Proof. The complete proof can be found in $`\mathrm{}`$ |
warning/0507/quant-ph0507050.html | ar5iv | text | # Generation of generalized coherent states with two coupled Bose-Einstein condensates.
## I Introduction
Combined advances in evaporative cooling techniques and magneto-optical trapping made it possible to create an atomic Bose-Einstein Condensate (BEC) experimentally, an important achievement of the last decade. Initially predicted by Einstein in 1925 Einstein (1925), it was produced in 1995 from a dilute gas of rubidium atoms Anderson et al. (1995). Other research groups produced condensates using sodium Davis et al. (1995), lithium Bradley et al. (1997) and hydrogen Fried et al. (1998). In a second generation of experiments, it was shown to be possible to create double condensates. Such a system can be constructed by trapping atoms in two different hyperfine sublevels of <sup>87</sup>Rb Myatt et al. (1997); Matthews et al. (1998).
Measurements of scattering lengths Hall et al. (1998a) and research on the dynamics of the relative phase of the condensates Hall et al. (1998b) and Rabi oscillations of the two BEC populations Matthews et al. (1999), can be carried out by using a laser-induced Raman transition in the <sup>87</sup>Rb experimental setup. The experimental production of the first BEC and the analogy between the behavior of the coherent matter waves and the electromagnetic ones, encouraged the development of Atom Optics Anderson and Meystre (2003). Nowadays, quantum optics tools are commonly used in the study of the BEC properties. Another consequence of this analogy is the study of problems already explored in quantum optics but mapped into the context of atomic systems. One example is the generation of “Schrödinger cat”-like states (SCS) whose creation in quantum optics via dynamical procedures involving nonlinear interaction, was proposed by Yurke and Stoler Yurke and Stoler (1986) and discussed by several authors Gerry et al. (2002); Gerry (1999); Agarwal et al. (1997); Agarwal and Banerji (1998). These schemes involve Kerr-like couplings and, in general, coherent states are used as initial states.
When two coupled BECs are analyzed by using a many-body Hamiltonian within the two-mode approximation (TMA), the terms that describe atomic collisions are analogous to a Kerr-like interaction Milburn et al. (1997). An additional Raman transition could be switched on so there is a Josephson-like coupling between the two modes. In this case, each mode corresponds to one of the BEC species and it is necessary to take into account inter- and intraspecies scattering processes. TMA was used by several authors to explore the possibility of creating quantum superposition states in BECs Cirac et al. (1998); Gordon and Savage (1999). Cirac et al. Cirac et al. (1998) calculated the ground state of the TMA Hamiltonian for various choices of coupling parameters. For certain sets of parameter values, the ground state is a SCS. Gordon and Savage Gordon and Savage (1999), among others, proposed the generation of SCS by exploiting the dynamical evolution of the system, in a similar fashion as has been done in the electromagnetic waves Yurke and Stoler (1986). Other aspects of BECs recently studied are the entanglement dynamics and the generation of entangled states Sørensen et al. (2001); You (2003); Micheli et al. (2003); Hines et al. (2003). A dynamical scheme was proposed by Micheli et al. Micheli et al. (2003) in order to generate a many-particle entangled state. In their approach, the entangled subsystems correspond to the individual atoms in BECs.
In the present contribution, we propose the generation of the Generalized Coherent State (GCS) in a system with two coupled BECs. The GCS was introduced by Titulaer and Glauber Titulaer and Glauber (1966) as a generalization of Glauber’s coherent state, defined as
$`|\text{GCS}`$ $`=`$ $`\text{exp}\left(\left|\gamma \right|^2/2\right){\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\gamma ^m}{\sqrt{m!}}}\text{exp}\left(i\upsilon _m\right)|m.`$ (1)
Here, the phases $`\upsilon _m`$ are functions of index $`m`$ which ensures Poisson excitation statistics. For the periodic case, when $`\upsilon _{m+l}=\upsilon _m`$, GCS can be written as a superposition of $`l`$ coherent states with equal mean value of excitation $`\left|\gamma \right|^2`$ Bialynicka-Birula (1968):
$$|\text{GCS}_{\text{periodic}}=\underset{k=1}{\overset{l}{}}c_k|\gamma \mathrm{exp}\left(i2\pi k/l\right)$$
Within the TMA we determine the conditions at wich two BECs, starting as a product of coherent states, first become entangled and, later, at certain specific times evolve to a product of the vacuum state and a GCS. We also show that the phases $`\upsilon _m`$ of the created GCS are periodic, and hence it can be rewritten as a superposition of $`l`$ coherent states. The period $`l`$ is fixed by both, the Josephson-like and nonlinear coupling strengths. We also explore numerically the dynamics of the system with the same initial state, in which the interspecies collision process is gradually inhibited. In these situations an exact GCS in no longer attained, but the evolved state has interesting properties such as sub-Poisson statistics, at the time GCS would have formed.
The paper is structured as follows. In Sec. II, we review the two-mode approximation, defining the parameters of interest in our calculation. Section III is reserved for the analysis of the necessary conditions to obtain a pure GCS. Also, we analyze the possibility of controlling the number of coherent states in the created superposition, by changing the coupling strengths. We estimate the evolution time necessary for the formation of the GCS. Section IV is devoted to the discussion of feasibility and sources of decoherence. Section V contains a numerical calculation of the dynamics of the system when the collisions between atoms of different species of BECs are inhibited. In Section VI, we summarize our results.
## II The two-mode model.
Our system consists of two atomic BECs of different atomic species labeled with suffixes $`a`$ and $`b`$, in a harmonic trap characterized by potentials $`V_{a,b}\left(r\right)`$. Interaction between atoms $`a`$ and $`b`$ are well described if we assume only two-body collisions. This can be done by considering three different scattering processes: $`aa`$, $`bb`$ and $`ab`$ atomic collisions. We are interested in the dynamics of this system when Josephson-like coupling between species $`a`$ and $`b`$ of BECs is switched on. The second quantized Hamiltonian which describes our system is given by Milburn et al. (1997); Cirac et al. (1998); Steel and Collett (1998); Villain and Lewenstein (1999)
$$\widehat{H}=\widehat{H}_a+\widehat{H}_b+\widehat{H}_{ab}+\widehat{H}_\text{c},$$
(2)
where
$`\widehat{H}_j`$ $`=`$ $`{\displaystyle }d^3\text{r}\widehat{\mathrm{\Psi }}_j^{}[{\displaystyle \frac{\mathrm{}^2}{2m}}^2+V_j\left(\text{r}\right)`$ (3)
$`+{\displaystyle \frac{4\pi \mathrm{}^2A_j}{2m}}\widehat{\mathrm{\Psi }}_j^{}\widehat{\mathrm{\Psi }}_j]\widehat{\mathrm{\Psi }}_j,`$
$`\widehat{H}_{ab}`$ $`=`$ $`{\displaystyle \frac{4\pi \mathrm{}^2A_{ab}}{m}}{\displaystyle d^3\text{r}\widehat{\mathrm{\Psi }}_a^{}\widehat{\mathrm{\Psi }}_b^{}\widehat{\mathrm{\Psi }}_a\widehat{\mathrm{\Psi }}_b},`$ (4)
$`\widehat{H}_\text{c}`$ $`=`$ $`{\displaystyle \frac{\mathrm{}\mathrm{\Omega }}{2}}{\displaystyle d^3\text{r}\left[\widehat{\mathrm{\Psi }}_a^{}\widehat{\mathrm{\Psi }}_be^{i\delta t}+\widehat{\mathrm{\Psi }}_b^{}\widehat{\mathrm{\Psi }}_ae^{i\delta t}\right]}`$ (5)
and $`j=a,b`$. Here, we have omitted spatial dependence in quantum field operators, $`\widehat{\mathrm{\Psi }}_{a,b}`$ ($`\widehat{\mathrm{\Psi }}_{a,b}^{}`$), which annihilate (create) atoms at position r. $`m`$ is the atomic mass and $`\widehat{V}_{a,b}\left(\text{r}\right)`$ are the harmonic trap potentials and $`A_{a,b}`$ are the scattering lengths associated with collisions between atoms of the same condensate (intraspecies collisions). Hamiltonian $`\widehat{H}_{ab}`$ describes the interaction between atoms of different species due to two-body collisions (interspecies collisions). $`H_\text{c}`$ is the Josephson-like coupling between the modes, $`\delta `$ being the detuning from Raman resonance and $`\mathrm{\Omega }`$ is the Rabi frequency.
Following a procedure similar to that described in Ref. Cirac et al. (1998), we obtain the TMA Hamiltonian. The field operators are written as $`\widehat{\mathrm{\Psi }}_a=\varphi _a\left(\text{r}\right)\widehat{a}`$ and $`\widehat{\mathrm{\Psi }}_b=\varphi _b\left(\text{r}\right)\widehat{b}`$, $`\varphi _{a,b}\left(\text{r}\right)`$ being the real spatial functions associated with each mode and $`\widehat{a}`$ and $`\widehat{b}`$ the standard bosonic operators. Additionally, we consider here $`\delta =0`$, to obtain the total Hamiltonian given by
$`\widehat{H}`$ $`=`$ $`\mathrm{}\omega _a\widehat{a}^{}\widehat{a}+\mathrm{}U_{aa}\widehat{a}^{}\widehat{a}^{}\widehat{a}\widehat{a}+\mathrm{}\omega _b\widehat{b}^{}\widehat{b}+\mathrm{}U_{bb}\widehat{b}^{}\widehat{b}^{}\widehat{b}\widehat{b}`$ (6)
$`+2\mathrm{}U_{ab}\widehat{a}^{}\widehat{a}\widehat{b}^{}\widehat{b}\mathrm{}\lambda \left(\widehat{a}^{}\widehat{b}+\widehat{a}\widehat{b}^{}\right),`$
with
$`\omega _j`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{}}}{\displaystyle d^3\text{r}\varphi _j\left(\text{r}\right)\left[\frac{1}{2}^2+\stackrel{~}{V}_j\left(\text{r}\right)\right]\varphi _j\left(\text{r}\right)},`$ (7a)
$`U_{jj}`$ $`=`$ $`{\displaystyle \frac{4\pi \mathrm{}A_j}{2m}}{\displaystyle d^3\text{r}\varphi _j^4\left(\text{r}\right)},`$ (7b)
$`U_{ab}`$ $`=`$ $`{\displaystyle \frac{4\pi \mathrm{}A_{ab}}{2m}}{\displaystyle d^3\text{r}\varphi _a^2\left(\text{r}\right)\varphi _b^2\left(\text{r}\right)},`$ (7c)
$`\lambda `$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }}{2}}{\displaystyle d^3\text{r}\varphi _a\left(\text{r}\right)\varphi _b\left(\text{r}\right)}.`$ (7d)
The TMA Hamiltonian (6) can be used in the description of two different experimental situations. The first one is the condensation of sodium, where atoms condense in hyperfine states localized in two different minima of the harmonic trap Andrews et al. (1997); Stenger et al. (1998). In this case, Josephson-like coupling describes tunneling. In some cases, a good approximation is obtained by neglecting the interspecies collisions. However, it is more general to assume that $`U_{ab}<U_{aa}=U_{bb}`$.
The second situation is connected with the experiments of the JILA group with condensation of atoms on two different hyperfine <sup>87</sup>Rb levels. In this context, the Josephson-like coupling is associated with a laser-induced Raman transition between the hyperfine levels. Reported scattering length values follow the relation $`A_a:A_{ab}:A_b1.03:1:0.97`$ Hall et al. (1998a, b). From Eqs.(7b) and (7c) it is clear that parameters $`U_{ij}`$ obey the same relations, for a fixed spatial mode function $`\varphi _{a,b}\left(\text{r}\right)`$. The latter is an important condition if we want to use the TMA: as we can see from Eqs.(7), the values of the strengths of the Hamiltonian depend on the spatial mode functions $`\varphi _{a,b}\left(\text{r}\right)`$. The approximation is valid only if these functions remain unaltered and the parameters in each term of Hamiltonian (6) can be considered as constants <sup>1</sup><sup>1</sup>1An estimative of validity of two-mode model can be found in Ref. Milburn et al. (1997). Also, in Section V of Ref. Gordon and Savage (1999), the authors discuss the different regimes in which this approximation is valid.. Several authors use $`A_a=A_b=A_{ab}`$ to simplify theoretical calculations with the Hamiltonian (2Park and Eberly (2000); Villain and Lewenstein (1999). In the TMA Hamiltonian (6), this situation corresponds to $`U_{ab}=U_{aa}=U_{bb}`$.
In this article, we assume that $`U_{aa}+U_{bb}=2U_{ab}`$ in order to extend the analytical solution of the Schrödinger equation in L. Sanz et al. (2003) and show how the GCS is exactly generated. Notice that this assumption applies for both cases: equal scattering lengths approximation ($`U_{aa}=U_{ab}=U_{bb}`$) and for the relation between experimental measured scattering lengths ($`U_{aa}:U_{ab}:U_{bb}1.03:1:0.97`$). Then, using numerical calculations, we explore the situation when $`U_{ab}<U=U_{aa}=U_{bb}`$. In this way, we are able to study the effect of the interspecies collision term on the dynamics and the transition between two different situations which can be related to the experimental contexts of rubidium and sodium ($`U_{ab}0`$) condensates.
## III Generation of Generalized Coherent States.
In this section, we show how the dynamical evolution associated with the TMA Hamiltonian can be exploited to produce a product of the vacuum state and the GCS. We first assume, reasoning by analogy with BECs in optical lattices Greiner et al. (2002a), that the system could be prepared as a product of coherent states $`|\mathrm{\Psi }(0)=|\alpha _a|\alpha _b`$ where $`\alpha _j`$ are the amplitude of the state thus $`\left|\alpha _j\right|^2`$ is the atomic population on mode-$`j`$. It is demonstrated Greiner et al. (2002a) that the manipulation of the Josephson-like coupling, by changing the potential depth between the $`\mathrm{}`$ local minima of the lattice, produces the state $`_{\mathrm{}}|\alpha _{\mathrm{}}`$. Another reason is that the coherent state satisfies the conditions for full coherence. In experiments, interference patterns between two BECs were observed Andrews et al. (1997) and collision-rate measurements Burt et al. (1997) probed the existence of third-order correlations. Although similar patterns could be obtained if BEC state is described either as Fock or coherent states You (2003); Castin and Dalibard (1999); Javanainen and Yoo (1996), studies about decoherence process due to three-body losses Jack (2002) supports the assumption that the state of a BEC is a coherent state with a well-defined phase. Also, phase and spatial dynamics were explored including the effect of fluctuations by Sinatra and Castin Sinatra and Castin (2000) and Ref. Villain and Lewenstein (1999). Results which are in agreement with the measure of relative phase between coupled condensate Hall et al. (1998b) were obtained by Li et al. Li et al. (2001) considering the initial state $`|\alpha _a|\alpha _b`$.
In this work we shall focus on the macroscopic superposition state resulting from the evolution of the system itself. Solving the Schrödinger equation associated with Hamiltonian (6), as shown in the Appendix A, the evolved state ($`\mathrm{}=1`$) is given by
$`|\mathrm{\Psi }(t)`$ $`=`$ $`e^{\frac{N}{2}}{\displaystyle \underset{n,m}{}}{\displaystyle \frac{\left[\alpha \left(t\right)\right]^n}{\sqrt{n!}}}{\displaystyle \frac{\left[\beta \left(t\right)\right]^m}{\sqrt{m!}}}e^{itU_{ab}\left(n^2m^2\right)}`$ (8)
$`\times e^{2itU_{ab}nm}e^{i\omega _0t\left(n+m\right)}|n,m`$
with
$`\alpha (t)`$ $`=`$ $`\alpha _a\mathrm{cos}\left(\lambda _1t\right)+i{\displaystyle \frac{\mathrm{sin}\left(\lambda _1t\right)}{\lambda _1}}\left(\lambda \alpha _b\omega _1\alpha _a\right),`$ (9a)
$`\beta (t)`$ $`=`$ $`\alpha _b\mathrm{cos}\left(\lambda _1t\right)+i{\displaystyle \frac{\mathrm{sin}\left(\lambda _1t\right)}{\lambda _1}}\left(\lambda \alpha _a+\omega _1\alpha _b\right)`$ (9b)
and
$`\omega _0`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[\omega _a+\omega _b2U_{ab}\right],`$ (10a)
$`\omega _1`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[\omega _a\omega _b+\left(U_{aa}U_{bb}\right)\left(N1\right)\right],`$ (10b)
$`\lambda _1`$ $`=`$ $`\sqrt{\lambda ^2+\omega _1^2},`$ (10c)
$`N`$ $`=`$ $`\widehat{N}=|\alpha _a|^2+|\alpha _b|^2,`$ (10d)
being $`N`$ the total excitation number of the system, which is a constant of motion and $`\lambda _1`$ the effective Rabi frequency.
We see that the state represented by Eq.(8) is an entangled state and there are only two situations where $`|\psi \left(t\right)`$ can be written as a direct product: The first is when the interaction parameter $`U_{ab}t`$ is a multiple of $`\pi `$ and the function $`\mathrm{exp}\left(2inmU_{ab}t\right)`$ in Eq.(8) is equal to unity. Thus, the disentanglement times associated with this first condition depend only on the value of the nonlinear coupling strength $`U_{ab}`$. At these times, $`|\psi \left(t\right)`$ can be rewritten as a direct product of new coherent states. The second case arises at those times such that either $`\alpha \left(t\right)`$ or $`\beta \left(t\right)`$ is zero, and the evolved state can be written as a product of the vacuum state and a superposition of Fock states.
The creation of the GCS given by Eq.(11) is restricted to the times associated with the second situation: if, for example, at certain evolution time $`t_e`$ the quantity $`\alpha (t_e)=0`$, the evolved state can be written as
$`|\psi \left(t_e\right)`$ $`=`$ $`|0e^{\frac{N}{2}}{\displaystyle \underset{n}{}}e^{iU_{ab}t_en^2}{\displaystyle \frac{\left[\beta \left(t_e\right)\right]^n}{\sqrt{n!}}}|n`$ (11)
$`=`$ $`|0|\text{GCS}.`$
Therefore, the GCS is a special superposition of Fock states and obeys Poisson statistics. From Eqs.(9a) and (9b) we conclude that the condition when either $`\alpha \left(t\right)`$ or $`\beta \left(t\right)`$ is zero can be written as
$$\alpha _j=i\mathrm{tan}\left(\lambda _1t_e\right)\left[\frac{\pm \omega _1\alpha _j\lambda \alpha _i}{\lambda _1}\right],$$
(12)
with $`ij=a`$ or $`b`$, depending on which quantity, $`\alpha (t)`$ or $`\beta (t)`$ goes to zero. Next, we analyze two particular choices of the interaction parameter, $`\lambda _1t_e`$, leading to the GCS:
1. At times given by
$$\lambda _1t_p=\frac{2p+1}{4}\pi $$
(13)
where $`p`$ is a positive or zero integer, the initial states satisfy the relation
$$\alpha _i=\frac{1}{\lambda }\left(\pm \omega _1+i\lambda _1\right)\alpha _j.$$
(14)
Note that in the particular case $`\omega _a=\omega _b`$ and $`U_{aa}=U_{ab}=U_{bb}`$ Park and Eberly (2000); Villain and Lewenstein (1999), we obtain, from Eqs.(10b) and (10c), that $`\omega _1=0`$ and $`\lambda _1=\lambda `$, and Eq.(14) is reduced to $`\alpha _j=i\alpha _i`$. Therefore, the initial mean number of atoms in each mode, $`\widehat{n}_a`$ and $`\widehat{n}_b`$, must be equal ($`\left|\alpha _a\right|^2=\left|\alpha _b\right|^2`$) with the relative phase $`\mathrm{\Delta }\varphi =\varphi _a\varphi _b`$ corresponding to $`\frac{\pi }{2}`$.
2. At times $`\lambda _1t_k=k\pi `$, we obtain the condition $`\alpha _j=0`$. Therefore, it is possible to generate the GCS if, for instance, initially all the $`N`$ atoms are condensed in the mode-$`a`$ and the second hyperfine level (mode-$`b`$) is used as an auxiliary mode. Thus, the atomic population leaves mode-$`a`$ and returns, not as a coherent state but as the GCS. For the initial state, $`|\mathrm{\Psi }(0)=|\sqrt{N}|0`$, we find that $`|\mathrm{\Psi }(t_k)=|\text{GCS}_{N,0}|0`$ with
$`|\text{GCS}_{N,0}`$ $`=`$ $`e^{\frac{N}{2}}{\displaystyle \underset{n}{}}{\displaystyle \frac{\left[\sqrt{N}e^{i\pi \frac{\omega _0}{\lambda _1}}\right]^n}{\sqrt{n!}}}`$ (15)
$`\times e^{ik\frac{U_{ab}}{\lambda _1}\pi n^2}|n.`$
This particular case is interesting because it is possible to create the GCS without the necessity to “imprint” any initial phase relation between the coupled BECs.
It is possible to rewrite the GCS as a superposition of coherent states Bialynicka-Birula (1968) if the phases given by $`\upsilon _n=U_{ab}t_en^2`$ on Eq.(11) are periodic. In our context, the necessary condition to obtain this kind of “Schrödinger cat”-like state is $`U_{ab}t_e=r/s`$, with $`r`$ and $`s`$ integers. This implies that the interspecies collision strength and effective Rabi frequency could also be written as a rational fraction. When this applies, we can use the discrete Fourier transform J. Banerji (2001) on Eq.(11). It is straightforward to rewrite $`|\text{GCS}`$ as the superposition
$$|C\left(t_p\right)=\underset{m=0}{\overset{l1}{}}a_m^{(r,s)}|\beta \left(t_e\right)e^{2\pi i\frac{m}{l}},$$
(16)
where $`l`$ is the number of coherent states present in the superposition. This value is defined by the condition below
$$l=\{\begin{array}{cc}2s\hfill & \text{if }r\text{ and }s\text{ are odd,}\hfill \\ s\hfill & \text{if }r\text{ is even and }s\text{ odd or vice versa.}\hfill \end{array}$$
(17)
The coefficients $`a_m^{(r,s)}`$ have the form
$$a_m^{(r,s)}=\frac{1}{l}\underset{k=0}{\overset{l1}{}}\mathrm{exp}\left(i\pi \frac{r}{s}k^2+2\pi i\frac{m}{l}k\right).$$
(18)
We see that the GCS corresponds to a superposition of coherent states with the same mean excitation number ($`|\beta |^2`$) and relative phases equal to $`e^{2\pi i\frac{m}{l}}`$. Note that the number of coherent states in the $`|\text{GCS}`$ depends on the ratio of nonlinear $`U_{ab}`$ to effective Rabi frequency $`\lambda _1`$.
In Figure 1 we plot the Husimi quasi-distribution function at time $`t_e=\frac{\pi }{4\lambda _1}`$ for $`\omega _a=\omega _b`$ and $`U_{aa}U_{bb}10^2U_{ab}`$. When $`U_{ab}`$ and $`\lambda _1`$ are chosen such that $`r/s=2/3`$, we obtain three distinguishable packets as shown in Fig. 1(a). Modifying the $`U_{ab}/\lambda _1`$ ratio it is possible to achieve a superposition of any number of coherent states. For instance, if $`U_{ab}/\lambda _1=8/5`$ we obtain superpositions of five coherent states as shown in Fig. 1(b). If $`U_{ab}/\lambda _1=1/2`$, we obtain eight packages, Fig. 1(c). Superposition of nine coherent states, shown in Fig. 1(d), is obtained when coupling strengths are set such that $`\frac{U_{ab}}{\lambda _1}=\frac{8}{9}`$. The last plot shows that the different gaussian packets, associated with different coherent states in $`|\text{GCS}`$, start to merge in phase space at high values of $`l`$. We also note that in Figs. 1(c-d) the deviation from a circular pattern of each gaussian packet in the superposition arises from interference between the packets, due to their proximity.
Because the effective Rabi frequency depends on the values of traps frequencies $`\omega _j`$ and collision parameters, all the examples above show that the formation of superpositions of coherent states in this scenario is highly sensitive to changes in these quantities and Josephson-like coupling. This means that, by controlling the values of scattering lengths, the effective harmonic potential, and the coupling between two species of condensates, it is possible to “build” a superposition of any desired number of coherent states, with a defined number of elements, mean excitation values and relative phases.
Next, we estimate the shortest time interval $`t_e`$ required to obtain the product state $`|0|\text{GCS}`$. The value of $`t_e`$ depends inversely on the effective Rabi frequency $`\lambda _1`$ ($`t_e\frac{\pi }{\lambda _1}`$). First, we assume Gaussian spatial mode functions $`\varphi _i\left(\text{r}\right)`$:
$$\varphi _i\left(\text{r}\right)=\left(\frac{1}{2\pi r_0^2}\right)e^{\text{r}^2/4r_0^2},$$
where $`r_0=\sqrt{\mathrm{}/2m\omega }`$, being $`\omega =\omega _a=\omega _b`$. Then, using the typical physical parameters of Rubidium experiments, $`\omega =50`$ s<sup>-1</sup>, $`m=1.4\times 10^{25}`$ Kg and the Rabi frequency $`\mathrm{\Omega }2\pi 600`$ s<sup>-1</sup> Hall et al. (1998a), we calculate the value of $`\lambda _1`$ from Eq.(10c). Thus, we obtain $`t_e10^3`$ s. It is important to note that $`t_e`$ can be set as short as possible by varying the Rabi frequency, $`\mathrm{\Omega }`$, of Raman transition.
## IV Discussion of Feasibility and Decoherence.
There are several questions about the feasiblity of the CGS arising from the results above. The first one is how to set the system in a convenient initial state. From all the possibilities suggested by Eq.(12), we conclude that the most reasonable initial condition is $`|\alpha _a|0`$. This state describes a condensate (in $`a`$-mode) and an empty auxiliary level ($`b`$-mode) described as a vacuum state. In this situation, the imprint of a relative phase between a pair of coupled BECs is not necessary.
Second question is the necessity of an efficient atomic population transference. If decoherence affects the process, we cannot guarantee the formation of the state $`|\text{GCS}|0`$. Following Ruostekoski and Walls Ruostekoski and Walls (1998) the effects of decoherence due to noncondensed atoms on BECs shows that purity decays fast, being lower than $`0.2`$ at $`U_{aa}t0.1`$. Hence, one must have the interaction parameter $`\lambda _1t`$ much smaller than the decoherence time scale associated with $`U_{aa}U_{ab}`$.
Once the GCS is created and Raman transition is switched off, it is necessary to check the effects of both, nonlinear interactions and decoherence. Because the fraction of noncondensed atoms is small, we can perform a simple calculation assuming that the interaction between those atoms and BEC induce phase-damping rather than atomic losses. Thus, we shall consider small the effect of decoherence due to condensate feeding and depleting. The following master equation for the density operator of $`a`$-mode, $`\widehat{\rho }_a`$, applies Anglin (1997); Louis et al. (2001)
$`{\displaystyle \frac{d\widehat{\rho }_a}{dt}}={\displaystyle \frac{i}{\mathrm{}}}[\widehat{H}_a,\widehat{\rho }_a]+\kappa \left(\{\widehat{n}_a^2,\widehat{\rho }_a\}2\widehat{n}_a\widehat{\rho }_a\widehat{n}_a\right).`$ (19)
with $`H_a=\mathrm{}\omega _a\widehat{a}^{}\widehat{a}+\mathrm{}U_{aa}\widehat{a}^{}\widehat{a}^{}\widehat{a}\widehat{a}`$. It is straightforward to calculate the solution of Eq.(19). They resemble the solutions for the phase-damped oscillator Gardiner and Zoller (2000):
$`\rho _a^{nm}\left(t\right)`$ $`=`$ $`e^{it\omega _a\left(nm\right)}e^{itU_{aa}\left[n\left(n1\right)m\left(m1\right)\right]}`$ (20)
$`\times e^{\kappa t\left(nm\right)^2}\rho _a^{nm}\left(0\right).`$
If $`\kappa =0`$, we are able to study the dynamics associated with atomic intraspecies collisions (nonlinear interaction term in $`\widehat{H}_a`$), assuming that we create successfully the GCS described by Eq.(15). From Eq.(20), the density matrix is given by
$$\widehat{\rho }_a(t)=|\text{GCS}(t)\text{GCS}(t)|$$
(21)
with
$`|\text{GCS}(t)`$ $`=`$ $`e^{\frac{N}{2}}{\displaystyle \underset{n}{}}{\displaystyle \frac{\mathrm{{\rm Y}}^n(t)}{n!}}e^{iU(t)n^2}|n,`$
$`\mathrm{{\rm Y}}(t)`$ $`=`$ $`\sqrt{N}\mathrm{exp}\left\{i\left[\pi {\displaystyle \frac{\omega _0}{\lambda _1}}+\left(\omega _aU_{aa}\right)t\right]\right\},`$
$`U(t)`$ $`=`$ $`\pi {\displaystyle \frac{U_{ab}}{\lambda _1}}+U_{aa}t.`$ (22)
Nonlinear collisions do not affect the character of the state and BEC is still in a GCS, with time-dependent amplitude $`\mathrm{{\rm Y}}(t)`$ and phase $`U(t)`$. From the analysis of Sec. III, we note that superpositions shown in Fig.1 can be destroyed as time progresses due to the changes on function $`U(t)`$ defined in Eq.(22). The effect of nonlinear interaction after the creation of the GCS could be reduced by manipulation of scattering length $`A_a`$ through a Feshbach resonance Vogels et al. (1997): controlling the scattering length, it is possible to change the value of $`U_{aa}t`$ so $`U(t)`$ varies smoothly with time.
For $`\kappa 0`$, we calculate $`\text{Tr}\left[\widehat{\rho }_a^2(t)\right]`$ in order to quantify the effects of the reservoir of noncondensed atoms on the BEC. We plot this quantity in Fig. 2 for different choices of $`\kappa `$ and number of atoms. From this results it is clear that, for large values of $`\kappa `$ the state is no longer a GCS neither a pure state. Additionally, the decay rate is sensitive to changes on $`\kappa `$ and $`N`$, decaying faster when both quantities increase. Hence, the phase damping is a serious limitation for manipulation of the GCS. Another source of decoherence is the three-body losses. Measurements of density-dependent losses demonstrate that three-body recombination is the dominant decoherence mechanism, which limits the lifetime and size of BECs Burt et al. (1997). In order to study the decoherence process due the three-body losses, a master equation is derived by M. Jack Jack (2002), where it is shown that a coherent state is a robust state in the limit of large-number of atoms. However, it is also shown that the superpositions defined in Eq.(16) are sensitive to the three-body losses.
Last concern is the effect of temperature and the ratio between Josephson and nonlinear couplings. A careful study of this effects on dephasing process was performed by Pitaevskii and Stringari Pittaevskii and Stringari (2001). They shown that coherence is strongly dependent on the ratio between Josephson coupling and collisions strength, $`U_{jj}/\lambda `$ and also with temperature, $`T/\lambda `$. One must have control over both ratios to keep them small in order to keep the phase coherence.
## V Inhibition of interspecies collisions.
In this section, we analyze the effect of the interspecies collisions on the dynamical evolution ($`U_{ab}<U=U_{aa}=U_{bb}`$). This is done by solving the Schrödinger equation numerically by direct diagonalization of the Hamiltonian (6) in a truncated Fock basis $`\left\{|n_a,n_b\right\}`$. In order to compare the results for $`U_{ab}<U`$ with those obtained when the condition $`U_{aa}+U_{bb}=2U_{ab}`$ is considered, we set the initial states as $`\left|\alpha _a\right|^2=\left|\alpha _b\right|^2`$ with relative phase $`\mathrm{\Delta }\varphi =\pi /2`$. We calculate the fraction of the total atom population in mode $`b`$, $`\widehat{n}_b/N=\widehat{b}^{}\widehat{b}/N`$, and its variance $`|\mathrm{\Delta }\widehat{n}_b|^2=\widehat{n}_b^2\widehat{n}_b^2`$. The “distance” between different states in the Fock basis can be analyzed using both $`|\mathrm{\Delta }\widehat{n}_b|^2`$ and the variance of operator $`\widehat{b}`$, defined as $`|\mathrm{\Delta }\widehat{b}|^2=\widehat{b}^{}\widehat{b}\widehat{b}^{}\widehat{b}`$ Kist et al. (1999). This last relation is useful to determine whether a given state can be considered as an eigenvalue of $`\widehat{n}_b`$ or $`\widehat{b}`$. The Mandel parameter
$$Q=\frac{|\mathrm{\Delta }\widehat{n}_b|^2\widehat{n}_b}{\widehat{n}_b},$$
(23)
and the linear entropy $`S_b=1\text{Tr}_b\left[\widehat{\rho }_b^2\left(t\right)\right]`$ are used, the first to characterize the statistics and the second to quantify the purity of the evolved state of mode $`b`$, respectively.
It is convenient to recall some well-known values for the definitions written above. For a coherent state ($`|\alpha `$), associated with $`\widehat{b}`$ and $`\widehat{b}^{}`$ operators, we obtain
$$\begin{array}{ccc}\widehat{n}_b& =\hfill & \left|\alpha \right|^2,\hfill \\ |\mathrm{\Delta }\widehat{n}_b|^2& =\hfill & \widehat{n}_b,\hfill \\ |\mathrm{\Delta }\widehat{b}|^2& =\hfill & 0,\hfill \\ Q& =\hfill & 0,\hfill \end{array}$$
(24)
indicating a Poisson statistics and that $`|\alpha `$ is an eigenstate of $`\widehat{b}`$. For a Fock state, $`|n`$, we obtain
$$\begin{array}{ccc}\widehat{n}_b& =\hfill & n,\hfill \\ |\mathrm{\Delta }\widehat{n}_b|^2& =\hfill & 0,\hfill \\ |\mathrm{\Delta }\widehat{b}|^2& =\hfill & n\hfill \\ Q& =\hfill & 1,\hfill \end{array}$$
(25)
indicating a sub-Poisson statistics and that $`|n`$ is an eigenstate of the $`\widehat{n}_b`$ operator. In order to compare these values with the numerical results, let us calculate the expressions above in the case of equal scattering lengths. Using the reduced density operator for mode $`b`$ extracted from Eq.(8), it is straightforward to obtain
$`\widehat{n}_b`$ $`=`$ $`\left|\beta \left(t\right)\right|^2,`$
$`|\mathrm{\Delta }\widehat{n}_b|^2`$ $`=`$ $`\left|\beta \left(t\right)\right|^2,`$
$`|\mathrm{\Delta }\widehat{b}|^2`$ $`=`$ $`\left|\beta \left(t\right)\right|^2\left\{1e^{2N\left[1\mathrm{cos}\left(2Ut\right)\right]}\right\},`$
$`Q`$ $`=`$ $`0.`$ (26)
Except for the Mandel parameter, which is time independent, all the functions depend on $`\left|\beta \left(t\right)\right|^2`$, which is the mean atom population in mode $`b`$, written as
$`\left|\beta \left(t\right)\right|^2`$ $`=`$ $`(\left|\alpha _a\right|^2+\left|\alpha _b\right|^2)\mathrm{cos}^2\lambda t`$ (27)
$`{\displaystyle \frac{i}{2}}\left(\alpha _a\alpha _b^{}\alpha _a^{}\alpha _b\mathrm{sin}2\lambda t\right),`$
where $`\left|\alpha \left(t\right)\right|^2+\left|\beta \left(t\right)\right|^2=N`$. From Eqs.(26), we can recover the result obtained from the analysis in Sec. III of the evolved state. The dynamics depends strongly on the relative phase and the initial population of both condensates. From the behavior of the partial population, $`\widehat{n}_b`$, assuming $`\mathrm{\Delta }\varphi =0`$ and $`\left|\alpha _a\right|=\left|\alpha _b\right|`$, we obtain $`\left|\beta \left(t\right)\right|^2=\left|\alpha _b\right|^2`$ and there is no transfer of population between the condensates. However, if $`\mathrm{\Delta }\varphi =\pi /2`$ we can see that $`\widehat{n}_b=\left|\alpha _b\right|^2\left[1\mathrm{sin}\left(2\lambda t\right)\right]`$ and the system undergoes Rabi oscillations with period equals to $`\pi /\lambda `$.
We also observe that the variance $`|\mathrm{\Delta }\widehat{b}|^2`$ is zero at times corresponding either to $`\pi /U`$ or when $`\left|\beta \left(t\right)\right|^2`$ goes to zero. Since the reduced linear entropy is also zero at these times, as we discussed in Sec.III, the variance $`|\mathrm{\Delta }\widehat{b}|^2`$ indicates that mode $`b`$ is in a coherent state. The Poisson statistics remains as time passes, independently of the entanglement dynamics of both modes.
In Fig. 3, we analyze the evolution of the atomic fraction in mode $`b`$, $`\widehat{n}_b/N`$, the Mandel parameter $`Q`$ and the linear entropy, $`\delta _b`$, for decreasing values of interspecies collision strength. We also plot the dynamics of each variable associated with the condition of equal scattering lengths, shown by a solid gray line. The vertical thick gray line indicates the time scale for formation of a GCS, $`t_e`$. The first aspect to be noticed in Fig.3(a) is a shift in the effective Rabi frequency of $`\widehat{n}_b/N`$ oscillations with decreasing $`U_{ab}`$. Also, there is an attenuation of Rabi oscillations if we compare both cases $`U_{ab}=U`$ and $`U_{ab}=0`$, shown in the inset. In particular, there are times at which the transfer of population is suppressed and atoms in each condensate are trapped. This “self-trapping” phenomenon was discussed elsewhere Milburn et al. (1997). We want to point out that the self-trapping can be associated with inhibition of interspecies collision and it is found even at slight differences between $`U`$ and $`U_{ab}`$.
The dynamics of the Mandel parameter, Fig.3(b), shows that the subsystem state presents sub-Poisson statistics at short times, with $`Q`$ becoming more negative as $`U_{ab}`$ decreases. Super-Poisson statistics are obtained at later times. Our results show that an appropriate manipulation of $`U_{ab}`$ leads the initial coherent state to a new one, with sub or super-Poisson statistics depending on the evolution time. It is interesting to note that, for $`U_{ab}=0`$, the time necessary to reach the minimum of Mandel parameter $`Q`$ is almost the same as that of the formation of GCS.
From our results for $`S_b`$, Fig. 3(c), we see that the entanglement process is now irreversible. A small change of $`U_{ab}`$ produces an increase in the linear entropy and the subsystems are unable to recover purity. In view of this result, we conclude that slight changes in our conditions for generation of the GCS destroys such a state. The variance of $`\widehat{b}`$ operator is shown in Fig. 4, for the same values of $`U_{ab}`$ as in Fig. 3. In all cases, there are no times at which $`|\mathrm{\Delta }\widehat{b}|^2=0`$, as in the case of equal scattering lengths, and mode $`b`$ never returns to a coherent state.
At this point, it is worth recalling that in the sodium condensate, $`\lambda `$ is associated with the tunnelling frequency between the two minima of potential, which depends on the width of the barrier. We can assume the values used by Milburn et al. in order to check the time scales in this case. With $`U`$ approximately $`53`$ s<sup>-1</sup> and $`\lambda 0.37\times 10^3`$ s<sup>-1</sup>, we estimate $`T_U6\times 10^2`$ s and $`t_e2\times 10^3`$. Again, the time scales for the formation of states with sub-Poisson statistics is of the order of milliseconds and shorter than the time scale associated with the internal collisions.
## VI Summary
Using the TMA Hamiltonian for the description of two coupled Bose-Einstein condensates, we demonstrate the possibility of creating a generalized coherent state in one of the condensate modes. The procedure presented here implies only dynamical evolution and requires the preparation of BECs in coherent states, which must follow the condition given by Eq.(12). The time necessary to obtain such a state depends only on the effective Rabi frequency $`\lambda _1`$, which is a function of Hamiltonian parameters. Also, it is shown that the ratio between the collision parameter $`U_{ab}`$ and $`\lambda _1`$ defines the number of coherent states contributing to the GCS. For $`U_{ab}<U_{aa}=U_{bb}`$, a new kind of non-classical statistics state is created. The analysis of fractional population, Mandel parameter and variance of the annihilation operator $`\widehat{b}`$ shows some interesting dynamical effects associated to this state. Such effects are, for instance, a shift of the effective Rabi frequency, some temporal regimes with sub-Poisson and super-Poisson statistics and irreversible entanglement.
###### Acknowledgements.
L. S. likes to thank E. I. Duzzioni, F. O. Prado and R. M. Angelo for helpful discussions. The authors also thanks to the referee for all the valuable critics. This work was supported by FAPESP (Fundação de Amparo à pesquisa do Estado de São Paulo) under grants 03/06307-9 and 00/15084-5 and CNPq (Instituto do Milênio de Informação Quântica).
## Appendix A Evolved state for TMA Hamiltonian with $`U_{aa}+U_{bb}=2U_{ab}`$.
In this appendix, we calculate a general solution of Schrödinger equation associated with Hamiltonian (6) by means of the unitary transformation
$$\widehat{V}\left(\gamma \right)=e^{\frac{\gamma }{2}\left(\widehat{a}^{}\widehat{b}\widehat{a}\widehat{b}^{}\right)}.$$
(28)
With this goal, we rewrite the Hamiltonian (2) using the number operator, $`\widehat{N}=\widehat{n}_a+\widehat{n}_b`$, and the unbalance population operator, $`\mathrm{\Delta }\widehat{n}=\widehat{n}_a\widehat{n}_b`$. If $`U_{aa}+U_{bb}2U_{ab}=0`$, we obtain
$$\widehat{H}=\omega _0\widehat{N}+\omega _1\mathrm{\Delta }\widehat{n}+U_{ab}\widehat{N}^2\lambda \left(\widehat{a}^{}\widehat{b}+\widehat{a}\widehat{b}^{}\right),$$
(29)
where $`\omega _0`$ and $`\omega _1`$ are the quantities defined in Eqs.(10a) and (10b). Using the relations
$`\widehat{V}^{}\widehat{a}\widehat{V}`$ $`=`$ $`\widehat{a}\mathrm{cos}\gamma /2+\widehat{b}\mathrm{sin}\gamma /2,`$
$`\widehat{V}^{}\widehat{a}^{}\widehat{V}`$ $`=`$ $`\widehat{a}^{}\mathrm{cos}\gamma /2+\widehat{b}^{}\mathrm{sin}\gamma /2,`$
$`\widehat{V}^{}\widehat{b}\widehat{V}`$ $`=`$ $`\widehat{b}\mathrm{cos}\gamma /2\widehat{a}\mathrm{sin}\gamma /2,`$
$`\widehat{V}^{}\widehat{b}^{}\widehat{V}`$ $`=`$ $`\widehat{b}^{}\mathrm{cos}\gamma /2\widehat{a}^{}\mathrm{sin}\gamma /2,`$ (30)
and choosing the unitary transformation parameter $`\gamma =\mathrm{arccos}\left(\omega _1/\lambda _1\right)`$, we obtain the transformed Hamiltonian
$$\widehat{H}_V=\omega _0\widehat{N}+U_{ab}\widehat{N}^2+\lambda _1\mathrm{\Delta }\widehat{n}.$$
(31)
The effective Rabi frequency $`\lambda _1=\sqrt{\lambda ^2+\omega _1^2}`$ depends on the differences between trap frequencies and collision strengths $`U_{jj}`$ as can be seen from Eqs.(10). It is straightforward to find the time propagator operator
$$|\mathrm{\Psi }(t)=\widehat{V}e^{i\widehat{H}_Vt}\widehat{V}^{}|\mathrm{\Psi }(0).$$
(32)
and, considering the initial state $`|\mathrm{\Psi }(0)=|\alpha _a|\alpha _b`$, we finally obtain the evolved state (8) with the quantities $`\alpha (t)`$ and $`\beta (t)`$ given by Eq.(9). |
warning/0507/cond-mat0507208.html | ar5iv | text | # Stimulated emission from the biexciton in a single quantum dot
## Abstract
Using two optical pulses of different frequencies, we demonstrate entanglement and disentanglement of the electronic states in Stranski-Krastanov quantum dots. Resonant two-photon excitation of the biexciton creates an entangled Bell-like state. The second pulse, being resonant to the exciton-biexciton transition, stimulates the emission from the biexciton and fully disentangles the two-bit system. By setting the polarization of the stimulation pulse, we control the recombination path of the biexciton and, by this, the state of the photons emitted in the decay cascade.
Two-photon absorbtion, Stimulated emission, Rabi oscillations, Quantum dots
Entangled states and their manipulation play an important role in quantum information processing book . A state of a pair of quantum systems is said to be entangled if it cannot be factored into the states of the individual subsystems. When aiming at practical devices, implementations in solid state are of particular interest. Here, a critical issue is decoherence that destroys the non-classical quantum correlations.
Semiconductor quantum dots (QDs) with electronic excitations localized on a nanometer length scale have recently attracted much attention as building blocks for quantum logic Steel ; Li . The two-exciton subspace of a QD is spanned by the ground-state $`|\text{g}`$ (no exciton), the single-exciton states, and the biexciton state $`|\text{b}`$. The optically active exciton is split by anisotropy in two linearly cross-polarized components $`|\text{x}`$ and $`|\text{y}`$ Gammon-Bacher . Thus, identifying $`|00=|\text{g}`$, $`|10=|\text{x}`$, $`|01=|\text{y}`$, and $`|11=|\text{b}`$, a two-bit system is built. The optical couplings between those states form the $`\text{V}\mathrm{\Lambda }`$ transition scheme of Fig. 1a, where the exciton-biexciton resonances are low-energy shifted from that of the single excitons by the exciton-exciton interaction energy $`\mathrm{\Delta }E_{\text{XX}}`$ in the biexciton.
Coherent optical coupling of the ground-state and the biexciton creates an entangled state of the type $`a_{00}|\text{00}+a_{11}|\text{11}`$ Hohenester ; Biolatti-Sham . A Bell state of maximum entanglement is achieved if $`|a_{ii}|=\frac{1}{\sqrt{2}}`$. Excitation of the biexciton $`|\text{g}|\text{b}`$ requires two photons. In Ref.Chen , two non-degenerate beams, each being resonant to one of the single-photon transitions, e.g. $`|\text{g}|\text{x}`$ and $`|\text{x}|\text{b}`$, have been utilized. The present study is based on resonant two-photon (TP) excitation TPCC . That approach, where the degenerate photons have half the energy of the biexciton, has the advantage that the spontaneous emission from the cascaded biexciton-exciton decay is clearly separated from the excitation stray light and can thus be used as a monitor of the quantum dynamics. The ability to track the QD emission is essential, as it enables to generate non-classical light states, like anti-bunched single-photon emission Michler-Santori or entangled photon pairs Benson .
In what follows we demonstrate that an entangled state in the two-exciton subspace of a QD can be indeed created by resonant TP excitation. The entanglement of formation reaches values of about 0.5. In a next step, we accomplish the conversion of the biexciton into an exciton by stimulated emission applying a second pulse that is tuned to the exciton-biexciton resonance. In this way, the recombination path is strictly defined by the polarization of the stimulation pulse, in contrast to the spontaneous emission cascade. In a quantum information sense, the stimulated emission corresponds to a disentanglement of the Bell-like state $`a_{00}|\text{00}+a_{11}|\text{11}(a_{00}|\text{0}+a_{11}|\text{1})|0`$. Disentanglement is a key step in conditional quantum dynamics and logical gates Barenco . We find efficiencies of close to 1 in our measurements.
The Stranski-Krastanov CdSe/ZnSe QD structures are grown by molecular beam epitaxy Litvinov . II-VI QDs are favored in the present context by the large $`\mathrm{\Delta }E_{\text{XX}}`$ 20 meV Kulakovski-Kreller , allowing the use of ultra-short pulses without loosing spectral selectivity of the excitation process. In order to study individual QDs, mesa structures with an area down to $`100\times 100`$ nm<sup>2</sup> are fabricated. The sum frequency of a Kerr-lens mode-locked Ti:sapphire laser and a synchronously pumped optical parametric oscillator are used to generate spectrally broad sub-ps pulses with 76 MHz repetition rate in the spectral region of interest. Two spectrally narrow pulses for selective TP excitation and stimulation (ST) of the biexciton are obtained by a pulse shaper based on a programmable reflective spatial light modulator (Fig. 1b). The outgoing pulses are spatially separated and passed through a delay line. The pulse durations are about 1.0 ps (TP) and 1.5 ps (ST) and the spectral full widths at half maximum are 1.7 meV and 1.2 meV, respectively (Fig. 1c). The secondary emission of the QD is collected in a confocal arrangement and dispersed in a triple spectrometer (0.23 nm/mm) equipped with a nitrogen cooled charged coupled device. For time-resolved measurements, only the two first stages are utilized in subtractive mode (0.7 nm/mm). A multi-channel-plate photomultiplier in conjunction with time-correlated single-photon counting unit provides an overall time resolution of 60 ps. Polarization control is achieved by quarter-wave or half-wave plates, placed in the path of both the linearly polarized excitation light as well as the emission signal. The polarization of the emission is analyzed by a Glan-Thomson prism introduced in front of the spectrometer. All measurements are carried out at temperatures of about 10 K.
The cascaded spontaneous emission of a QD subsequent to TP excitation is depicted in Fig. 2. For linearly polarized excitation, distinct exciton and biexciton features, placed symmetrically to the excitation photon energy, are present. In full accord with the transition scheme of Fig. 1 a, both features consist of a fine structure doublet of linearly cross-polarized lines, however, with a reversed sequence of the polarization in the exciton and biexciton emission. The TP transition of the biexciton is forbidden for circular excitation polarization. While the emission yield decreases indeed by one order of magnitude, a weak rest emission at the exciton survives. This background originates from electron-hole pairs off-resonantly excited in the energy continuum of the hetero-structure and captured by the QD. On the other hand, under linearly polarized excitation, the emission signal at the exciton lines is entirely insensitive on the polarization direction, excluding that the TP pulse addresses directly the exciton states to a measurable extent. A flip between $`|\text{x}`$ and $`|\text{y}`$ takes place on a time-scale markedly longer than the life-time and can be thus ignored Flissi . The time-resolved emission shown in Fig. 2b clearly confirms the existence of an emission cascade, with the biexciton recombining first and with a respective rise time for the exciton. Double-exponential fits yield that the radiative life-time of the biexciton ($`\tau _{\text{XX}}`$ = 120 ps) is about two times shorter than for the exciton ($`\tau _\text{X}`$ = 200 ps). Since the biexciton has two recombination channels this means that the exciton and the exciton-biexciton transition posses almost the same dipole moment. Using $`d^2=3\pi \epsilon _0\mathrm{}c^3/n_\text{r}\omega _0^3\tau `$, we find $`d=`$ 27 Debye.
Now we focus on the stimulated emission of the biexciton. In these measurements, the QD is excited by a sequence of TP and ST pulses, linearly co-polarized under an angle $`\phi `$ relative to the intrinsic polarization $`\stackrel{}{\pi }_\text{x}`$ and with tuneable time-delay $`\tau `$. The result of the stimulation process is a photon as well as an exciton. While the photon can hardly be detected, the stimulated exciton is manifested by the photon that it emits subsequently. The data in Fig. 3a directly verify this scenario. While both exciton emission components have equal intensity for TP excitation only, the line polarized along the ST polarization is amplified at expense of the cross-polarized line when both pulses are present. During pulse delay, the biexciton state is increasingly emptied by spontaneous decay. Consistent with the radiative life-time, the transition is not longer capable of stimulated emission after about 250 ps. A measure to what extend a certain exciton state can be selected by the ST pulse is given by the induced linear polarization degree $`\rho _\text{L}=(I_\text{x}I_\text{y})/(I_\text{x}+I_\text{y})`$, $`I_i`$ being the spectrally integrated signal of $`|i`$. As seen in Fig. 3c, $`\rho _\text{L}=0`$ for $`\phi =\pi /4`$, while, exciting along the intrinsic polarization axis, $`\rho _\text{L}`$ has the same absolute value, but is positive for $`\phi =0`$ and negative for $`\phi =\pi /2`$. In an incoherent regime, the polarization degree is limited to $`\rho _\text{L}=0.5`$, because the inversion between the biexciton and exciton population, addressed for a given ST polarization, can not exceed zero. In contrast, $`\rho _\text{L}`$ clearly exceeds this limit and reaches values of up to 0.8 for the maximum pulse densities available by our set-up. Even the onset of Rabi oscillations for the exciton-biexciton transition is evidenced in Fig. 3b. The relatively large fine structure splitting of the exciton allows us to spectrally separate spontaneous and stimulated emission. For QDs where the splitting is within the homogeneous width the stimulation process can be also accomplished, but shows only up in the polarization degree.
In order to draw quantitative conclusions on the quantum dynamics behind the experimental observations, the full density matrix of the two-exciton subspace has to be considered. Its time evolution is determined by the Master equation $`ı\mathrm{}\dot{\varrho }=[H,\varrho ]+ı\mathrm{}\mathrm{\Gamma }[\rho ]`$ with the Hamiltonian given in units of $`\mathrm{}`$ and rotating-wave approximation by
$`H`$ $`=`$ $`\omega _\text{g}|\text{g}\text{g}|+\omega _\text{x}|\text{x}\text{x}|+\omega _\text{y}|\text{y}\text{y}|)+\omega _\text{b}|\text{b}\text{b}|`$
$`{\displaystyle \frac{1}{2}}\{\mathrm{\Omega }(t)\text{e}^{ı\omega t}[\mathrm{cos}(\phi )(|\text{x}\text{g}|+|\text{b}\text{x}|)`$
$`+\mathrm{sin}(\phi )(|\text{y}\text{g}|+|\text{b}\text{y}|)]+\text{h.c.}`$
$`\mathrm{\Omega }=\frac{d}{\mathrm{}}(t)`$ is the time-dependent Rabi frequency of the field amplitude $`(t)`$, whereby the same $`d`$ is taken for all transitions. Expressed in the two-exciton basis, the Master equation defines a complete set of differential equations for the 10 independent density matrix elements $`\varrho _{ij}=\varrho _{ji}^{}`$. For the relaxation terms, we use $`i|\mathrm{\Gamma }[\varrho ]|j=\mathrm{\Gamma }_{ij}\varrho _{ij}(ij),\text{b}|\mathrm{\Gamma }[\varrho ]|\text{b}=\varrho _{\text{bb}}/\tau _{\text{XX}},i|\mathrm{\Gamma }[\varrho ]|i=\varrho _{\text{bb}}/2\tau _{\text{XX}}\varrho _{ii}/\tau _\text{X}(i=\text{x,y})`$, and $`\dot{\varrho }_{\text{bb}}+\dot{\varrho }_{\text{xx}}+\dot{\varrho }_{\text{yy}}+\dot{\varrho }_{\text{gg}}=0`$. Except the off-diagonal damping rates, all parameters $`(\mathrm{\Delta }E_{\text{xx}},E_\text{x}E_\text{y},d,\tau _{\text{XX}},\tau _\text{X})`$ entering the equations are known experimentally. Therefore, though the set is relatively large, the population dynamics is fully defined and the $`\mathrm{\Gamma }_{ij}`$, describing the coherence decay, can be deduced by comparing with the experimental data. The dynamical response is independent on whether $`|\text{x}`$ or $`|\text{y}`$ is addressed. Hence, $`\mathrm{\Gamma }_{\text{xg}}=\mathrm{\Gamma }_{\text{yg}}`$, $`\mathrm{\Gamma }_{\text{bx}}=\mathrm{\Gamma }_{\text{by}}`$, and making the reasonable assumption $`\mathrm{\Gamma }_{\text{bx}}=\mathrm{\Gamma }_{\text{xg}}+\mathrm{\Gamma }_{\text{bg}}`$, the only two free parameters left are the exciton ($`\mathrm{\Gamma }_{\text{xg}}`$) and biexciton ($`\mathrm{\Gamma }_{\text{bg}}`$) decoherence rates.
We have numerically solved the equations for Gaussian pulse shapes assuming $`\varrho _{\text{gg}}=1`$ and all other $`\varrho _{ij}=0`$ before the TP pulse and taking the resultant $`\varrho _{ij}`$ as initial values for the interaction with the ST pulse. In a time-integrated detection mode, the linear polarization degree is given by $`\rho _\text{L}=(\delta \varrho _{\text{xx}}\delta \varrho _{\text{yy}})/(\delta \varrho _{\text{bb}}+\delta \varrho _{\text{xx}}+\delta \varrho _{\text{yy}})`$, $`\delta \varrho _{ii}`$ denoting the change of the population generated by the pulses. Calculating $`\rho _\text{L}`$ as a function of the pulse delay yields indeed perfect agreement with experimental curves in Fig. 3c, verifying that the population dynamics is governed by the life-times $`\tau _{\text{XX}}`$ and $`\tau _\text{X}`$. The decoherence rates follow from the density dependence of $`\rho _\text{L}`$ (Fig. 3b). Too large rates spoil rapidly the purity of the TP excitation, resulting in a significant polarization degree even without the ST pulse, while the maximum level of only $`\rho _\text{L}=0.8`$ signifies the presence of pure dephasing beyond the radiative damping. The fit to the data yields $`\mathrm{\Gamma }_{\text{xg}}^1\mathrm{\Gamma }_{\text{bg}}^1=6`$ ps. A careful analysis of the spectral line-shape of the exciton emission provides that the radiative Lorentzian is superimposed to a weak but broad acoustic phonon background, which translates in non-exponential damping in the time domain with a short component consistent with the rates deduced from the ST data (see also Borri-Besombes ). Distinct non-exponential damping is also indicated by TP coherent control measurements TPCC . Here, a significant drop of the contrast occurs right after pulse separation, whereas the subsequent decay evolves on a much longer time-scale. Fig. 4 represents plots of the density matrix elements versus pulse area $`\theta =_{\mathrm{}}^+\mathrm{}\mathrm{\Omega }(t)𝑑t`$ in the experimentally relevant pulse-density range. Despite of the relatively fast dephasing, the TP pulse creates a significant biexciton coherence $`\varrho _{\text{bg}}`$ that can be manipulated with the succeeding ST pulse. Damped Rabi oscillations have been already observed previously Steel ; Li ; Rabi . However, unlike our study, where we monitor an absolute quantity, no conclusions about the true coherence degree could be made.
Knowing the whole density matrix, the entanglement of formation $`E(\varrho )`$ can be calculated. Following Ref.Wootters , $`E=E(C)=h(\frac{1+\sqrt{1C^2}}{2})`$ with the concurrence $`C`$ and $`h(x)=x\mathrm{log}_2(x)(1x)\mathrm{log}_2(1x)`$. For a purely coherent regime, it is straightforward to show that $`C=2\sqrt{\varrho _{\text{gg}}\varrho _{\text{bb}}}`$. While we can come with the TP pulse experimentally close to a situation $`\varrho _{\text{bb}}\varrho _{\text{gg}}\frac{1}{2}`$, the actual entanglement is markedly lower than 1. Decoherence establishes a mixed state, where now $`C(\varrho )=\text{max}[0,\lambda _1\lambda _2\lambda _3\lambda _4]`$, the $`\lambda `$’s being in decreasing order the square roots of $`\varrho \sigma _y\sigma _y\varrho ^{}\sigma _y\sigma _y`$ with the time inversion operator $`\sigma _y`$ Wootters . The calculation provides a maximum entanglement of about 0.5. However, this limited value given, complete disentanglement of the 2-quantum-bit system is accomplished by the ST pulse in a wide range of pulse areas (Fig. 4).
In conclusion, we have demonstrated entanglement and disentanglement of the electronic states of a QD by a sequence of two optical pulses. Being able to control the recombination path of the biexciton, the cascade emission can be adjusted from a pair of correlated or even entangled photons to a single photon of defined polarization or in an entangled polarization state. While an extrinsic background spoiling the figure-of-merit can be eliminated by appropriate sample design, our measurements uncover also rather large decoherence rates, probably related to a non-Markovian acoustic phonon contribution. This point deserves further investigations beyond the scope of this work. Our results are also of relevance for the use of QD in lasers, as they show that the exciton level is rapidly emptied avoiding reabsorption.
The authors thank S. Rogaschewski for the lithographic etching. This work was supported by the Deutsche Forschungsgemeinschaft within Project No. He 1939/18-1. |
warning/0507/cond-mat0507047.html | ar5iv | text | # Thermodynamics of An Ideal Generalized Gas: I Thermodynamic Laws
## I incomparable thermodynamic laws
Thermodynamics distinguishes itself on being able to take into account processes involving heat transfer. According to J. J. Thomson JJ (1888), heat is the ‘uncontrollable’ form of work, and temperature is its measure. J. P. Joule appreciated that heat and work were interconvertible, depending only on a constant for the units chosen to measure heat and mechanical work.
The first law was formulated as an expression for the conservation of internal energy even in the presence of heat transfer. The second law placed limitations on the amount of heat that could be converted into mechanical work. According to William Thomson (Lord Kelvin) Kelvin (1851)
> it is impossible to construct an engine which when operated in a cycle will produce no effect other than the extraction of heat from a reservoir and the performance of an equivalent amount of work.
The second law also succinctly sums up observations of nature like “heat always passes from hotter to colder, and never in the reverse direction”Carathéodory (1909), at constant volume. It also asserts that “heat will be absorbed as a gas expands to keep its temperature constant”.
Nicolas-Léonard-Sadi Carnot Carnot (1824) not only discovered that there is an upper bound to the efficiency of engines operating in a closed cycle, but, also called attention to a truly mathematical invariant quantity that inspired Thomson’s later researches into the discovery of that entity. This invariant would be the same for all susbstances at the same temperature, and Thomson, equated it with the temperature measured on the ideal gas absolute scale. This factor proved to be an integrating denominator for the quantity of heat, and became a matter of contention between R. Clausius and Thomson for its true authorship Tait (1868). Thus, the entropy, like the internal energy, and in contrast to the quantity of heat, was identified as a point function which had the advantage of depending only on the instantaneous state of the body, and not on the process by which the body arrived in that state.
To the best of our knowledge there has never been any question of the compatibility of the first and second laws of thermodynamics. Certainly, the two laws are compatible for a classical ideal gas (ICG) because of the separability of the temperature and volume afforded by the logarithmic function. However, this is not true, in general, for an ideal generalized gas (IGG), or a ‘quantum’ gas (IQG) Landsberg (1961), which is the low temperature extension of an ICG. Here, we make the distinction between an IGG and an ICG in that the former can be valid at all temperatures and not only in the high temperature limit. This means that the transition between and IGG and an ICG occurs in a manner different than taking the high temperature limit, and one in which we shall explore in the last section of this paper.
For an IGG, as well as an IQG, power laws are involved that contain products of different powers of the temperature and volume. And, in fact, one finds that for all processes other than those involving pure heat conduction, the power means derived from the first and second laws are incomparable.
In order to rectify this incompatibility, we will prove a corollary to Carnot’s theorem, asserting that the ratio of the work done in a complete cycle to the heat absorbed on expansion at constant temperature is “the same for all substances at the same temperature”. The corollary states that the ratio of the work done in a complete cycle to the heat absorbed on volume expansion at constant temperature is “the same for all substances at the same volume”. This, in effect, replaces the isotherms of the Carnot cycle by isochores without affecting the efficiency of the cycle.
In addition to the inverse temperature, another integrating factor for the quantity of heat will be shown to exist Einbinder (1948), and lead to a new adiabatic potential, which will be comparable to the entropy, and differ from it by a power of the adiabatic variable. The difference between this new potential and the entropy is that it is not a first-order homogeneous function of the volume. This will enable a comparison of means involving processes involving work where, otherwise, the first and second laws would be mute since both the internal energy and the entropy are first-order homogeneous functions of the volume for an IGG.
The existence of a new adiabatic potential raises the question as to the actual content and predictive power of the second law. The irresistible increase in the entropy during an irreversible process will now be confronted with the same inevitable decrease in the new adiabatic potential. Irreversibility will have to be detached from quasi-static processes which occur as a passage through a sequence of equilibrium states Born (1949). Irreversiblity will apply to those processes where the initial and final states differ either with respect to their temperatures and/or volumes. An equilibration resulting in the conservation of one of the two adiabatic potentials will depend on whether the temperature ratio of the two states, or their inverse volume ratios raised to a characteristic power, is equal to the ratio of quantities of heat absorbed or rejected at these temperatures or volumes. One conserving equilibration will be necessary in order to secure a final uniform state with a common mean value. General statements can be made about such conserving equilibrations: an entropy conserving equilibration has the lowest common final mean temperature or volume, implying the maximum amount of work has been performed. In addition, thermal efficiency can never be inferior to mechanical efficiency.
In a companion paper Lavenda (2005), these statements will be translated into comparable power means, and the lack of absoluteness of these potentials as a class of equivalent means. The metrizability of such a space of equivalent means will also be studied; the notion of a metric is entirely foreign to classical thermodynamics. Probability distributions and probabilistic notions will enter naturally whenever processes inside the system occur in an uncontrollable manner, like heat transfer and deformations. Evolution criteria will be shown to eminate from the fundamental property that power means are monotonically increasing functions of their order. There will then be shown numerous mathematical inequalities, like the Tchebychef and Jensen inequalities, all which predict an increase in entropy on the average, or to the decrease on the average of the complementary adiabatic potential. These mathematical inequalities will eliminate the need to have recourse to experiment, albeit one single experiment, to determine the sign of the entropy change, or to that of the complementary adiabatic potential.
## II power laws
If the absolute temperature, $`T`$, and volume, $`V`$, are chosen as the independent variables, the integrability condition for the entropy is
$$T\left(\frac{p}{T}\right)_V=\left(\frac{E}{V}\right)_T+p,$$
(1)
where $`p`$ is the pressure, and $`E`$ is the internal energy. If the product $`pV`$ measures the absolute temperature scale then
$$\left(\frac{E}{V}\right)_T=0.$$
(2)
In other words, in order for Boyle’s law to be identical with the absolute temperature, it is necessary and sufficient that the internal energy be a *linear* function of the temperature alone Buchdahl (1966).
If the pressure were to vary as some power of the temperature, say $`T^\alpha `$, then from (1) we would have
$`\left({\displaystyle \frac{E}{V}}\right)_T`$ $`=`$ $`\left[\left({\displaystyle \frac{\mathrm{ln}p}{\mathrm{ln}T}}\right)_V1\right]p`$ (3)
$`=`$ $`(\alpha 1)p.`$
Thus, the ICG condition (2) that $`pV`$ measure the absolute temperature, separates two domains: one in which (3) is positive, $`\alpha >1`$, and the attractive nature of the gas implies that it will condense, and another region in which (3) is negative, $`\alpha <1`$, and repulsion prevails implying that the system has a zero-point energy Einbinder (1948).
Every gas comprised of mechanically noninteracting particles obeys an equation of state of the form Einbinder (1948)
$$pV=sE(V,T),$$
(4)
where $`s>0`$ is proportional to the adiabatic exponent. Introducing (4) into (1) converts it into the differential equation
$$E=T\left(\frac{E}{T}\right)_V\frac{V}{s}\left(\frac{E}{V}\right)_T.$$
(5)
Treating (5) as a Lagrange equation, the auxiliary equations are
$$\frac{dT}{T}=\frac{sdV}{V}=\frac{dE}{E}.$$
There are two independent solutions, $`TV^s=c`$, and either $`E/T=a`$, or $`EV^s=b`$, where $`a`$, $`b`$, and $`c`$ are arbitrary constants. The general solution is $`\mathrm{\Psi }_1(a,c)=\mathrm{\Psi }_1(E/T,TV^s)=0`$, i.e.,
$$E=T\psi _1\left(TV^s\right),$$
(6)
or $`\mathrm{\Psi }_2(b,c)=\mathrm{\Psi }_2(EV^s,TV^s)=0`$, i.e.,
$$E=V^s\psi _2\left(TV^s\right),$$
(7)
where $`\mathrm{\Psi }_i`$ and $`\psi _i`$ are arbitrary functions. The coefficients of $`E`$, namely, $`1/T`$ and $`V^s`$, will later be appreciated as integrating factors for the quantity of heat.
The functions $`\psi _i`$ are solutions to the adiabatic equation
$$T\left(\frac{\psi _i}{T}\right)_V\frac{V}{s}\left(\frac{\psi _i}{V}\right)_T=0.$$
(8)
The auxiliary equations to (8) are
$$\frac{dT}{T}=\frac{sdV}{V}=\frac{d\psi _i}{0}.$$
There are again two independent solutions, $`\psi _i=a`$ and $`TV^s=b`$, so that the general solution to (8) is $`\psi _i=\psi _i(z)`$, in which $`T`$ and $`V^s`$ appear only through the combination $`z=TV^s`$.
For $`\psi _1=\text{const}`$., (6) is the thermal equation of state for an ICG, while for $`\psi _2=\text{const}`$., (7) is the zero-point energy. Whereas $`E=aT`$ is an approximate relation, valid in the high temperature limit, so too $`E=bV^s`$ can be considered an approximate relation, this time valid in the low temperature limit Einbinder (1948). These are the extreme cases where the internal energy is a function of either the absolute temperature, in the high temperature limit, or of the volume, in the low temperature limit. The zero-point expresses the fact that particle interactions are repulsive, $`(E/V)_T<0`$, and such a system would be entirely mechanical since $`dQ=dE+pdV=0`$. However, there is another possibility in which the internal energy, as well as the entropy, tends to zero monotonically with the temperature.
The thermal equation of state (6) can be considered as an IGG, i.e., one for which the particle number, $`N=\psi _1(V,T)`$, is variable, being a function of the temperature. This will provide a very profound analogy between a two phase classical system, like a Carnot engine, that was analyzed by Clapeyron by an equation bearing his name, and an IGG which does not conserve the number of particles.
In the case that $`E`$ tends to zero monotonically with $`T`$, $`\psi _2(z)`$ may be approximated by $`L(z)=cT^\alpha V^{s\alpha }`$ at low temperatures where all we demand for the present is that $`\alpha >0`$. According to (4), the internal energy, $`E=cT^\alpha V^{(\alpha 1)s}`$, gives a pressure $`p=scT^\alpha V^{(\alpha 1)s1}`$. For dynamic stability we require
$$\left(\frac{p}{V}\right)_T=[(\alpha 1)s1]\frac{p}{V}<0.$$
However, if the internal energy is to retain its property of being a first-order homogeneous function, we must have $`(\alpha 1)s=1`$, implying
$$\left(\frac{p}{V}\right)_T=0,$$
(9)
or a phase equilibrium Einbinder (1948). Varying the volume at constant temperature leaves the vapor pressure constant by having the liquid either evaporate or condense. This is precisely the condition under which the Clapeyron equation is valid. For an IGG, the volume of the second phase is nil since there is no longer particle conservation, $`N=\psi _1(V,T)\text{const}`$. This fine balance keeps the pressure independent of the volume and the internal energy a first-order homogeneous finction
$$p=scT^{q/r}E=cT^{q/r}V.$$
(10)
In the low temperature limit $`\psi _2(z)`$ can be replaced by
$$L(z)=a+cz^{q/r},$$
(11)
where $`a>0`$. Then, since $`EV^s=L(z)`$
$$E=V^s\left(a+cz^{q/r}\right),$$
and $`(E/V)_T<0`$ if $`a`$ is finite, or $`(E/V)_T>0`$ if $`a`$ vanishes, and $`q>r`$. In the later case, the pressure must satisfy (9) so that $`s=r/(qr)`$. Furthermore, $`E/T=\psi _1(z)=L(z)/z`$, and $`d\left(EV^s\right)/z=dL(z)/z=:dS`$, where
$$S(z)=\frac{q}{qr}cz^{1/s}=\frac{q}{qr}\psi _1(z)$$
(12)
is the entropy.
Hence, no decision can be made between the Planck ($`S=0`$ at $`T=0`$) and the Nernst $`(S=\text{const.}`$ at $`T=0`$) formulations of the third law because the concept of absolute entropy is meaningless. Moreover, we will see in the following paper that an absolute $`L(z)`$ is also meaningless since only its difference is measurable.
The difference between the enthalphy,
$$H=U+pV=(1+s)cT^{q/r}V,$$
and $`T`$ times the entropy, (12), is
$$G=\left(1+s\frac{q}{qr}\right)cT^{q/r}V.$$
(13)
The Gibbs free energy (13) presents itself as a measure of non-extensivity, and vanishes when $`s=r/(qr)`$ Einbinder (1948).
## III the second laws
The quantity of heat, $`dQ`$, which must be absorbed by a body to makes its temperature rise to $`T+dT`$ and its volume expand to $`V+dV`$ is
$$dQ=MdV+NdT.$$
(14)
This caloric equation, familiar to all the early thermodynamicists, undoubtedly fell out of favor due to the fact Clausius (1850)
> that $`Q`$ cannot be a function of $`V`$ and $`T`$, if these variables are independent of each other. For if it were, then by the well-known law of the differential calculus, that if a function of two variables is differentiated with respect to both of them, the order of differentiation is indifferent
and this was definitely not so with (14). In fact, the exactness condition of the internal energy,
$$dU=dQpdV=(Mp)dV+NdT,$$
shows that their difference gave Thomson (1851)
$$\frac{dp}{dT}=\frac{M}{T}\frac{N}{V}.$$
(15)
This led Planck Planck (1945) to remark that the notation $`dQ`$
> has frequently given rise to misunderstanding, for $`dQ`$ has been repeatedly regarded as the differential of a known finite quantity $`Q`$. This faulty reasoning may be illustrated by the following example.
And the example Planck gave resulted in (15), which Thomson took merely as a statement of the first law.
The statement that he took as the second law was Clapeyron’s equation
$$\frac{dp}{dT}=\frac{M}{C(T)},$$
(16)
where $`C(T)`$ is the reciprocal of the Carnot function, which Thomson showed was equal to the absolute temperature, $`T`$. That is, the ratio of the total work done in an infinitesimal cycle, $`\frac{dp}{dT}dTdV`$, to the ratio of the heat absorbed in the first branch of the Carnot cycle, $`MdV`$,
$$\frac{\frac{dP}{dT}dT}{M}=\frac{dT}{C(T)}$$
(17)
must be the product of $`dT`$ and a function of $`T`$ only. The “very remarkable theorem that $`dp/dT/M`$ must be the same for all substances at the same temperature was first given (although not in precisely the same terms) by Carnot”Thomson (1851).
For an ICG, $`M=p`$, because the internal energy is a function of the absolute temperature alone, and $`N`$, it was realized by Clausius, “can be a function of $`T`$ only. It is even probable that this magnitude \[$`N`$\], which represents the specific heat of the gas at constant volume, is a constant.”
In constrast, for an IGG,
$$M=\left(1+\frac{1}{s}\right)cT^{q/r}=\frac{q}{r}p,$$
(18)
and
$$sN=\frac{q}{r}cT^{1/s}V.$$
(19)
Expression (18) invalidates Clausius’s conclusion that “a permanent gas, when expanded at constant temperature, takes up only so much heat as is consumed doing external work during the expansion”. Clausius’s conclusion is based on the fact that the working substance was an ICG, obeying $`pV=T`$, while (18) shows that an IGG absorbs more since $`q>r`$.
To convert $`dQ`$ into the total differential of a certain function, we introduce the integrating factor $`\lambda `$, and require
$$\frac{\lambda M}{T}=\frac{\lambda N}{V}.$$
(20)
The exactness condition (20) can be rearranged to read
$$\left(\frac{M}{T}\frac{N}{V}\right)=M\frac{\mathrm{ln}\lambda }{T}+N\frac{\mathrm{ln}\lambda }{V}.$$
(21)
By a composite system argument, $`\lambda `$ can only be a function of $`T`$ or $`V`$. For consider two simple fluids in thermal contact; such a system will have three independent variables $`T`$, $`V_1`$ and $`V_2`$. The integrating factor can only be a function of the common variable $`T`$ Carathéodory (1909). Hence,
$$\frac{\left(\frac{M}{T}\frac{N}{V}\right)}{M}=\frac{d\mathrm{ln}\lambda }{dT}.$$
On the strength of the exactness condition for the internal energy (15), this is equivalent to the Clapeyron equation (17), and on the strength of Carnot’s theorem $`C(T)=T`$ is the integrating denominator for the increment in the heat, $`dQ`$, giving the entropy, (12), which depends on $`V`$ and $`T`$ only through the combination $`z=TV^s`$.
Now consider two simple fluids in mechanical contact, for which the independent variables are $`V`$, $`T_1`$, and $`T_2`$. Again, the integrating factor $`\lambda `$ can only be a function of the variable in common to both subsystems so that the exactness condition (21) now reduces to
$$\frac{\frac{dp}{dT}dTdV}{NdT}=d\mathrm{ln}\lambda (V).$$
(22)
By interchanging isochores for isotherms, $`NdT`$ is the quantity of heat absorbed in the first segment of an equivalent Carnot cycle, as we shall discuss in the next section. The very remarkable fact, equivalent to Carnot’s theorem, is that the right-hand side is the product of $`dV`$ and a function only of the volume. Thus, *$`dp/dT/N`$ is the same for all substances at the same volume*. For an IGG, $`\lambda (V)=V^s`$ Einbinder (1948), which reduces to $`\lambda (V)=V^{R/C_v}`$ for an ICG, where $`C_v`$ is the heat capacity at constant volume, and $`R`$ is the gas constant, which will only be introduced in conjunction with $`C_v`$.
The integrating factor $`\lambda (V)`$ for the increment in the heat, $`dQ`$, now gives the point function (11), which, again, depends on $`V`$ and $`T`$ only through the combination, $`z`$. However, whereas the entropy (12) is extensive, the potential (11) is not.
Consider any two states $`1`$ and $`2`$, with $`T_1>T_2`$ for concreteness. Any process connecting the two states will be said to be irreversible if
$$z_1>z_2.$$
(23)
Rearranging (23) we get $`V_1^s/V_2^s>T_2/T_1`$ so that the thermal efficiency,
$$\eta _t=1\frac{T_2}{T_1}1\frac{V_1^s}{V_2^s}=\eta _v,$$
can never be inferior to the mechanical efficiency, $`\eta _v`$.
If $`Q_1`$ and $`Q_2`$ are quantities of heat absorbed and rejected at $`T_1`$ and $`T_2`$, respectively, then we say that there is *thermal equilibration* if
$$\frac{T_2}{T_1}=\frac{Q_2}{Q_1}\frac{V_1^s}{V_2^s},$$
whereas there is *mechanical equilibration* if
$$\frac{V_1^s}{V_2^s}=\frac{Q_2}{Q_1}\frac{T_2}{T_1}.$$
More work can be accomplished during a thermal equilibration than a mechanical one since the system achieves a lower final temperature, $`T_2`$.
For any infinitesimal narrow Carnot cycle, the ratio $`Q_2/Q_1`$ may be replaced by its differential, $`dQ_2/dQ_1`$. Thermal equilibration results when
$$\frac{dQ}{T}=0,$$
(24)
and
$$V^s𝑑Q0,$$
(25)
while mechanical equilibration requires
$$V^s𝑑Q=0,$$
(26)
and
$$\frac{dQ}{T}0,$$
(27)
since any irreversible cycle is necessarily less efficient than a Carnot cycle, $`dQ_2/dQ_1<T_2/T_1`$.
For processes of pure thermal conduction, we would identify (26) with the first law, and the conservation of energy over the entire cycle, and (27) as the statement of the second law. However, for purely mechanical interactions neither (27) nor the first law would give an evolutionary criterion since the internal energy and entropy are first-order homogeneous functions of the volume. The criteria of mechanical equilibration would fix the final volume as the weighted arithmetic mean of the partial volumes, according to (25), while (24) would show that $`L`$ would tend to decrease since the arithmetic mean is inferior to the power mean of order $`q/(qr)`$.
If we split the cycle up into two segments, $`AB`$, and $`BA`$, where the former contains all the irreversibility, then (25) gives
$$_A^BV^s𝑑Q>L_BL_A,$$
while (27) becomes
$$S_BS_A>_A^B\frac{dQ}{T}.$$
If the infinitesimal increment in the heat is the sum of the infinitesimal heat $`dQ_e`$ introduced into the system, or extracted from it, and the sum of irreversible heat transfers within the system,
$$dQ=dQ_e+\underset{i=1}{\overset{n}{}}dQ_i,$$
then for an isolated system $`dQ_e=0`$,
$$L_BL_A<\underset{i=1}{\overset{n}{}}_A^BV^s𝑑Q_i,$$
(28)
under the condition that $`_{i=1}^n𝑑Q_i/T=0`$, or
$$S_BS_A>\underset{i=1}{\overset{n}{}}_A^B\frac{dQ_i}{T},$$
(29)
under the condition $`_{i=1}^nV^s𝑑Q_i=0`$. Inequality (29) is referred to as Clausius’s inequality, while inequality (28) appears to be novel.
For each individual process of heat transfer, $`dQ_i`$ will appear twice: once as a positive quantity and once as a negative quantity. The second law (29) asserts that the integrating denominator in the former is smaller than in the latter, since ‘heat flows spontaneously from a hotter to a colder body at constant volume.’ Alternatively, according to (28), ‘as heat is absorbed the gas expands at constant temperature’, so that the integrating factor of the former will be smaller than the latter.
Either (28) or (29) can be taken as statements of the second law, depending on the constraints imposed. It is quite remarkable that both criteria can be formulated in terms of *comparable* means, and their fundamental property that the power mean is a monotonic increasing function of its order will be investigated thoroughly in the following paper.
## IV An equivalent Carnot cycle
An equivalent cycle to that of Carnot can be obtained by replacing isothermals by isochores. This will prove a very remarkable theorem that $`dp/dT/N`$ must be the same for all substances at the same volume, which is precisely analogous to Carnot’s theorem that $`dp/dT/M`$ must be the same for all substances at the same temperature.
The cycle consists of:
1. Absorption of a quantity of heat $`Q_1`$, at a constant volume $`V_1`$, by compression which raises the absolute temperature from $`T_1`$ to $`T_2`$, and, consequently, increases the pressure.
2. An adiabatic expansion to a state of larger volume $`V_2`$, and lower temperature $`T_3`$.
3. The rejection of a quantity of heat $`Q_2`$, by expansion which lowers the temperature to $`T_4`$ at constant volume $`V_2`$.
4. An adiabatic compression which restores the system to volume $`V_1`$, and temperature $`T_1`$.
The adiabatic branches provide the following ratios
$$\frac{V_1^s}{V_2^s}=\frac{T_3}{T_2}=\frac{T_4}{T_1}.$$
The ratio of the heat rejected, $`|Q_{34}|`$, to the heat absorbed, $`Q_{12}`$ is
$$\frac{|Q_{34}|}{Q_{12}}=\frac{V_1^s}{V_2^s},$$
and the efficiency of the engine
$$\eta _v=1\frac{V_1^s}{V_2^s},$$
is the same as the Carnot, thermal, efficiency, $`\eta _t`$.
Now, for an infinitesimal cycle, the ratio of the total work performed to the heat absorbed at constant volume is (22), which is the same for all substances at the same volume. This, as we have shown in the last section, is equivalent to the Carnot theorem (17).
Both (22) and (17) gives the same Clapeyron equation, which, when integrated, gives the pressure in (10), independent of volume for a homogeneous system. And while it was realized that $`dp/dT`$ was a function of the temperature alone Clausius (1850), it was not appreciated that the mechanical ICG equation of state, $`pV=T`$, could not be used to evaluate the work.
## V The transition IGG$``$ICG
The absolute temperature, $`T`$, and the empirical temperature, $`t`$, will coincide only for an ICG. For an ICG, $`pV`$ reads the temperature, and $`E`$ is a function of the temperature alone, independent of the volume. In contrast for an IGG, $`p`$ will be independent of the volume, and $`E`$ will be a function of it. In addition, $`E`$ will no longer be a linear function of the temperature.
The only demand made by the zeroth law is that when two identical systems are placed in thermal contact their empirical temperatures be the same when a state of mutual thermal equilibrium has been reached. Once the empirical scale has been chosen, the absolute temperature must be a monotonically increasing function of it, viz.,
$$T(t)=t^r,$$
with $`r1`$.
On the empirical temperature scale, $`E`$ will vary as $`t^q`$, where $`qr`$ with the equality sign pertaining to the ICG limit. Whereas the exponent $`r`$ is related to the average kinetic energy of the particles, the exponent $`q`$ is related to additional forms of energy, like the energy required to create or annihilate particles, since for an IGG the particle number is not conserved. As a matter of fact, the vanishing of the difference of the two exponents will signal the transition from an IGG to an ICG, as we shall now show. This is in distinction to the transition between an IQG and an ICG which occurs only in the high temperature limit.
It is notable that in the limit as $`qr`$, the entropy (12) is of the indeterminate form $`0/0`$. Applying L’Hôpital’s rule, we get
$$\underset{qr}{lim}S(z)=c\mathrm{ln}z,$$
where $`z=TV^{R/C_v}`$. Identifying $`c`$ as $`C_v`$ gives
$$\underset{qr}{lim}S(z)=C_v\mathrm{ln}T+R\mathrm{ln}V.$$
In the same limit,
$$\underset{qr}{lim}L(z)=C_vT^{R/C_v},$$
and consequently, either
$$dE=TdS(z)pdV=C_vdT$$
or
$$dE=\frac{1}{V^{R/C_v}}dL(z)pdV=C_vdT,$$
which is the thermal equation of state of an ICG. |
warning/0507/hep-th0507086.html | ar5iv | text | # Gauge-invariant perturbation theory for trans-Planckian inflation
## I Introduction
Cosmological inflation is currently considered to be the best paradigm for describing the early stages of the universe Linde (1990); Liddle and Lyth (2000). Inflation leads to the existence of a causal mechanism for producing fluctuations on cosmological scales (when measured today), which at the time of matter-radiation equality had a physical wavelength larger than the Hubble radius. Thus, it solves several conceptual problems of standard cosmology and leads to a predictive theory of the origin of cosmological fluctuations.
During inflation, the physical wavelength corresponding to a fixed comoving scale decreases exponentially as time decreases whereas the Hubble radius is constant. Thus, as long as the period of inflation is sufficiently long, all scales of cosmological interest today originate inside the Hubble radius during inflation. Recently, it has been realized that if inflation lasts slightly longer than the minimal time (i. e. the time it needs to last in order to solve the horizon problem and to provide a causal generation mechanism for CMB fluctuations), then the corresponding physical wavelength of these fluctuations at the beginning of inflation will be smaller than the Planck length. This is commonly referred to as the trans-Planckian problem of inflation Brandenberger and Martin (2001); Martin and Brandenberger (2001); Niemeyer (2001).
Naturally, considerable amount of attention has been devoted to examine the possibility of detecting trans-Planckian imprints on the CMB Tanaka (2000); Niemeyer and Parentani (2001); Kempf and Niemeyer (2001); Starobinsky (2001); Easther et al. (2001); Lemoine et al. (2002); Bastero-Gil et al. (2002); Easther et al. (2003); Shankaranarayanan (2003); Danielsson (2002); Brandenberger and Ho (2002); Hassan and Sloth (2003); Easther et al. (2002); Kaloper et al. (2002); Martin and Brandenberger (2003); Bastero-Gil et al. (2003); Shankaranarayanan and Sriramkumar (2004a, b); Brandenberger and Martin (2005); Tsujikawa et al. (2003); Sriramkumar and Padmanabhan (2005); Greene et al. (2005); Calcagni (2004). Broadly, there have been two approaches in the literature in order to study these effects. In the first approach, the specific nature of trans-Planckian physics is not presumed, but is rather described by the boundary conditions imposed on the mode at the cut-off scale. In the second approach, which is of interest in this work, one incorporates quantum gravitational effects by introducing the fundamental length scale into the standard field theory in a particular fashion.
In almost all the analyses performed so-far in the literature, the trans-Planckian effects are introduced into the scalar and tensor perturbation equations in an ad hoc manner. In the standard inflation, the scalar and tensor perturbation equations are derived, from the first principles, in a gauge-invariant manner Kodama and Sasaki (1984); Mukhanov et al. (1992). However in the trans-Planckian inflationary scenario, to our knowledge, such a calculation has never been performed and the scalar/tensor power spectrum, with the trans-Planckian corrections, were obtained from the presumed perturbation equations.
With the possibility of detecting the trans-Planckian signatures in the current and future CMB experiments<sup>1</sup><sup>1</sup>1For the current status and developments of WMAP and PLANCK, see the following URLs: http://map.gsfc.nasa.gov/, http://astro.estec.esa.nl/SA-general/Projects/Planck, it becomes imperative to obtain the scalar and tensor perturbation equation from the first principles. (For recent work on the trans-Planckian constraints from the CMB, see Refs. Martin and Ringeval (2004a, b); Martin and Ringeval (2005); Easther et al. (2005a, b).) In this work, we perform the gauge-invariant cosmological perturbation theory for the single-scalar field inflation with the trans-Planckian corrections. The model we shall consider in this work is a self-interacting scalar field in (3 + 1)-dimensional space-time satisfying a linear wave equation with higher spatial derivative terms. The dispersion relation \[$`\omega =\omega (k)`$\] thus differs at high wave-vectors from that of the ordinary wave equation. Such a model breaks the local Lorentz invariance explicitly while preserving the rotational and translational invariance. The particular dispersion relation we shall study in detail is
$$\omega ^2(k)=|\stackrel{}{k}|^2+b_{11}|\stackrel{}{k}|^4,$$
(1)
where $`b_{11}`$ is a dimensionfull parameter. $`b_{11}<0`$ implies subluminal group velocity, while $`b_{11}>0`$ implies superluminal group velocity. The above dispersion relation is a subset of a general class of the form $`\omega ^2=|\stackrel{}{k}|^2[1+g(|\stackrel{}{k}|/k_0)]`$, where $`g`$ is a function which vanishes as $`k_0\mathrm{}`$, and $`k_0`$ is a constant which sets the scale for the deviation from Lorentz invariance. It has been suggested that these general modified dispersion relation might arise in loop quantum gravity Gambini and Pullin (1999); Alfaro et al. (2000), or more generally from an unspecified modification of the short distance structure of space-time (see for example Refs. Amelino-Camelia (2000); Jacobson (1999)). Possible observational consequences have also been studied. For an up-to-date review, see Refs. Mattingly (2005); Jacobson et al. (2005).
Even though the above dispersion relation breaks the local Lorentz invariance explicitly it has been shown that it is possible to write a Lagrangian for the above field in a generally covariant manner, consistent with spatial translation and rotation invariance, by introducing a unit time-like Killing vector field which defines a particular direction Jacobson and Mattingly (2001a); Lemoine et al. (2002); Lim (2005) (see Sec. (III) for more details). In this work, we use such a framework in-order to obtain the scalar/tensor perturbation equations for such a model during inflation.
Using the covariant Lagrangian used in Ref. Jacobson and Mattingly (2001a), we obtain the perturbed stress-tensor for the scalar and tensor perturbations about the FRW background. We show that: (i) The non-linear effects introduce corrections to the perturbed energy density while the other components of the stress-tensor remains unchanged. (ii) The non-linear terms contributing to the stress-tensor are proportional to $`k^2`$. Hence in the super-Hubble scales the trans-Planckian contributions to the perturbed energy density, as expected, can be ignored. (iii) The spatial higher derivative terms appear only in the equation of the motion of the perturbed inflation field ($`\delta \phi `$) and not in the equation of motion of the scalar perturbations ($`\mathrm{\Phi }`$). (iv) Unlike the canonical scalar field inflation, the perturbations, in general, are not purely adiabatic. The entropic perturbations generated during the inflation, however, vanish at the super-Hubble scales. The speed of propagation of the perturbations is a constant and is less than the speed of light. (v) The tensor perturbation equation remain unchanged indicating that the well-know consistency relation between the scalar and tensor ratio will also be broken in this model.
We obtain the equation of motion of the Mukhanov-Sasaki variable corresponding to the inflaton field. We show that the equation of motion derived from the gauge-invariant perturbation theory is not same as those assumed in the earlier analyzes Brandenberger and Martin (2001); Martin and Brandenberger (2001); Niemeyer (2001). Later, we combine the system of differential equations into a single differential equation in $`\mathrm{\Phi }`$ and obtain the solutions for the power-law inflation in different regimes. We also obtain the spectrum of scalar perturbations in a particular limit and compare with the earlier results.
This paper is organized as follows: In the following section, the theory of cosmological perturbations for the canonical single scalar field inflation is discussed and essential steps leading to the perturbation equation are reviewed. In Sec. (III), we discuss the general covariant formulation of the Lagrangian describing the scalar field with modified dispersion relation and derive the corresponding stress-tensor. In Sec. (IV), we obtain the perturbed stress-tensor for the scalar and tensor perturbations. In Sec. (V), we obtain the scalar perturbation equation and the equation of motion of the Mukhanov-Sasaki variable. In Sec. (VI), we perform the classical analysis and obtain the form of the scalar perturbations ($`\mathrm{\Phi }`$) in the various regimes. In Sec. (VII), we solve the perturbation equations in a particular limit and obtain the power-spectrum of the perturbations. Our results are summarized and discussed in the last section. In Appendices (A, B) we derive the equations of motion of the fields in the FRW and perturbed FRW backgrounds. In Appendix (C), we obtain the equation of motion of the Bardeen potential in our model.
Through out this paper, the metric signature we adopt is $`(+,,,)`$ Landau and Lifshitz (1975), we set $`\mathrm{}=c=1`$ and $`1/(8\pi G)=M_{_{\mathrm{Pl}}}^2`$. The various physical quantities with the over-line refers to the values evaluated for the homogeneous and isotropic FRW background. A dot denotes derivative with respect to the cosmic time ($`t`$), a prime stands for a derivative with respect to conformal time ($`\eta `$) and <sub>,i</sub> denotes a derivative w.r.t space components. We follow the notation of Ref. Mukhanov et al. (1992) to provide easy comparison.
## II Gauge-invariant perturbation: Canonical single field inflation
In this section, we obtain the scalar and tensor perturbation equations for the canonical single scalar field inflation. In the following subsection, we discuss key properties of the perturbed FRW metric and the “gauge problem” of the scalar perturbations. In the subsequent subsections, we discuss the matter Lagrangian and provide key steps in obtaining the scalar/tensor perturbation equations.
### II.1 Perturbed FRW metric
We consider perturbations about a spatially flat $`(3+1)`$-dimensional FRW line element
$`ds^2`$ $`=`$ $`\overline{g}_{\mu \nu }dx^\mu dx^\nu `$ (2)
$`=`$ $`dt^2a^2(t)d𝐱^2=a^2(\eta )\left(d\eta ^2d𝐱^2\right),`$
where $`t`$ is the cosmic time, $`a(t)`$ is the scale factor and $`\eta =\left[dt/a(t)\right]`$ denotes the conformal time. $`H`$ is the Hubble parameter given by $`H\dot{a}/a`$ while $`a^{}/a`$ is related to the Hubble parameter by the relation $`=Ha`$.
At the linear level, for the canonical single scalar field inflation, the metric perturbations ($`\delta g_{\mu \nu }`$) can be categorized into two distinct types — scalar and tensor perturbations. Thus, the perturbed FRW line-element can be written as
$`\mathrm{d}s^2`$ $`=`$ $`a^2(\eta )[(1+2\varphi )\mathrm{d}\eta ^22_iB\mathrm{d}x^i\mathrm{d}\eta `$
$`[(12\psi )\delta _{ij}+2_i_jE+h_{ij}]\mathrm{d}x^i\mathrm{d}x^j],`$
where the functions $`\varphi `$, $`B`$, $`\psi `$ and $`E`$ represent the scalar sector whereas the tensor $`h_{ij}`$, satisfying $`h_i^i=^ih_{ij}=0`$, represent gravitational waves. Note that all these first-order perturbations are functions of $`(\eta ,𝐱)`$. For convenience, we do not write the dependence explicitly.
The tensor perturbations do not couple to the energy density ($`\delta \rho `$) and pressure ($`\delta p`$) inhomogeneities. However, the scalar perturbations couple to the energy density and pressure which lead to the growing inhomogeneities. At the linear level, the two types of perturbations decouple and can be treated separately.
The scalar and tensor perturbations have four and two degrees of freedom respectively. In the case of tensor perturbations, the two degrees of freedom correspond to the two polarizations of the gravitational waves and hence are physical. The scalar perturbations suffer from the gauge problem. (For a detailed discussion, see Refs. Kodama and Sasaki (1984); Mukhanov et al. (1992).) However, it is possible to construct two gauge invariant variables, which characterize the perturbations completely, from the metric variables alone i. e.,
$$\mathrm{\Phi }\varphi +\frac{1}{a}[(BE^{})a]^{},\mathrm{\Psi }\psi (BE^{}).$$
(4)
Physically, $`\mathrm{\Phi }`$ corresponds to the Newtonian gravitational potential and is commonly referred to as the Bardeen potential while $`\mathrm{\Psi }`$ is related to the perturbations of the $`3`$-space. For the single canonical scalar field scenario, we have $`\mathrm{\Phi }=\mathrm{\Psi }`$.
$`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ are related to the pressure and density perturbations of a generic perfect fluid via the perturbed Einstein’s equations. The pressure perturbations, in general, can be split into adiabatic and entropic (non-adiabatic) parts, by writing
$$\delta p=c_\mathrm{s}^2\delta \rho +\overline{p}{}_{}{}^{}\mathrm{\Gamma },$$
(5)
where $`c_s^2\overline{p}{}_{}{}^{}/\overline{\rho }^{}`$ is the adiabatic sound speed Wands et al. (2000); Malik et al. (2003); Malik and Wands (2005). The non-adiabatic part is $`\delta p_{\mathrm{nad}}\overline{p}{}_{}{}^{}\mathrm{\Gamma }`$, and
$$\mathrm{\Gamma }\frac{\delta p}{\overline{p}^{}}\frac{\delta \rho }{\overline{\rho }^{}}.$$
(6)
The entropic perturbation $`\mathrm{\Gamma }`$, defined in this way, is gauge-invariant, and represents the displacement between hyper-surfaces of uniform pressure and uniform density. In the context of canonical single scalar field inflation, only adiabatic perturbations are present, i. e., $`\delta p=c_s^2\delta \rho `$ (where $`c_s^2=1`$) or $`\mathrm{\Gamma }=0`$. However, as we shall see in the later sections, the trans-Planckian inflationary scenario introduces entropic perturbations as well.
### II.2 Canonical scalar field
The dominant matter component during inflation is a spatially homogeneous canonical scalar field $`\overline{\phi }(\eta )`$ (inflaton). The Lagrangian density for the canonical scalar field $`\phi (\eta ,\overline{x})`$ propagating in a general curved background is given by
$$_\phi =\frac{1}{2}g^{\mu \nu }_\mu \phi _\nu \phi V(\phi ),$$
(7)
where $`V(\phi )`$ is the self-interacting scalar field potential. The equation of motion and the stress tensor of the scalar field ($`\phi `$) in the conformally flat FRW background (2) are given by
$`\phi ^{\prime \prime }+2\phi ^{}^2\phi +a^2V_{,\phi }(\phi )=0,`$ (8)
$`T_{\nu }^{\mu }{}_{}{}^{(\phi )}=^\mu \phi _\nu \phi \left[{\displaystyle \frac{1}{2}}^\alpha \phi _\beta \phi V(\phi )\right]\delta _\nu ^\mu ,`$ (9)
where $`V_{,\phi }=\left(dV(\phi )/d\phi \right)`$ and $`^2`$ refers to the Laplacian in the flat space.
Let us consider a small inhomogeneous quantum fluctuations on top of a homogeneous and isotropic classical background. For the scalar field, we have
$$\phi (\eta ,𝐱)=\overline{\phi }(\eta )+\delta \phi (\eta ,𝐱),$$
(10)
where one assumes that the perturbation $`\delta \phi `$ is small. The perturbed scalar field ($`\delta \phi `$) and the perturbed stress-tensor of the scalar field ($`\delta T_{\nu }^{\mu }{}_{}{}^{(\phi )}`$), like the other scalar-type perturbation functions $`(\varphi ,B,\psi ,E)`$, suffer from the gauge problem. (For a detailed discussion, see Refs. Kodama and Sasaki (1984); Mukhanov et al. (1992).) Similar to Eq. (4), it is possible to define a gauge-invariant quantity for the perturbed scalar field and the perturbed stress-tensor, i. e.,
$`\delta \phi ^{\left(gi\right)}\delta \phi +\overline{\phi }^{}\left(BE^{}\right),`$ (11)
$`\delta T_{0}^{0}{}_{}{}^{\left(\mathrm{gi}\right)}\delta T_0^0+\overline{T_0^0}^{}\left(BE^{}\right),`$
$`\delta T_{i}^{j}{}_{}{}^{\left(\mathrm{gi}\right)}\delta T_i^j+\overline{T_i^j}^{}\left(BE^{}\right),`$ (12)
$`\delta T_{0}^{i}{}_{}{}^{\left(\mathrm{gi}\right)}\delta T_0^j+\left(\overline{T_0^0}{\displaystyle \frac{1}{3}}\overline{T_i^j}\right)\left(BE^{}\right)_{,i}.`$
Separating the homogeneous and perturbed part from Eq. (8), we have
$`\overline{\phi }^{^{\prime \prime }}+2\overline{\phi }^{^{}}+a^2V_{,\phi }=0,`$ (13)
$`\delta \phi _{}^{(gi)}{}_{}{}^{\prime \prime }+2\delta \phi _{}^{(gi)}{}_{}{}^{}^2\left(\delta \phi ^{(gi)}\right)`$ (14)
$`+V_{,\phi \phi }a^2\delta \phi ^{(gi)}4\overline{\phi }^{}\mathrm{\Phi }^{}+2V_{,\phi }a^2\mathrm{\Phi }=0.`$
Similarly, separating the homogeneous and perturbed part in the stress-tensor (9), we get
$`\overline{T_{0}^{0}{}_{}{}^{\left(\phi \right)}}={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\overline{\phi }^2}{a^2}}+V\left(\overline{\phi }\right)\right);\overline{T_{j}^{i}{}_{}{}^{\left(\phi \right)}}={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\overline{\phi }^2}{a^2}}V\left(\overline{\phi }\right)\right)\delta _j^i`$
$`{}_{}{}^{\left(gi\right)}\delta T_{0}^{0}{}_{}{}^{\left(\phi \right)}=a^2\left[\overline{\phi ^{}}^2\mathrm{\Phi }+\overline{\phi ^{}}\delta \phi _{}^{\left(gi\right)}{}_{}{}^{}+V_{,\phi }a^2\delta \phi ^{\left(gi\right)}\right],`$
$`{}_{}{}^{\left(gi\right)}\delta T_{i}^{0}{}_{}{}^{\left(\phi \right)}=a^2\overline{\phi ^{}}\delta \phi _{,i}^{\left(gi\right)},`$
$`{}_{}{}^{\left(gi\right)}\delta T_{j}^{i}{}_{}{}^{\left(\phi \right)}=a^2\left[\overline{\phi ^{}}^2\mathrm{\Phi }\overline{\phi ^{}}\delta \phi _{}^{\left(gi\right)}{}_{}{}^{}+V_{,\phi }a^2\delta \phi ^{\left(gi\right)}\right]\delta __j^i.`$
### II.3 Scalar and Tensor Perturbation equations
In the earlier subsections, we obtained the gauge invariant variables related to the scalar field $`(\delta \phi ^{(gi)},^{(gi)}\delta T_{\nu }^{\mu }{}_{}{}^{(\phi )})`$ and metric perturbations $`(\mathrm{\Phi },\mathrm{\Psi })`$. In this subsection, we outline the essential steps leading to the scalar and tensor equations of motion. Even though, this is a standard result, and can be found in numerous review articles (see, for example, Refs. Kodama and Sasaki (1984); Mukhanov et al. (1992)), we have given here the key steps for future reference.
From Eq. (II.2), it is easy to see that the non-diagonal space-space components of the stress-tensor are absent. This leads to the condition that $`\mathrm{\Phi }=\mathrm{\Psi }`$. Thus, the Einstein’s equations for the perturbed FRW metric (II.1) in terms of the gauge invariant quantities are:
$`^2\mathrm{\Phi }3\mathrm{\Phi }^{}3^2\mathrm{\Phi }={\displaystyle \frac{1}{2M_{_{\mathrm{Pl}}}^2}}a^2\delta T_{0}^{0}{}_{}{}^{(\mathrm{gi})}`$ (17a)
$`{\displaystyle \frac{(a\mathrm{\Phi })_{,i}^{}}{a}}={\displaystyle \frac{1}{2M_{_{\mathrm{Pl}}}^2}}a^2\delta T_{i}^{0}{}_{}{}^{(\mathrm{gi})}`$ (17b)
$`\mathrm{\Phi }^{\prime \prime }+3\mathrm{\Phi }^{}+(2^{}+^2)\mathrm{\Phi }={\displaystyle \frac{1}{2M_{_{\mathrm{Pl}}}^2}}a^2\delta T_{i}^{i}{}_{}{}^{(\mathrm{gi})},`$ (17c)
where $`\delta T_{\nu }^{\mu }{}_{}{}^{(\mathrm{gi})}`$ is given by Eq. (II.2). The three perturbed Einstein’s equations can be combined to form a single differential equation in $`\mathrm{\Phi }`$:
$$\mathrm{\Phi }^{^{\prime \prime }}^2\mathrm{\Phi }+2\left(\frac{\overline{\phi }^{^{\prime \prime }}}{\overline{\phi }^{^{}}}\right)\mathrm{\Phi }^{^{}}+2\left(^{^{}}\frac{\overline{\phi }^{^{\prime \prime }}}{\overline{\phi }^{^{}}}\right)\mathrm{\Phi }=0.$$
(18)
The system of perturbation equations (18, 17b, 14) is quite complex. In order to extract the physical content more transparent, these equations are expressed in terms of two new variables – $`Q`$ and $`u`$ – which are linearly related to $`\mathrm{\Phi }`$ and $`\delta \phi ^{(gi)}`$, i. e.,
$$Q=a\left[\delta \phi +\overline{\phi }^{}\frac{\psi }{}\right];u=\frac{a}{\overline{\phi }^{}}\mathrm{\Phi }.$$
(19)
$`Q`$ is a gauge-invariant (Mukhanov-Sasaki) Kodama and Sasaki (1984); Mukhanov et al. (1992) variable whose equation of motion is homogeneous. This ensures that one can quantize $`Q`$ in the standard way using the Lagrangian associated with its equation of motion. At the early stages of inflation where the quantum effects are important, the equation of motion of $`Q`$ helps in quantizing the fields and fixing the initial conditions. However, at the end stages of inflation where the relevant modes have crossed the Hubble radius and behave classically, it is easier to analyze the equation of motion of $`u`$.
The equation of motion of $`Q`$ is derived as follows: Substituting $`\delta \phi `$ in terms of $`Q`$ in Eq. (14), and using the relations (18, 13, 17b), we get:
$$Q^{\prime \prime }^2Q\frac{z^{\prime \prime }}{z}Q=0,$$
(20)
where
$$z=\frac{a(^2^{})^{1/2}}{c_s}=a\frac{\overline{\phi }^{}}{}.$$
(21)
The equation of motion of $`u`$ is obtained by substituting the transformation (19) in Eq. (18):
$$u^{\prime \prime }^2u\frac{\theta ^{\prime \prime }}{\theta }u=0,$$
(22)
where
$$\theta =\frac{}{a}\left[\frac{2}{3}(^2^{})\right]^{1/2}=\frac{}{a\overline{\phi }^{}}.$$
(23)
$`Q`$ is related to another gauge-invariant quantity $`\zeta ((/\rho ^{^{}})\delta \rho +\psi )`$ by the relation $`Q=2c_sz\zeta `$. The quantity $`\zeta `$ is time-independent on scales larger than the Hubble radius and, more importantly, related to the large scale CMB anisotropies (via the Sachs-Wolfe effect) Liddle and Lyth (2000).
Decomposing $`Q`$ into Fourier space, we have
$$\mu __S^{\prime \prime }+\left[k^2\frac{z^{\prime \prime }}{z}\right]\mu __S=0,$$
(24)
where $`k=|𝐤|`$ and $`\mu __S=Q_k=2c_sz\zeta `$. The above equation is similar to a time-independent Schrödinger equation where the usual role of the radial coordinate is now played by the conformal time whose effective potential is $`U__\mathrm{S}z^{\prime \prime }/z`$ Wang et al. (1997); Martin and Schwarz (2003). The scalar perturbation spectrum per logarithmic interval can then be written in terms of the modes $`\mu __S`$ as
$$\left[k^3𝒫_S(k)\right]=\left(\frac{k^3}{2\pi ^2}\right)\left(\frac{|\mu __S|}{z}\right)^2,$$
(25)
and the expression on the right hand side is to be evaluated when the physical wavelength $`(k/a)^1`$ of the mode corresponding to the wavenumber $`𝐤`$ equals the Hubble radius $`H^1`$.
Before proceeding to the next section, we obtain the tensor perturbation equation in the FRW background. As we mentioned earlier, the tensor perturbations $`h_{ij}(\eta ,𝐱)`$ do not couple to the energy density. These represent free gravitational waves and satisfy the equation:
$$\mu __T^{^{\prime \prime }}+\left(k^2\frac{a^{\prime \prime }}{a}\right)\mu __T=\mathrm{\hspace{0.17em}0}.$$
(26)
where $`\mu _Tah_k`$. This equation is very similar to the corresponding equation (24) for scalar gravitational inhomogeneities, except that in the effective potential $`(U__\mathrm{T}a^{\prime \prime }/a)`$ the scale factor $`a(\eta )`$ is replaced by $`z(\eta )`$.
## III Modified Dispersion relation Lagrangian
In this section, we briefly discuss the general covariant formulation describing a scalar field with modified dispersion relation and derive the corresponding stress-tensor.
As discussed in the introduction, to keep the calculations tractable, we will assume that the scalar field with the high frequency dispersion relation is of the form
$$\omega ^2=|\stackrel{}{k}|^2+b_{_{11}}|\stackrel{}{k}|^4,$$
(27)
where $`b_{_{11}}>0`$. The above dispersion relation breaks the local Lorentz invariance explicitly while it preserves rotational and translational invariance. Even though, the modified dispersion relation breaks the local Lorentz invariance explicitly, it was shown in Ref. Jacobson and Mattingly (2001b) that a covariant formulation of the corresponding theory can be carried out by introducing a unit time-like vector field $`u^\mu `$ which defines a preferred rest frame.
### III.1 Covariant Lagrangian
The action for a scalar field with the modified dispersion relation takes the form Jacobson and Mattingly (2001b); Lemoine et al. (2002)
$`S`$ $`=`$ $`{\displaystyle \mathrm{d}^4x\sqrt{g}(_\phi +_{_{\mathrm{cor}}}+_u)},`$ (28)
where $`_\phi `$ is the standard Lagrangian of a minimally coupled scalar field given by Eq. (7). The last two terms — $`_{_{\mathrm{cor}}}`$ and $`_u`$ — contribute to the modified dispersion relation of the scalar field. $`_{_{\mathrm{cor}}}`$ corresponds to the non-linear part of the dispersion relation while $`_u`$ describes the dynamics of the vector field $`u^\mu `$. The two corrective Lagrangians have the form
$`_{_{\mathrm{cor}}}`$ $`=`$ $`b_{11}\left(𝒟^2\phi \right)^2,`$ (29a)
$`_u`$ $`=`$ $`\lambda (g^{\mu \nu }u_\mu u_\nu 1)d_1F^{\mu \nu }F_{\mu \nu },`$ (29b)
where
$`F_{\mu \nu }`$ $``$ $`_\mu u_\nu _\nu u_\mu ,`$ (30)
$`𝒟^2\phi `$ $`=`$ $`^{\alpha \beta }_\alpha _\beta \phi +u^\alpha _\alpha \phi _\beta u^\beta ,`$ (31)
$`_{\mu \nu }`$ $``$ $`g_{\mu \nu }+u_\mu u_\nu .`$ (32)
The covariant derivative associated with the metric $`g_{\mu \nu }`$ is $`_\mu `$ while $`b_{11}`$ and $`d_1`$ are arbitrary (dimensional) constants. The tensor $`_{\mu \nu }`$ gives the metric on a slice of fixed time while $`𝒟^2`$ is proportional to the Laplacian operator on the same surface. The fact that $`u^\mu `$ is a unit time-like vector ($`u^\mu u_\mu =1`$) is enforced by the Lagrange multiplier $`\lambda `$. In the above, $`b_{11},d_1`$ have the dimensions of inverse mass square and mass square, respectively while $`u_\mu `$ is dimensionless.
The equation of motion for $`\phi `$ and $`u_\mu `$ obtained by varying the action (28), respectively, are
$`^\mu _\mu \phi +V_{,\phi }`$ $`=`$ $`2b_{_{11}}[_\mu \left(𝒟^2\phi u^\mu _\nu u^\nu \right)`$
$`_\mu _\nu \left(𝒟^2\phi ^{\mu \nu }\right)],`$
$`2d_1_\nu F^{\nu \mu }\lambda u^\mu `$ $`=`$ $`b_{_{11}}[^\mu \left(𝒟^2\phi u^\nu _\nu \phi \right)`$
$``$ $`^\mu _\nu \phi u^\nu 𝒟^2\phi _\nu \left(^\mu \phi u^\nu \right)𝒟^2\phi ].`$
From the above equation, we get
$`\lambda `$ $`=`$ $`b_{_{11}}u_\mu [^\mu \left(𝒟^2\phi u^\nu _\nu \phi \right)^\mu _\nu \phi u^\nu 𝒟^2\phi `$
$`_\nu \left(^\mu \phi u^\nu \right)𝒟^2\phi ]+2d_1u_\mu _\nu F^{\nu \mu }.`$
Using the results of Appendix (A), it is easy to show that for the FRW background $`\overline{\phi }`$ satisfies Eq. (13) which is same as that of the canonical scalar field. The field equation for $`u_\mu `$ gives $`\overline{\lambda }=0`$.
Before we proceed to the computation of the stress-tensor, it is important to know the current astrophysical constraints on the parameters of the model Mattingly (2005): The constraints of the parameter $`d_1`$ comes from the big-bang nucleosynthesis Carroll and Lim (2004) and the solar system tests of general relativity Graesser et al. (2005). These give:
$$0<\frac{d_1}{M_{_{\mathrm{Pl}}}^2}<\frac{1}{7}.$$
(36)
The constraints on the parameter $`b_{11}`$ comes from the observations of highest energy cosmic rays Gagnon and Moore (2004). Using effective field theory with higher-dimensional operators, resulting in the modified field theory with a dispersion relation, it was shown that for various standard model particles
$$b_{11}M_{_{\mathrm{Pl}}}^2<5\times 10^5.$$
(37)
### III.2 Stress tensor
In this subsection, we obtain the stress-tensor for the scalar field defined in Eq. (28). Formally, the stress-tensor for a general Lagrangian containing up to first order derivative in the metric is given by (cf. Ref. Landau and Lifshitz (1975), p. 272)
$$T_{\mu \nu }=g_{\mu \nu }+2\frac{}{g^{\mu \nu }}\frac{2}{\sqrt{g}}_\rho (\sqrt{g}\frac{}{(_\rho g^{\mu \nu })}).$$
(38)
For simplicity, we will separate the contributions from the three different Lagrangians defined in Eq. (28), i. e.,
$$T_{\mu \nu }=T_{\mu \nu }^{(\phi )}+T_{\mu \nu }^{(u)}+T_{\mu \nu }^{(cor)}.$$
(39)
The stress-tensor $`T_{\mu \nu }^{^{(\phi )}}`$ corresponding to the canonical scalar field Lagrangian is given in Eq. (9). The stress-tensor corresponding to the Lagrangian $`__u`$ can easily be obtained and is given by
$`T_{\mu \nu }^{^{(u)}}`$ $`=`$ $`d_1g_{\mu \nu }F_{\epsilon \varkappa }F^{\epsilon \varkappa }4d_1F_{\mu \epsilon }F_{\nu \varkappa }g^{\epsilon \varkappa }`$ (40)
$`+\lambda \left[g_{\mu \nu }(g_{\epsilon \varkappa }u^\epsilon u^\varkappa 1)2u_\mu u_\nu \right].`$
However, the stress-tensor corresponding to $`_{_{cor}}`$ is much more involved. For the sake of continuity we give below only the final result while the steps leading to the result is given in Appendix (A). We get,
$$T_{\mu \nu }^{^{\left(\mathrm{cor}\right)}}=b_{_{11}}g_{_{\mu \nu }}\left(𝒟^2\phi \right)^2+b_{_{11}}E_{_{\mu \nu }}𝒟^2\phi +4b_{_{11}}C_{_{\mu \nu }}^{\rho \varkappa }_\varkappa \phi _\rho \left[𝒟^2\phi \right],$$
(41)
where,
$`E_{\mu \nu }`$ $`=`$ $`4[_\rho [\mathrm{ln}(\sqrt{g})]C_{\mu \nu }^{\rho \varkappa }_\varkappa \phi +_\rho \left(C_{\mu \nu }^{\rho \varkappa }_\varkappa \phi \right)`$ (42)
$`A_{\mu \nu }^{\epsilon \varkappa }_\epsilon _\varkappa \phi B_{\mu \nu }^\varkappa _\varkappa \phi ]`$
$`C_{\mu \nu }^{\rho \varkappa }`$ $`=`$ $`{\displaystyle \frac{1}{2}}[g_{\mu \nu }g^{\varkappa \rho }\delta _\mu ^\rho \delta _\mu ^\varkappa \delta _\mu ^\varkappa \delta _\nu ^\rho +u__\mu \delta _\nu ^\rho u^\varkappa u__\mu u__\nu g^{\varkappa \rho }`$ (43)
$`+`$ $`u__\mu \delta _\nu ^\varkappa u^\rho +u__\nu \delta _\mu ^\varkappa u^\rho g_{\mu \nu }u^\rho u^\varkappa +\delta _\mu ^\rho u_\nu u^\varkappa ].`$
Before proceeding with the evaluation of the perturbed stress-tensor we would like to mention the following point: Using the results of Appendix (A), it is clear that $`T_{\mu \nu }^{(u)}`$ and $`T_{\mu \nu }^{(cor)}`$ vanishes in the unperturbed FRW background. Hence, the equations determining the evolution of the scale factor, i. e.,
$$3^2=\frac{a^2}{M_{_{\mathrm{Pl}}}^2}\overline{T_0^0}^{(\phi )};2^{^{}}+^2=\frac{a^2}{3M_{_{\mathrm{Pl}}}^2}\overline{T_i^i}^{(\phi )},$$
(44)
remain the same as in the canonical scalar field inflation. It is also worth mentioning that the trans-Planckian corrections do not play any role on the expansion of the FRW background while, as we will see in the next section, the trans-Planckian corrections affect the metric and inflaton perturbations.
In this work, we will focus on the power-law inflation, for which, the scale factor is given by
$$a(t)=\left(a_0t^p\right)\mathrm{or}a(\eta )=\left(\frac{\eta }{\eta _0}\right)^{(\beta +1)},$$
(45)
where $`p>1`$, $`\beta 2`$ ($`\beta =2`$ corresponds to de Sitter), $`a_0`$ is a constant,
$$\beta =\left(\frac{2p1}{p1}\right)\mathrm{and}(\eta _0)=\frac{a_0^{1/p}}{(p1)}.$$
(46)
The scalar field potential and other background field parameters are given by ($`q=\sqrt{2/p}`$)
$`V=v_oM_{_{\mathrm{Pl}}}^4\mathrm{exp}\left(q{\displaystyle \frac{\overline{\phi }}{M_{_{\mathrm{Pl}}}}}\right)`$ ; $`={\displaystyle \frac{(1+\beta )}{(\eta )}};`$ (47)
$`\overline{\phi }=\sigma _oM_{_{\mathrm{Pl}}}\mathrm{ln}\left({\displaystyle \frac{\eta }{\eta _0}}\right)`$ ; $`\sigma _o=\sqrt{2\beta (\beta +1)};`$
$`a_0={\displaystyle \frac{\sqrt{(\beta +1)(1+2\beta )}}{\sqrt{v_o}\eta _0M_{_{\mathrm{Pl}}}}}`$ ; $`{\displaystyle \frac{z^{\prime \prime }}{z}}={\displaystyle \frac{a^{\prime \prime }}{a}}={\displaystyle \frac{\beta (\beta +1)}{(\eta )^2}}.`$
## IV Perturbed stress tensor
In this section, we will obtain the perturbed stress-tensor for the scalar field with modified dispersion relation (28) in the perturbed FRW background (II.1). As in the previous section, we will separate the contributions to the perturbed stress-tensor from the three different Lagrangians defined in Eq. (28), i. e.<sup>2</sup><sup>2</sup>2The mixed stress-tensor $`\delta T_\nu ^\mu `$ is given by
$$\delta T_\nu ^\mu \delta \left(g^{\mu ϵ}T_{ϵ\nu }\right)=\overline{g^{\mu ϵ}}(\delta T_{ϵ\nu })+(\delta g^{\mu ϵ})\overline{T_{ϵ\nu }}$$
(48) ,
$$\delta T_\nu ^\mu =\delta T_{\nu }^{\mu }{}_{}{}^{(\phi )}+\delta T_{\nu }^{\mu }{}_{}{}^{(u)}+\delta T_{\nu }^{\mu }{}_{}{}^{(cor)}.$$
(49)
The first term in the RHS of the above expression corresponds to the perturbed stress-tensor of the canonical scalar field Lagrangian and is given by Eqs. (II.2). In the rest of the section, we will obtain contributions from the other two terms.
Perturbing Eq. (40) and using the fact that $`F_{\mu \nu }`$ vanishes for the FRW background \[See Appendix (A)\], we get
$$\delta T_{\nu }^{\mu }{}_{}{}^{(u)}=2\delta _0^\mu \delta _\nu ^0(\delta \lambda ),$$
(50)
where $`\delta \lambda `$ is given by Eq. (150). Using the fact that $`\overline{𝒟^2(\phi )}`$ and $`\overline{_\rho 𝒟^2(\phi )}`$ vanish for the FRW background \[See Appendix (A)\], the perturbation of Eq. (41) takes the following simple form:
$$\delta T_{\mu \nu }^{(cor)}=4b_{11}\left[\overline{E_{\mu \nu }}\delta (𝒟^2\phi )+\overline{C_{\mu \nu }^{\rho 0}}\overline{\phi }^{}_\rho (\delta 𝒟^2\phi )\right].$$
(51)
Substituting the relation (147) for $`\delta (𝒟^2\phi )`$ and Eqs. (143) in the above expression, we get
$`\delta T_{\nu }^{\mu }{}_{}{}^{(cor)}={\displaystyle \frac{2b_{11}}{a^4}}[5{\displaystyle \frac{}{a}}\overline{\phi }_{}^{^{}}{}_{}{}^{2}^2\xi ^{(gi)}{\displaystyle \frac{1}{a}}\overline{\phi }_{}^{}{}_{}{}^{2}^2\xi _{}^{(gi)}{}_{}{}^{}`$ (52)
$`(\overline{\phi }^{^{\prime \prime }}+4\overline{\phi }^{^{}})^2(\delta \phi )+\overline{\phi }^{^{}}^2(\delta \phi )^{}]\delta ^\mu __0\delta ^0__\nu ,`$
where $`\xi `$ is defined in Eq. (158). Substituting Eqs. (II.2, 50, 52) in Eq. (49), we obtain the perturbed gauge-invariant stress-tensor to be:
$`\delta T_{0}^{0}{}_{}{}^{(gi)}`$ $`=`$ $`{\displaystyle \frac{1}{a^2}}\left[\overline{\phi ^{}}^2\mathrm{\Phi }+\overline{\phi ^{}}\delta \phi _{}^{(gi)}{}_{}{}^{^{}}+V_{,\phi }a^2\delta \phi ^{(gi)}\right]`$ (53a)
$`+{\displaystyle \frac{4d__1}{a^2}}\left[^2\mathrm{\Phi }{\displaystyle \frac{1}{a}}^2\xi _{}^{(gi)}{}_{}{}^{^{}}\right],`$
$`\delta T_{j}^{i}{}_{}{}^{(gi)}`$ $`=`$ $`\left[\overline{\phi ^{}}^2\mathrm{\Phi }\overline{\phi ^{}}\delta \phi ^{(gi)}+V_{,\phi }a^2\delta \phi ^{(gi)}\right]\delta __j^i,`$ (53b)
$`\delta T_{i}^{0}{}_{}{}^{(gi)}`$ $`=`$ $`a^2\overline{\phi ^{}}\delta \phi _{,i}^{(gi)}.`$ (53c)
Following points are worth-noting regarding the above result: Firstly, the two corrective Lagrangians – $`_u`$ and $`_{_{cor}}`$ – contributes only to the perturbed energy density $`(\delta \rho )`$. This implies that the non-diagonal space-space components of the stress-tensor are absent leading to the condition that $`\mathrm{\Phi }=\mathrm{\Psi }`$. This also implies that the constraint equation (17b) remains unchanged even for trans-Planckian inflation. Secondly, since the trans-Planckian corrections do not change the pressure perturbations, the perturbation equations for the tensor modes do not change. Hence, the tensor perturbation equations remain unchanged. Recently, Lim Lim (2005) had show that general Lorentz violating models (with out taking into account the higher derivatives of the scalar field) can modify the pressure perturbations and hence the tensor perturbation equations. However, in our specific Lorentz violating model, this is not the case. This indicates that the well-know consistency relation between the scalar and tensor ratio will also be broken in this model Hui and Kinney (2002); Ashoorioon and Mann (2005); Ashoorioon et al. (2005).
Thirdly, it is interesting to note that the trans-Planckian contributions to the energy density go as $`k^2`$. Hence as one would expect, in the super-Hubble scales, only the canonical scalar field contributes significantly in these scales. Lastly, these can have two significant implications on the perturbation spectrum: (a) The speed of propagation of the perturbations ($`c_s^2`$) can be different from that of the standard single-scalar field inflation. In the case of single scalar field inflation, we know that $`c_s^2=1`$. However, due to the extra contributions to the energy density, this can no longer be true. (b) The perturbations need not be purely adiabatic. $`\xi `$ can act as an extra scalar field during the inflation and hence can act as a source. This can introduce non-adiabatic (entropic) perturbations. (See, for example, Ref. Groot Nibbelink and van Tent (2002); van Tent (2004)) We will discuss more on these in the following sections.
## V Scalar Perturbation equation
Substituting Eqs. (53) in (17), the first-order perturbed Einstein’s equations take the following form,
$`^2\mathrm{\Phi }3\mathrm{\Phi }^{}3^2\mathrm{\Phi }={\displaystyle \frac{2d__1}{M_{_{\mathrm{Pl}}}^2}}\left[^2\mathrm{\Phi }{\displaystyle \frac{1}{a}}^2\xi _{}^{\left(gi\right)}{}_{}{}^{^{}}\right]`$ (54a)
$`+{\displaystyle \frac{1}{2M_{_{\mathrm{Pl}}}^2}}\left[\overline{\phi ^{}}^2\mathrm{\Phi }+\overline{\phi ^{}}\delta \phi _{}^{\left(gi\right)}{}_{}{}^{^{}}+V_{,\phi }a^2\delta \phi ^{\left(gi\right)}\right],`$
$`\mathrm{\Phi }+\mathrm{\Phi }^{}={\displaystyle \frac{1}{2M_{_{\mathrm{Pl}}}^2}}\overline{\phi ^{}}\delta \phi ^{\left(gi\right)},`$ (54b)
$`\mathrm{\Phi }^{\prime \prime }+3\mathrm{\Phi }^{}+(2^{}+^2)\mathrm{\Phi }={\displaystyle \frac{1}{2M_{_{\mathrm{Pl}}}^2}}[\overline{\phi ^{}}^2\mathrm{\Phi }\overline{\phi ^{}}\delta \phi ^{\left(gi\right)}`$
$`+V_{,\phi }a^2\delta \phi ^{\left(gi\right)}].`$ (54c)
As in the canonical scalar field inflation, the three perturbed Einstein’s (54) can be combined to give
$$\begin{array}{c}\mathrm{\Phi }^{^{\prime \prime }}\left(1\frac{2d__1}{M_{_{\mathrm{Pl}}}^2}\right)^2\mathrm{\Phi }+2\left(\frac{\overline{\phi }^{^{\prime \prime }}}{\overline{\phi }^{^{}}}\right)\mathrm{\Phi }^{^{}}\hfill \\ \hfill +2\left(^{^{}}\frac{\overline{\phi }^{^{\prime \prime }}}{\overline{\phi }^{^{}}}\right)\mathrm{\Phi }=\frac{2d__1}{M_{_{\mathrm{Pl}}}^2}\frac{1}{a}^2\xi _{}^{(gi)}{}_{}{}^{^{}}.\end{array}$$
(55)
Perturbing the field equations (III.1, III.1), we get
$`\delta \phi _{}^{\left(gi\right)}{}_{}{}^{\prime \prime }+2\delta \phi _{}^{\left(gi\right)}{}_{}{}^{}^2\left(\delta \phi ^{\left(gi\right)}\right)+V_{,\phi \phi }a^2\delta \phi ^{\left(gi\right)}`$ (56)
$`4\overline{\phi }^{}\mathrm{\Phi }^{}+2V_{,\phi }a^2\mathrm{\Phi }+{\displaystyle \frac{2b_{_{11}}}{a^2}}\left[^4\delta \phi ^{\left(gi\right)}{\displaystyle \frac{\overline{\phi }^{}}{a}}^4\xi ^{\left(gi\right)}\right]=0,`$
$`_m\left[\left(1{\displaystyle \frac{c_1}{M_{_{\mathrm{Pl}}}^2}}{\displaystyle \frac{^2}{a^2}}\right)\mathrm{\Phi }2\left(1+{\displaystyle \frac{c_1}{2M_{_{\mathrm{Pl}}}^2}}{\displaystyle \frac{^2}{a^2}}\right)\mathrm{\Phi }^{}\right]`$
$`{\displaystyle \frac{1}{a}}_m\left[\xi _{}^{\left(gi\right)}{}_{}{}^{^{\prime \prime }}3\xi _{}^{\left(gi\right)}{}_{}{}^{^{}}{\displaystyle \frac{b_{11}}{2d_1}}{\displaystyle \frac{\overline{\phi }^2}{a^2}}^2\xi ^{\left(gi\right)}\right]=0,`$ (57)
where $`c_1=(M_{_{\mathrm{Pl}}}^4b_{11})/d_1`$ is a dimensionless constant.
Following points are interesting to note regarding the above results: (i) The spatial higher derivatives appear only in the equation of motion of $`\delta \phi `$ and not in metric perturbation equation $`\mathrm{\Phi }`$. Unlike the standard inflation, the perturbations are not purely adiabatic and the speed of propagation of the perturbations is less than unity i. e. $`c_s^2=12d_1/M_{_{\mathrm{Pl}}}^2`$. (ii) In the case of canonical scalar-field inflation, the two dynamical variables $`\mathrm{\Phi }`$ and $`\delta \phi ^{(gi)}`$ are related by the constraint equation (17b). In our model, $`\mathrm{\Phi }`$ and $`\delta \phi ^{(gi)}`$ are again related by the same constraint equation, however $`\xi `$ is related to the other fields via the equations of motion. Hence, unlike the standard inflation, we have two sets of independent variables. (iii) In the case of canonical scalar-field inflation, $`\mathrm{\Phi }`$ and $`\delta \phi ^{(gi)}`$ can be combined into a single variable — Mukhanov-Sasaki ($`Q`$) variable — in terms of which we can obtain the perturbation spectrum. However, in this model, as in the case of multi-field inflation models Groot Nibbelink and van Tent (2002), the equations of motion in terms of the Mukhanov-Sasaki variables are coupled. In the rest of the section, we derive the equation of motion of the Mukhanov-Sasaki variables corresponding to $`\delta \phi ^{(gi)}`$ and $`\xi `$.
Substituting $`\delta \phi `$ in-terms of $`Q`$ in Eq. (56), and using the relations (55, 13, 17b), we get
$`Q^{\prime \prime }\left(1{\displaystyle \frac{2b_{11}}{a^2}}^2\right)^2Q{\displaystyle \frac{z^{\prime \prime }}{z}}Q={\displaystyle \frac{2d_1}{M_{_{\mathrm{Pl}}}^2}}^2𝒮(\eta ),`$ (58)
where
$$𝒮=\overline{\phi }^{^{}}\left[\frac{Q_\xi ^{(2)}}{}+\frac{c_1}{a^{1/2}}^2Q_\xi ^{(1)}\right],$$
(59)
and $`Q_\xi ^{(1)},Q_\xi ^{(2)}`$ are the gauge-invariant variables associated with $`\xi `$ and are given by:
$$Q_\xi ^{(1)}=a^{3/2}\left[\xi +\frac{a}{}\psi \right];Q_\xi ^{(2)}=\xi ^{^{}}a\varphi .$$
(60)
Substituting for $`\xi `$ in-terms of $`Q_\xi ^{(1)}`$ in (57), we get,
$`[a^{3/2}(Q_{\xi }^{\left(1\right)}{}_{}{}^{^{\prime \prime }}+[{\displaystyle \frac{3}{2}}^{^{}}{\displaystyle \frac{9}{4}}^2{\displaystyle \frac{b_{_{11}}\overline{\varphi }_{}^{^{}}{}_{}{}^{2}}{2d_1a^2}}^2]Q_\xi ^{\left(1\right)})`$ (61)
$`{\displaystyle \frac{\overline{\varphi }^{^{}}}{M_{_{\mathrm{Pl}}}^2}}({\displaystyle \frac{\overline{\varphi }^{^{\prime \prime }}}{\overline{\varphi }^{^{}}}}{\displaystyle \frac{^{^{}}}{}}{\displaystyle \frac{c_1}{2M_{_{\mathrm{Pl}}}^2a^2}}^2)Q{\displaystyle \frac{a}{}}^2\mathrm{\Phi }],m=0.`$
Decomposing $`Q,Q_\xi ^{(1)},Q_\xi ^{(2)}`$ into Fourier space, we have
$`\mu __S^{^{\prime \prime }}+\left[k^2+{\displaystyle \frac{2b_{_{11}}}{a^2\left(\eta \right)}}k^4{\displaystyle \frac{z^{^{\prime \prime }}}{z}}\right]\mu __S={\displaystyle \frac{2d_1}{M_{_{\mathrm{Pl}}}^2}}k^2𝒮_k\left(\eta \right)`$ (62)
$`\mu __\xi ^{^{\prime \prime }}+[{\displaystyle \frac{3}{2}}^{^{}}{\displaystyle \frac{9}{4}}^2+{\displaystyle \frac{b_{_{11}}\overline{\varphi }_{}^{^{}}{}_{}{}^{2}}{2d_1a^2}}k^2]\mu __\xi =a^{3/2}[{\displaystyle \frac{ak^2}{}}\mathrm{\Phi }_k`$ (63)
$`+{\displaystyle \frac{\overline{\varphi }^{^{}}}{M_{_{\mathrm{Pl}}}^2}}({\displaystyle \frac{\overline{\varphi }^{^{\prime \prime }}}{\overline{\varphi }^{^{}}}}{\displaystyle \frac{^{^{}}}{}}+{\displaystyle \frac{c_1}{2M_{_{\mathrm{Pl}}}^2a^2}}k^2)\mu __S],`$
where the Fourier transform of $`Q`$, $`Q_\xi ^{(1)}`$, $`Q_\xi ^{(2)}`$, respectively, are $`\mu __S,\mu __\xi ,Q_k^{(2)}`$ and
$$𝒮_k(\eta )=\overline{\phi }^{^{}}\left[\frac{Q_k^{(2)}}{}\frac{c_1}{a^{1/2}}k^2\mu __\xi \right].$$
(64)
Eqs. (62, 63) are the main results of our paper, regarding which we would like to stress the following points: Firstly, in the earlier analyses, the equation of motion of the Mukhanov-Sasaki variable ($`Q`$) was assumed to be satisfy the differential equation Eq. (62) in which the source term was assumed to be zero. We have shown explicitly from the gauge invariant perturbation theory that, in general, this is not true. The RHS of (62) vanishes in the super-Hubble scales (i.e $`k0`$) where the perturbations can be treated classical. Hence, as expected, the trans-Planckian effects are negligible. Secondly, it is clear from Eq. (62) that the terms in the RHS will dominate during the trans-Planckian regime and can have interesting consequences on the primordial spectrum. Lastly, the perturbations (in general) are not purely adiabatic, i. e., it contains isocurvature perturbations. However, these perturbations does not contribute significantly in the super-Hubble scales. Taking the Fourier transformation of the non-adiabatic part of the pressure perturbation ($`\delta p_{\mathrm{nad}}`$), we have
$$(\delta p_{\mathrm{nad}})=4d_1\frac{k^2}{a^3}Q_k^{(2)}.$$
(65)
From the above expression, it is straight forward to see that, in the super-horizon scales, the entropic perturbations vanish. Following Refs. Wands et al. (2000); Malik et al. (2003), we can assume that, on large scales, the total curvature perturbation $`\zeta `$ is conserved. As mentioned earlier, in the FRW background only the canonical scalar field contributes to the stress-tensor. Following Ref.Malik and Wands (2005), it is possible to show that, on large scales, only the curvature perturbation associated to $`\delta \varphi `$ contributes to the total curvature perturbation. Hence, it is sufficient to calculate the power-spectrum associated to the scalar-field perturbation ($`\delta \phi `$). This will be discussed in Sec. (VII)
## VI Classical Analysis
In this section, we combine Eqs. (55,56,57) to obtain a single differential equation of $`\mathrm{\Phi }`$. We show that the resultant differential equation of $`\mathrm{\Phi }`$ is different from that of the standard canonical scalar field driven inflation. More importantly, the differential equation of $`\mathrm{\Phi }`$ in our model is fourth order while in the standard canonical scalar field it is second order. We obtain the solutions of $`\mathrm{\Phi }`$ in three regimes — trans-Planckian (I), linear (II) and super-Hubble (III) — for the power-law inflation. In the following section, we obtain the power-spectrum of the scalar perturbations, in a particular limit, for the power-law inflation.
### VI.1 The Power law inflation.
Let us decompose the fields in their Fourier modes:
$`\mathrm{\Phi }(\eta ,\stackrel{}{x})`$ $`=`$ $`\mathrm{\Phi }_k(\eta )e^{i\stackrel{}{k}.\stackrel{}{x}},\delta \phi (\eta ,\stackrel{}{x})=\delta \phi _k(\eta )e^{i\stackrel{}{k}.\stackrel{}{x}}`$
$`\delta \xi (\eta ,\stackrel{}{x})`$ $`=`$ $`\xi _k(\eta )e^{i\stackrel{}{k}.\stackrel{}{x}}.`$ (66)
We have dropped the superscript indicating that the quantities are gauge invariant. Combining the equations, we end up with a fourth order differential equation in the Bardeen potential. To keep the presentation light, the derivation of the equation is given in Appendix C. This equation, for the power-law inflation (45,III.2), reads
$`\mathrm{\Phi }_k^{\left(4\right)}`$ $`+`$ $`{\displaystyle \frac{\mathrm{\Gamma }_1}{\eta }}\mathrm{\Phi }_k^{\left(3\right)}`$
$`+`$ $`\left[{\displaystyle \frac{\mathrm{\Gamma }_2}{\eta ^2}}+\left(1+{\displaystyle \frac{\mathrm{\Gamma }_3}{\left(\eta \right)^{\left(5+3\beta \right)}}}\right)k^2+{\displaystyle \frac{\mathrm{\Gamma }_4}{\left(\eta \right)^{2\left(1+\beta \right)}}}k^4\right]\mathrm{\Phi }_k^{\left(2\right)}`$
$`+`$ $`\left[{\displaystyle \frac{\mathrm{\Gamma }_5}{\eta ^3}}+\left({\displaystyle \frac{\mathrm{\Gamma }_6}{\left(\eta \right)^{3\left(2+\beta \right)}}}+{\displaystyle \frac{\mathrm{\Gamma }_7}{\eta }}\right)k^2+{\displaystyle \frac{\mathrm{\Gamma }_8}{\left(\eta \right)^{3+2\beta }}}k^4\right]\mathrm{\Phi }_k^{^{}}`$
$`+`$ $`\left[{\displaystyle \frac{\mathrm{\Gamma }_9}{\eta ^4}}+{\displaystyle \frac{\mathrm{\Gamma }_{10}}{\eta ^2}}k^2+\left({\displaystyle \frac{\mathrm{\Gamma }_{11}}{\left(\eta \right)^{2\beta +4}}}+{\displaystyle \frac{\mathrm{\Gamma }_{12}}{\left(\eta \right)^{5+3\beta }}}\right)k^4\right]\mathrm{\Phi }_k=0.`$
The constants $`\mathrm{\Gamma }_i`$ depend on the background and the fundamental constants in the following way
$`\mathrm{\Gamma }_1`$ $`=`$ $`4(2+\beta ),\mathrm{\Gamma }_2=12+13\beta +3\beta ^26\sigma _o^2,`$
$`\mathrm{\Gamma }_3`$ $`=`$ $`{\displaystyle \frac{3}{2}}{\displaystyle \frac{(1)^{3\beta }\eta _o^{3+3\beta }\sigma _o^2}{a_o^3}}{\displaystyle \frac{b_{11}}{d_1}}M_{pl}^2,`$
$`\mathrm{\Gamma }_4`$ $`=`$ $`2{\displaystyle \frac{(1)^{2\beta }\eta _o^{2+2\beta }}{a_o^2}}b_{11},\mathrm{\Gamma }_5=18(1+\beta )\sigma _o^2,`$
$`\mathrm{\Gamma }_6`$ $`=`$ $`3{\displaystyle \frac{(1)^{3\beta }\eta _o^{3+3\beta }\sigma _o^2}{a_o^3}}{\displaystyle \frac{b_{11}}{d_1}}M_{pl}^2,\mathrm{\Gamma }_7=6+4\beta ,`$
$`\mathrm{\Gamma }_8`$ $`=`$ $`4(2+\beta ){\displaystyle \frac{(1)^{2\beta }\eta _o^{2+2\beta }}{a_o^2}}b_{11},`$
$`\mathrm{\Gamma }_9`$ $`=`$ $`6a_o^2v_o\eta _o^2{\displaystyle \frac{q^2(24+q^2)}{(6+q^2)^2}}M_{pl}^2,\mathrm{\Gamma }_{10}=4+7\beta +3\beta ^2,`$
$`\mathrm{\Gamma }_{11}`$ $`=`$ $`2(1+\beta )(2+\beta ){\displaystyle \frac{(1)^{2\beta }\eta _o^{2+2\beta }}{a_o^2}}b_{11},`$
$`\mathrm{\Gamma }_{12}`$ $`=`$ $`{\displaystyle \frac{3}{2}}{\displaystyle \frac{(1)^{3\beta }\eta _o^{3+3\beta }\sigma _o^2}{a_o^3}}{\displaystyle \frac{b_{11}}{d_1}}(2d_1M_{pl}^2).`$ (68)
### VI.2 The zeroth order approximation
We can find approximate solutions for the power law inflation in the following way. Introducing the quantity $`ϵ`$ by
$$\beta =2ϵ,$$
(69)
($`ϵ`$ vanishes on the de Sitter space) we can make Taylor expansions of the coefficients $`\mathrm{\Gamma }_i`$ and postulate the same for the Bardeen potential
$$\mathrm{\Phi }_k(\eta )=\underset{m=0}{}ϵ^m\mathrm{\Phi }_{k,m}(\eta ).$$
(70)
The outcome is the following. Each component $`\mathrm{\Phi }_{k,m}(\eta )`$ obeys a differential equation which is inhomogeneous, the source term depending on the preceding components $`\mathrm{\Phi }_{k,0}(\eta ),\mathrm{},\mathrm{\Phi }_{k,m1}(\eta )`$.
As discussed in Ref. Martin and Brandenberger (2003), we have to be careful in the limit of $`ϵ0`$ in the sense that it does not give the perturbation corresponding to the de Sitter space. This is related to the fact that on this background the inflaton field is constant so that the quantities $`X_i,Y_i,Z_i`$ are undetermined. To work out what happens for the de Sitter space, one would have to consider the earlier equations where divisions by the derivatives of the inflaton was not performed yet.
The zeroth order contribution is a solution of the equation
$`\mathrm{\Phi }_{k,0}^{(4)}(\eta )`$ $`+`$ $`\left(k^2{\displaystyle \frac{2}{\eta ^2}}+\gamma ^2k^4\eta ^2\right)\mathrm{\Phi }_{k,0}^{(2)}(\eta )2{\displaystyle \frac{k^2}{\eta }}\mathrm{\Phi }_{k,0}^{^{}}(\eta )`$ (71)
$`+`$ $`2{\displaystyle \frac{k^2}{\eta ^2}}\mathrm{\Phi }_{k,0}^{^{}}(\eta )=0.`$
We now introduce the dimensionless variable $`x`$ defined by $`x=k\eta `$ and the function $`f(x)`$ given by
$`\mathrm{\Phi }_{k,0}(\eta )=f(x).`$ (72)
The fourth order equation takes the form
$$f^{\left(4\right)}\left(x\right)+\left(1\frac{2}{x^2}+\gamma ^2x^2\right)f^{\prime \prime }\left(x\right)\frac{2}{x}f^{}\left(x\right)+\frac{2}{x^2}f\left(x\right)=0,$$
(73)
where,
$$\gamma =\sqrt{2b_{11}v_0}M_{pl},$$
As mentioned earlier, we obtain approximate solutions to the above differential equation in three different regions.
#### VI.2.1 The first region
In this region, the term $`\gamma ^2x^2`$ dominates. In other words, the trans-Planckian effects are dominant and we are dealing with large values of $`x`$. Using the fact that $`f(x)=x`$ is a solution of the full equation, one introduces the function $`h(x)`$ by the relation
$$f(x)=x^xh(\zeta )𝑑\zeta $$
(74)
and obtains the third order equation
$$2(1+\gamma ^2x^2)h(x)+\gamma ^2x^3h{}_{}{}^{}(x)+4h{}_{}{}^{\prime \prime }(x)+xh^{^{\prime \prime \prime }}(x)=0$$
(75)
We can get rid of the second derivative by the change of function
$`h(x)`$ $`=`$ $`{\displaystyle \frac{1}{x^{4/3}}}R(x);`$ (76)
using the fact that we are in the region given by large values of $`x`$, one ends up with the differential equation
$`R^{^{\prime \prime \prime }}(x)+\gamma ^2x^2R^{^{}}(x)+{\displaystyle \frac{2}{3}}\gamma ^2xR(x)=0`$ (77)
whose solution is a combination of generalized hypergeometric functions multiplied by polynomials:
$`R(x)`$ $`=`$ $`C_1^{^{}}F_{pq}[\left\{{\displaystyle \frac{1}{6}}\right\},\{{\displaystyle \frac{1}{2}},{\displaystyle \frac{3}{4}}\},{\displaystyle \frac{1}{16}}\gamma ^2x^4]`$ (78)
$`+`$ $`C_2^{^{}}\sqrt{\gamma }xF_{pq}[\left\{{\displaystyle \frac{5}{12}}\right\},\{{\displaystyle \frac{3}{4}},{\displaystyle \frac{5}{4}}\},{\displaystyle \frac{1}{16}}\gamma ^2x^4]`$
$`+`$ $`C_3^{^{}}\gamma x^2F_{pq}[\left\{{\displaystyle \frac{2}{3}}\right\},\{{\displaystyle \frac{5}{4}},{\displaystyle \frac{3}{2}}\},{\displaystyle \frac{1}{16}}\gamma ^2x^4],`$
where $`C_i^{}`$’s are constants to be determined and $`F_{pq}`$ are the generalized Hypergeometric functions. Thus, the solution to the differential equation (73) is
$`f(x)`$ $`=`$ $`C_1(k)x`$
$`+`$ $`C_2(k)x^{2/3}F_{pq}[\{{\displaystyle \frac{1}{12}},{\displaystyle \frac{1}{6}}\},\{{\displaystyle \frac{1}{2}},{\displaystyle \frac{3}{4}},{\displaystyle \frac{11}{12}}\},{\displaystyle \frac{1}{16}}\gamma ^2x^4]`$
$`+`$ $`C_3(k)x^{5/3}F_{pq}[\{{\displaystyle \frac{1}{6}},{\displaystyle \frac{5}{12}}\},\{{\displaystyle \frac{3}{4}},{\displaystyle \frac{7}{6}},{\displaystyle \frac{5}{4}}\},{\displaystyle \frac{1}{16}}\gamma ^2x^4]`$
$`+`$ $`C_4(k)x^{8/3}F_{pq}[\{{\displaystyle \frac{5}{12}},{\displaystyle \frac{2}{3}}\},\{{\displaystyle \frac{5}{4}},{\displaystyle \frac{17}{12}},{\displaystyle \frac{3}{2}}\},{\displaystyle \frac{1}{16}}\gamma ^2x^4]`$
where $`C_i`$’s are related to $`C_i^{}`$’s. These generalized hypergeometric functions have a few properties which are worth mentioning. First, they are highly oscillating. For example, the function
$`F_{pq}[\{{\displaystyle \frac{1}{12}},{\displaystyle \frac{1}{6}}\},\{{\displaystyle \frac{1}{2}},{\displaystyle \frac{3}{4}},{\displaystyle \frac{11}{12}}\},{\displaystyle \frac{1}{16}}\gamma ^2x^4]`$ (80)
goes from $`\mathrm{7.7\hspace{0.17em}10}^{32}`$ to $`\mathrm{1.4\hspace{0.17em}10}^{194}`$ when $`x`$ goes from $`x=50`$ to $`x=100`$, fixing $`\gamma =1/10`$ for illustrative purposes.
Let us now say a few words about these generalized hypergeometric functions; this will help us to quantify their oscillatory behavior. They are special cases of Meijer functions which can be defined by integrals on the complex plane Mathai and Saxena (1973):
$$G_{p,q}^{m,n}\left(z|\begin{array}{cccc}a_1& a_2& \mathrm{}& a_p\\ b_1& b_2& \mathrm{}& b_q\end{array}\right)=\frac{1}{2\pi i}_C\chi (s)z^s𝑑s$$
(81)
where
$$\chi (s)=\frac{\mathrm{\Pi }_{j=1}^m\mathrm{\Gamma }(b_j+s)\mathrm{\Pi }_{j=1}^n\mathrm{\Gamma }(1a_js)}{\mathrm{\Pi }_{j=m+1}^q\mathrm{\Gamma }(1b_js)\mathrm{\Pi }_{j=n+1}^p\mathrm{\Gamma }(a_j+s)}$$
(82)
and three possibilities are allowed for the contour $`C`$, according to some conditions on the parameters $`a_i,b_j,m,n,p,q`$ Mathai and Saxena (1973). Our solutions correspond to $`m=n=0`$.
The asymptotic behavior which is relevant here is the following. For large values with $`(\nu ^{}+1)\pi <argz<0`$, the dominant part is roughly given by
$$G_{p,q}^{m,n}\left(z|\begin{array}{cccc}a_1& a_2& \mathrm{}& a_p\\ b_1& b_2& \mathrm{}& b_q\end{array}\right)H_{pq}(ze^{i\pi \mu ^{}}).$$
(83)
where
$`\mu ^{}`$ $`=`$ $`qmn,\nu ^{}=p+m+n,`$
$`H_{pq}(z)`$ $`=`$ $`\mathrm{exp}\left((pq)z^{\frac{1}{qp}}\right)z^\rho ^{},\mathrm{and}`$
$`\rho ^{}`$ $`=`$ $`{\displaystyle \frac{1}{qp}}\left({\displaystyle \underset{j=1}{\overset{q}{}}}b_j{\displaystyle \underset{j=1}{\overset{p}{}}}a_j+{\displaystyle \frac{pq+1}{2}}\right)`$ (84)
In our case, one has to make the replacements
$$z=\frac{1}{16}\gamma ^2k^4\eta ^4,q=3,,p=2$$
(85)
so that
$`\mathrm{\Phi }_{k,0}(\eta )`$ $`=`$ $`C_0(k)k\eta +{\displaystyle \underset{i=1}{\overset{3}{}}}C_i(k)(\kappa \eta )^{\sigma _i}`$
$`\times `$ $`\left({\displaystyle \frac{1}{16}}\gamma ^2k^4\eta ^4\right)^{\rho _i^{}}\mathrm{exp}\left({\displaystyle \frac{1}{16}}\gamma ^2k^4\eta ^4\mathrm{exp}(i\pi \mu _i^{})\right),`$
where $`\sigma _1=2/3,\sigma _2=5/3,\sigma _3=8/3`$. Fomr the above expressiom, it is easy to see that the solution in the the Bardeen potential $`\mathrm{\Phi }`$ is oscillating in this region. The choice of the constants $`C_i(k)`$ correspond to different choices of initial conditions and thus, in principle, to different choices of vacua. We will come back to this later.
#### VI.2.2 The second region
In the intermediary region, $`1`$ dominates over $`\gamma ^2x^2`$. The solution in this region is
$`f(x)`$ $`=`$ $`D_1(k)x+D_2(k)x^2`$ (87)
$`+`$ $`D_3(k)\left[e^{ix}(1+ix)x^2Ei(ix)\right]`$
$`+`$ $`D_4(k)\left[e^{ix}(ix)+ix^2Ei(ix)\right]`$
where $`Ei(x)`$ refers to the exponential integral. Using the asymptotic behavior of the exponential integral (cf. Ref. Abramowitz and Stegun (1964), p. 231), we get
$`f(x)`$ $`=`$ $`D_1(k)x+D_2(k)x^2`$ (88)
$`+`$ $`D_3(k)\left[e^{ix}+2x\mathrm{sin}(x)\right]`$
$`+`$ $`D_4(k)\left[ie^{ix}2x\mathrm{sin}(x)\right]`$
As we can see, the Bardeen potential is a sum of plane-waves.
#### VI.2.3 The third region
When the term $`2/x^2`$ dominates in the coefficient of the second derivative, the solution can be found and is given by
$`f(x)`$ $`=`$ $`G_1(k)+G_4(k)x+G_3(k)x^4+G_2(k)x\mathrm{ln}x.`$ (89)
From the above expression, we see that in the super-Hubble scales the scalar perturbations has a constant term which is identical to the canonical scalar field inflation.
To finish this section, let us remark that in the non trans-Planckian region, i.e when
$$1\frac{2}{x^2}>>\gamma ^2x^2,$$
the solution to the differential equation (73) can be obtained and is given by
$`f(x)`$ $`=`$ $`H_1(k)x+{\displaystyle \frac{1}{2}}H_2(k)(22x+x^2)`$ (90)
$`+`$ $`{\displaystyle \frac{1}{2}}H_3(k)x\left({\displaystyle \frac{e^{ix}}{x}}iEi(ix)\right)`$
$`+`$ $`{\displaystyle \frac{1}{2}}H_4(k)x\left({\displaystyle \frac{e^{ix}}{x}}iEi(ix)\right);`$
this approximation covers the $`II`$ and $`III`$ region simultaneously.
### VI.3 The first order approximation
The first order contribution obeys the equation
$`\mathrm{\Phi }_{k,1}^{(4)}(\eta )+\left(k^2{\displaystyle \frac{2}{\eta ^2}}+\gamma ^2k^4\eta ^2\right)\mathrm{\Phi }_{k,1}^{(2)}(\eta )2{\displaystyle \frac{k^2}{\eta }}\mathrm{\Phi }_{k,1}^{^{}}(\eta )`$
$`+2{\displaystyle \frac{k^2}{\eta ^2}}\mathrm{\Phi }_{k,1}(\eta )=S_k(\eta ),`$ (91)
where
$`S_k(\eta )`$ $`=`$ $`{\displaystyle \frac{4}{\eta }}\mathrm{\Phi }_{k,0}^{(3)}(\eta )+[{\displaystyle \frac{5}{\eta ^2}}+{\displaystyle \frac{b_{_{11}}M_{pl}^5v_o^{3/2}}{d_1}}k^2\eta `$
$``$ $`{\displaystyle \frac{10}{3}}b_{_{11}}M_{pl}^2v_ok^4\eta ^2+4b_{11}M_{pl}^2v_o\eta ^2\mathrm{log}\left({\displaystyle \frac{\eta }{\eta _o}}\right)]\mathrm{\Phi }_{k,0}^{(2)}(\eta )`$
$`+`$ $`\left[{\displaystyle \frac{12}{\eta ^3}}4{\displaystyle \frac{k^2}{\eta }}4b_{11}M_{pl}^2v_ok^4\eta \right]\mathrm{\Phi }_{k,0}^{^{}}(\eta )`$
$`+`$ $`[{\displaystyle \frac{24}{\eta ^4}}+{\displaystyle \frac{5k^2}{\eta ^2}}+2b_{11}M_{pl}^2v_o`$
$`+{\displaystyle \frac{b_{11}M_{pl}^3v_o^{3/2}(2d_1+M_{pl}^2)}{d_1}}\eta ]\mathrm{\Phi }_{k,0}(\eta ).`$
This equation is exactly the one obeyed by the zeroth order contribution, except for the source term which is known since we obtained the approximations of the zeroth order in the three regions. Let us specialize to one of the regions and call $`Y_1(\eta ),Y_2(\eta ),Y_3(\eta ),Y_4(\eta )`$ the four different solutions of the homogeneous equation given in Eq (71):
The equation being linear and knowing the complete solution of the homogeneous equation ($`\mathrm{\Phi }=_{a=1}^4L_aY_a`$ with $`L_a`$ constants) solution, we can solve it using the method of the variation of the constants. One can show that this can be achieved by the following system of equations
$`{\displaystyle \underset{a=1}{\overset{4}{}}}L_a^{^{}}Y_a=0`$ , $`{\displaystyle \underset{a=1}{\overset{4}{}}}L_a^{^{}}Y_a^{^{}}=0,`$
$`{\displaystyle \underset{a=1}{\overset{4}{}}}L_a^{^{}}Y_a^{(2)}=0`$ , $`{\displaystyle \underset{a=1}{\overset{4}{}}}L_a^{^{}}Y_a^{(3)}=S(\eta ).`$ (92)
Let us concentrate on the second and third region for example (the non trans-Planckian zone). One has
$`\mathrm{\Phi }_{k,1}(\zeta )`$ $`=`$ $`k\zeta {\displaystyle _0^\zeta }d\eta [i{\displaystyle \frac{(2i+2k\eta ik^2\eta ^2)}{2k^4\eta }}Ei(ik\eta )`$
$`+`$ $`{\displaystyle \frac{(2i+2k\eta +ik^2\eta ^2)}{2k^4\eta }}Ei(ik\eta )]S_k(\eta )`$
$`+`$ $`\left[1k\zeta +{\displaystyle \frac{1}{2}}k^2\zeta ^2\right]{\displaystyle _0^\zeta }𝑑\eta \mathrm{\hspace{0.17em}2}{\displaystyle \frac{1}{k^4\eta }}S_k(\eta )`$
$`+`$ $`\left[{\displaystyle \frac{1}{2}}ie^{ik\zeta }+{\displaystyle \frac{1}{2}}xEi(ik\zeta )\right]`$
$`\times {\displaystyle _0^\zeta }d\eta i{\displaystyle \frac{(2+2ik\eta +k^2\eta ^2)}{k^4\eta }}e^{ik\eta }S_k(\eta )`$
$`+`$ $`\left({\displaystyle \frac{1}{4}}ie^{ik\zeta }{\displaystyle \frac{1}{4}}ik\zeta Ei(ik\zeta )\right)`$
$`\times {\displaystyle _0^\zeta }d\eta \mathrm{\hspace{0.17em}2}{\displaystyle \frac{(22ik\eta +k^2\eta ^2)}{k^4\eta }}e^{ik\eta }S_k(\eta ).`$
A similar treatment can be applied to the trans-Planckian region but the formulas are too lengthy and will not be recorded here.
Using the analysis discussed in this section, the power spectrum of the perturbations can be obtained upto a $`k`$ dependent constant factor. In order to obtain the exact power spectrum, we need to quantize the theory and fix the initial state of the field Mukhanov et al. (1992). In the following section, we obtain exact power spectrum of the perturbations in a particular limit.
## VII Power-spectrum of the perturbations – Quantum Analysis
In this section, we calculate the power-spectrum corresponding to $`\mu __S`$ during the power-law inflation using the following approach: (i) We assume that the quantum field $`\mu __S`$ is coupled to an external, classical source field $`𝒮_k(\eta )`$ which is determined by solving the coupled differential equations (55, 57). (ii) We solve the equation of motion of $`\mu __S`$ in three regions – Trans-Planckian (I), linear (II) and super-Hubble (III) – separately Martin and Brandenberger (2001). We further assume that $`𝒮_k(\eta )`$ will contribute significantly in the trans-Planckian region while it can be neglected in the linear and super-Hubble region. (iii) The power-spectrum at the super-Hubble scales is determined by performing the matching of the modes and its derivatives at the times of transition between regions I and II \[$`(\eta _{\mathrm{Pl}})^{1+\beta }(\omega k^2)^{1/2}`$\] and regions II and III \[$`\eta __H(1+\beta )/k`$\]. We assume that the quantum field $`\mu __S`$ is in a minimum energy state at $`\eta =\eta _i`$ Brown and Dutton (1978).
Region (I) corresponds to the limit where the non-linearities of the dispersion relation play a dominant role, i. e. $`2b_{11}[k/a(\eta )]^21`$ and $`k\eta 1`$. Region (II) corresponds to the limit where the non-linearities of the modes are negligible i. e. $`\omega k`$ and $`k\eta 1`$. Region (III) corresponds to the limit where $`k\eta 1`$. In the three regions, the equation of motion of $`\mu __S`$ (62) reduces to:
$`\mu _{_S}^{(I)}{}_{}{}^{^{\prime \prime }}+\omega ^2(\eta )\mu __S^{(I)}`$ $``$ $`{\displaystyle \frac{2d_1}{M_{_{\mathrm{Pl}}}^2}}k^2𝒮_k(\eta ),`$ (94a)
$`\mu _{_S}^{(II)}{}_{}{}^{^{\prime \prime }}+k^2\mu __S^{(II)}`$ $``$ $`0,`$ (94b)
$`\mu _{_S}^{(III)}{}_{}{}^{^{\prime \prime }}{\displaystyle \frac{\beta (\beta +1)}{\eta ^2}}\mu __S^{(III)}`$ $``$ $`0,`$ (94c)
where
$$\omega (\eta )=\omega _0\frac{k^2}{(\eta )^{(1+\beta )}};\omega _0=(2b_{_{11}})^{1/2}(\eta _0)^{(1+\beta )},$$
(95)
and $`𝒮_k`$ is given by Eq. (64). The general solution to the differential equation (94a) is given by
$`\mu __S^{(\mathrm{I})}(\eta )`$ $`=`$ $`A_1(k)(\eta )^{1/2}H_\nu ^{(1)}[\alpha (\eta )]`$ (96)
$`+`$ $`A_2(k)(\eta )^{1/2}H_\nu ^{(2)}[\alpha (\eta )]+\mu __P(\eta )`$
where $`\mu __P(\eta )`$ is the particular solution to the inhomogeneous part of the differential equation and is given by (cf. Ref. Morse and Feshbach (1953), p. 529)
$$\mu __P\left(\eta \right)=\frac{i\pi }{2}\frac{d_1k^2}{\beta M_{_{\mathrm{Pl}}}^2}\left(\eta \right)^{1/2}\left[H_\nu ^{\left(1\right)}\left[\alpha \left(\eta \right)\right]_{\eta _l}^\eta \left(s\right)^{1/2}H_\nu ^{\left(2\right)}\left[\alpha \left(s\right)\right]𝒮_k\left(s\right)𝑑s+H_\nu ^{\left(2\right)}\left[\alpha \left(\eta \right)\right]_{\eta _l}^\eta \left(s\right)^{1/2}H_\nu ^{\left(1\right)}\left[\alpha \left(s\right)\right]𝒮_k\left(s\right)𝑑s\right]$$
(97)
$$\mathrm{and}\nu =\frac{1}{2\beta };\alpha (\eta )=\alpha _0(\eta )^\beta ;\alpha _0=\frac{\omega _0k^2}{\beta },$$
(98)
$`\eta _l`$ ($`<\eta _i`$) is the epoch in which the integrals in (97) vanish. The quantities $`H_\nu ^{(1)}`$ and $`H_\nu ^{(2)}`$ in the above solution are the Hankel functions of the first and the second kind (of order $`\nu `$), respectively, and the $`k`$-dependent constants $`A_1(k)`$ and $`A_2(k)`$ are to be fixed by the initial conditions for the modes at $`\eta _\mathrm{i}`$. Unlike the canonical scalar field inflation, where one assumes that the field is in a Bunch-Davies vacuum at $`\eta _i`$, it is not possible to assume such an initial condition due to the non-linearities of the modes. As mentioned earlier, we assume that the field is in the minimum energy vacuum state at $`\eta _i`$, i. e.,
$$\mu __S(\eta _\mathrm{i})=\frac{1}{\sqrt{2\omega (\eta _\mathrm{i})}};\mu __S^{}(\eta _\mathrm{i})=\pm i\sqrt{\frac{\omega (\eta _\mathrm{i})}{2}}.$$
(99)
We thus get
$`A_1(k)`$ $`=`$ $`{\displaystyle \frac{i\pi \alpha (\eta _i)}{4}}(\eta _\mathrm{i})^{1/2}\stackrel{~}{\mu }__S(\eta _i)H_{\nu 1}^{(2)}[\alpha (\eta _i)]`$ (100a)
$`\times `$ $`[1+{\displaystyle \frac{(\eta _\mathrm{i})^{\beta +1}}{(\alpha _0\beta )}}{\displaystyle \frac{\stackrel{~}{\mu }_{_S}^{}{}_{}{}^{}(\eta _\mathrm{i})}{\stackrel{~}{\mu }__S(\eta _\mathrm{i})}}{\displaystyle \frac{H_\nu ^{(2)}[\alpha (\eta _i)]}{H_{\nu 1}^{(2)}[\alpha (\eta _i)]}}],`$
$`A_2(k)`$ $`=`$ $`{\displaystyle \frac{i\pi \alpha (\eta _i)}{4}}(\eta _\mathrm{i})^{1/2}\stackrel{~}{\mu }__S(\eta _i)H_{\nu 1}^{(1)}[\alpha (\eta _i)]`$ (100b)
$`\times `$ $`[1+{\displaystyle \frac{(\eta _\mathrm{i})^{\beta +1}}{(\alpha _0\beta )}}{\displaystyle \frac{\stackrel{~}{\mu }_{_S}^{}{}_{}{}^{}(\eta _\mathrm{i})}{\stackrel{~}{\mu }__S(\eta _\mathrm{i})}}{\displaystyle \frac{H_\nu ^{(1)}[\alpha (\eta _i)]}{H_{\nu 1}^{(1)}[\alpha (\eta _i)]}}],`$
where $`\stackrel{~}{\mu }__S(\eta )=\mu __S^{(I)}(\eta )\mu __P(\eta )`$. From the above expressions it is evident that the particular solution to the inhomogeneous differential equation effectively changes the initial state of the field. Hence, $`\stackrel{~}{\mu }__S(\eta _\mathrm{i})`$ can be treated as the new effective initial state of the field. It is worth mentioning that the particular solution has not been fixed and is arbitrary.
In Region II, the solution to the differential equation (94b) is (Minkowski) plane waves, i. e.,
$$\mu __S^{(\mathrm{II})}(\eta )=B_1(k)\mathrm{exp}[ik\eta ]+B_2(k)\mathrm{exp}[ik\eta ],$$
(101)
where $`B_1(k),B_2(k)`$ are $`k`$dependent constants and are obtained by the junction conditions of the mode functions $`\mu __S^{(\mathrm{I})},\mu __S^{(\mathrm{II})}`$ and their derivatives at $`\eta =\eta _{\mathrm{Pl}}`$. This gives,
$`{\displaystyle \frac{\mathrm{exp}\left(ik\eta _{\mathrm{Pl}}\right)}{\left(\eta _{\mathrm{Pl}}\right)^{1/2}}}B_1`$ $`=`$ $`{\displaystyle \frac{A_1}{2}}H_\nu ^{\left(1\right)}\left[\alpha \left(\eta _{\mathrm{Pl}}\right)\right]\left[1+{\displaystyle \frac{i\beta \alpha _0}{k\left(\eta _{\mathrm{Pl}}\right)^{\beta +1}}}{\displaystyle \frac{H_{\nu 1}^{\left(1\right)}\left[\alpha \left(\eta _{\mathrm{Pl}}\right)\right]}{H_\nu ^{\left(1\right)}\left[\alpha \left(\eta _{\mathrm{Pl}}\right)\right]}}\right]+{\displaystyle \frac{A_2}{2}}H_\nu ^{\left(2\right)}\left[\alpha \left(\eta _{\mathrm{Pl}}\right)\right]\left[1+{\displaystyle \frac{i\beta \alpha _0}{k\left(\eta _{\mathrm{Pl}}\right)^{\beta +1}}}{\displaystyle \frac{H_{\nu 1}^{\left(2\right)}\left[\alpha \left(\eta _{\mathrm{Pl}}\right)\right]}{H_\nu ^{\left(2\right)}\left[\alpha \left(\eta _{\mathrm{Pl}}\right)\right]}}\right],`$ (102a)
$`{\displaystyle \frac{\mathrm{exp}\left(ik\eta _{\mathrm{Pl}}\right)}{\left(\eta _{\mathrm{Pl}}\right)^{1/2}}}B_2`$ $`=`$ $`{\displaystyle \frac{A_1}{2}}H_\nu ^{\left(1\right)}\left[\alpha \left(\eta _{\mathrm{Pl}}\right)\right]\left[1{\displaystyle \frac{i\beta \alpha _0}{k\left(\eta _{\mathrm{Pl}}\right)^{\beta +1}}}{\displaystyle \frac{H_{\nu 1}^{\left(1\right)}\left[\alpha \left(\eta _{\mathrm{Pl}}\right)\right]}{H_\nu ^{\left(1\right)}\left[\alpha \left(\eta _{\mathrm{Pl}}\right)\right]}}\right]+{\displaystyle \frac{A_2}{2}}H_\nu ^{\left(2\right)}\left[\alpha \left(\eta _{\mathrm{Pl}}\right)\right]\left[1{\displaystyle \frac{i\beta \alpha _0}{k\left(\eta _{\mathrm{Pl}}\right)^{\beta +1}}}{\displaystyle \frac{H_{\nu 1}^{\left(2\right)}\left[\alpha \left(\eta _{\mathrm{Pl}}\right)\right]}{H_\nu ^{\left(2\right)}\left[\alpha \left(\eta _{\mathrm{Pl}}\right)\right]}}\right].`$ (102b)
In region III, the solution is
$$\mu _{(\mathrm{III})}(\eta )=C(k)a(\eta ),$$
(103)
where $`C(k)`$ is constant (not to be confused with the constants used in the previous section) whose modulus square gives the power spectrum of the density perturbations and is determined by performing the matching of the modes $`\mu __S^{(\mathrm{II})},\mu __S^{(\mathrm{III})}`$ at $`\eta __H(1+\beta )/k`$. We thus get,
$$C\left(k\right)=\left[\frac{\eta __0k}{1+\beta }\right]^{1+\beta }\left[B_1\left(k\right)\mathrm{exp}\left(ik\eta __H\right)+B_2\left(k\right)\mathrm{exp}\left(ik\eta __H\right)\right].$$
(104)
The spectrum of the perturbations (25) reduce to
$$\left[k^3𝒫_S(k)\right]=\left(\frac{1}{4\pi ^2M_{_{\mathrm{Pl}}}^2}\frac{(\beta +1)}{\beta }\right)k^3\left|C(k)\right|^2.$$
(105)
We are interested in the leading order behavior of the primordial power-spectrum and the possible modifications to the primordial spectrum due to the trans-Planckian effects. In order to do that, we need to obtain the leading order behavior of the constants $`A_1,A_2,B_1`$ and $`B_2`$. Using the fact that $`k\eta 1`$ and the asymptotic behavior of the Hankel functions, viz. (cf. Ref. Abramowitz and Stegun (1964), p. 364)
$`\underset{z\mathrm{}}{lim}H_\nu ^{(1/2)}(z)\left({\displaystyle \frac{2}{\pi z}}\right)^{1/2}\mathrm{e}^{\pm i\left[z(\pi \nu /2)(\pi /4)\right]},`$ (106)
we get
$`A_1(k)`$ $``$ $`A_0\stackrel{~}{\mu }__S(\eta _i)\mathrm{exp}(ix_\mathrm{i})\left(1\right),`$ (107a)
$`A_2(k)`$ $``$ $`A_0\stackrel{~}{\mu }__S(\eta _i)\mathrm{exp}(ix_\mathrm{i})\left(1\pm \right),`$ (107b)
where
$`A_0=\left({\displaystyle \frac{\pi \alpha _0}{8}}\right)^{1/2}(\eta _\mathrm{i})^{(\beta +1)/2}`$ ; $`x_\mathrm{i}=\alpha (\eta _i){\displaystyle \frac{\pi \nu }{2}}{\displaystyle \frac{\pi }{4}}`$
$`={\displaystyle \frac{1\mu __P^{}(\eta _i)/\mu __S^{}(\eta _i)}{1\mu __P(\eta _i)/\mu __S(\eta _i)}}`$ $`.`$ (108)
Having obtained $`A_1,A_2`$ in the limit of $`k\eta 1`$. Our next step is to evaluate $`B_1(k),B_2(k)`$ in the same limit. In order to do that, we need to know the correct matching time $`\eta _{\mathrm{Pl}}`$. Demanding $`\omega ^2(\eta _{_{\mathrm{Pl}}})=k^2`$ gives $`(\eta _{\mathrm{Pl}})^{1+\beta }=\omega ^{1/2}/k`$. We thus get,
$`B_1`$ $``$ $`A_1\left({\displaystyle \frac{2\beta }{\pi k}}\right)^{1/2}\mathrm{exp}(ix_{\mathrm{Pl}})`$ (109)
$`B_2`$ $``$ $`A_2\left({\displaystyle \frac{2\beta }{\pi k}}\right)^{1/2}\mathrm{exp}(ix_{\mathrm{Pl}})`$ (110)
where $`x_{\mathrm{Pl}}=k\eta _{\mathrm{Pl}}(\beta +1)/\beta \pi \nu /2\pi /4`$. Thus, we get,
$`\left[k^3𝒫_S(k)\right]`$ $``$ $`C_0k^{2(\beta +2)}\left|1{\displaystyle \frac{\mu __P(\eta _i)}{\mu __S(\eta _i)}}\right|^2`$
$`\times `$ $`\left[1+2\mathrm{cos}(x__H)2Im[]\mathrm{sin}(x__H)\right]`$
where
$`C_0`$ $`=`$ $`\left({\displaystyle \frac{1}{16\pi ^2M_{_{\mathrm{Pl}}}^2}}{\displaystyle \frac{(\beta +1)}{\beta }}\right)\left({\displaystyle \frac{\eta _0\eta _i}{1+\beta }}\right)^{2(1+\beta )};`$
$`x__H`$ $`=`$ $`2(1+\beta x_{_{Pl}}+x_i)`$ (112)
and we have neglected higher order terms like $`||^2`$. It is interesting to note that in the limit of $`𝒮_k(\eta )0`$, the power-spectrum is same as that of the standard power-law inflation spectrum with small oscillations. In this limit, we recover the result of Refs. Martin and Brandenberger (2001); Niemeyer and Parentani (2001). In order to obtain the exact form of the power-spectrum, we need to evaluate $`\mu __P`$ which requires the knowledge of $`𝒮_k(\eta )`$.
In the rest of this section, we evaluate the power-spectrum in a particular limit ($`1/c_10`$). We, first, obtain the form of $`𝒮_k(\eta )`$ by solving the system of coupled differential equations (55, 57) in two – sub-Hubble and super-Hubble – regimes. As mentioned earlier, the two differential equations (55, 57) do not contain higher order spatial derivatives. Hence, it is sufficient to obtain solutions in these two regimes. Performing the following transformations
$$u=\frac{a(\eta )}{\overline{\phi }^{}(\eta )}\mathrm{\Phi };\xi ^{(gi)}=a^{3/2}(\eta )\stackrel{~}{\xi },$$
(113)
and taking the Fourier transform, Eqs. (55, 57), reduce to
$`u_k^{\prime \prime }+\left(c_s^2k^2{\displaystyle \frac{\theta ^{\prime \prime }}{\theta }}\right)u_k={\displaystyle \frac{2d_1}{M_{_{\mathrm{Pl}}}^2}}{\displaystyle \frac{k^2}{\overline{\phi }^{}}}\left(a^{3/2}\stackrel{~}{\xi }_k\right)^{^{}}`$
$`\stackrel{~}{\xi }^{\prime \prime }+\left[{\displaystyle \frac{9}{4}}^2+{\displaystyle \frac{3}{2}}^{}+{\displaystyle \frac{c_1}{a^2}}{\displaystyle \frac{\phi ^2}{2M_{_{\mathrm{Pl}}}^2}}{\displaystyle \frac{k^2}{M_{_{\mathrm{Pl}}}^2}}\right]\stackrel{~}{\xi }`$
$`=`$ $`a^{1/2}\left[\left(1+{\displaystyle \frac{c_1}{a^2}}{\displaystyle \frac{k^2}{M_{_{\mathrm{Pl}}}^2}}\right)\mathrm{\Phi }_k^{}2\left(1{\displaystyle \frac{c_1}{2a^2}}{\displaystyle \frac{k^2}{M_{_{\mathrm{Pl}}}^2}}\right)\mathrm{\Phi }_k\right].`$
In the limit of $`1/c_10`$ (i. e. $`d_1/M_{_{\mathrm{Pl}}}^2b_{_{11}}M_{_{\mathrm{Pl}}}^2`$), the above differential equations can be solved exactly. In this limit, the above differential equations become:
$`u_k^{\prime \prime }+\left(k^2{\displaystyle \frac{\theta ^{\prime \prime }}{\theta }}\right)u_k=0`$
$`{\displaystyle \frac{\phi ^2}{2M_{_{\mathrm{Pl}}}^2}}\xi _k^{(gi)}=a(\overline{\phi }^{}u_k)^{}.`$ (116)
For the sub-Hubble scales, during the power-law inflation, we get
$`u_k`$ $`=`$ $`D_1(k)\mathrm{exp}(ik\eta )+D_2(k)\mathrm{exp}(ik\eta )`$
$`\xi _k^{(gi)}`$ $`=`$ $`ik{\displaystyle \frac{(\eta )^{(\beta +2)}}{(\eta _0)^{(\beta +1)}}}\sqrt{{\displaystyle \frac{2M_{_{\mathrm{Pl}}}^2}{\beta (\beta +1)}}}`$ (117)
$`\left[D_1(k)\mathrm{exp}(ik\eta )D_2(k)\mathrm{exp}(ik\eta )\right],`$
where we have neglected the terms of the order $`1/(k\eta )`$ and $`D_1(k),D_2(k)`$ are $`k`$-dependent constants with the dimensions of length squared ($`k^2`$). Using the condition that the modes are outgoing, we set $`D_2(k)=0`$. In the super-Hubble scales, we have
$$u_kD_3(k)a(\eta );\xi _k^{(gi)}=D_3(k)\sqrt{\frac{2\beta }{\beta +1}}a^2(\eta ),$$
(118)
where $`D_3(k)`$ is a constant. In the sub-Hubble scales, we have
$$𝒮_k(\eta )=4iM_{_{\mathrm{Pl}}}^2b_{_{11}}D_1(k)k^5\left(\frac{\eta _0}{\eta }\right)^{\beta +1}\mathrm{exp}(ik\eta ).$$
(119)
Our next task is to obtain $`\mu __P`$ and the power-spectrum of the scalar perturbations. From Eq. (97) using the asymptotic limit of Hankel functions, we get
$$\begin{array}{c}\mu __S(\eta )=\left(\frac{b_{_{11}}}{\beta ^2}\right)^{\frac{\beta +1}{4\beta }}\frac{(\eta _0)^{\frac{1\beta ^2}{2\beta }}}{(\eta )^{\frac{\beta +1}{2}}}M_{_{\mathrm{Pl}}}^2k^{1/\beta }\hfill \\ \hfill \times \mathrm{cos}\left[\alpha (\eta )+\mathrm{ln}\left(\mathrm{\Gamma }[\frac{\beta 1}{2\beta },i\alpha (\eta )]\right)\right],\end{array}$$
(120)
where we have set $`D_11/k^2`$. Substituting the above expression in Eq. (VII), we get,
$`\left[k^3𝒫_S(k)\right]`$ $`=`$ $`C_0k^{2(\beta +2)}\left|1C_1k^{(1+1/\beta )}\right|^2`$
$`\times `$ $`\left[1+2\mathrm{cos}(x__H)2Im[]\mathrm{sin}(x__H)\right],`$
where $`C_1`$ depends on $`b_{_{11}}`$ and parameters of the power-law inflation.
In Fig. (1), we have plotted the standard and the trans-Planckian inflationary power spectra. As can be seen the trans-Planckian power spectrum has oscillations. We would like to caution the readers that the oscillations are small, here we have magnified the effect for illustrative purposes. We would also like to point the following: The power-spectrum (VII) we have obtained becomes significantly different compared to that of Ref. Martin and Brandenberger (2001) for very large $`k`$. For example, for a particular wave-vector this vanishes. However such an effect is not observable.
Following points are to be noted regarding the above results: (i) We have obtained the general power-spectrum (VII) of the scalar perturbations assuming that the scalar field is in the minimum energy state and that the contribution of the unit vector field to the energy density can be neglected. We have shown that the power spectrum depends on the form of the source term $`𝒮_k`$ which can be solved analytical in some particular limits. (ii) We have computed the power-spectrum of perturbations in a particular limit i. e. $`(1/c_10)`$. In this limit, we recover the result of Refs. Martin and Brandenberger (2001); Niemeyer and Parentani (2001).
## VIII Discussion and Conclusion
In this work, we have computed the gauge-invariant cosmological perturbation for the single scalar field inflation with the trans-Planckian effects introduced via the Jacobson-Corley dispersion relation. Even though the dispersion relation breaks the local Lorentz invariance, a covariant formulation of the corresponding theory can be carried out by introducing a unit time-like vector field.
Using the covariant Lagrangian, we have obtained the perturbed stress-tensor for the scalar and tensor perturbations around the FRW background. We have shown the following: (i) The non-linear effects introduce corrections to the perturbed energy density while the other components of the perturbed stress-tensor remains unchanged. Thus, for the trans-Planckian scenario, we have shown that $`\mathrm{\Phi }=\mathrm{\Psi }`$ and the constraint equation (17b) remains unchanged. (ii) The non-linear terms contributing to the stress-tensor are proportional to $`k^2`$ and hence in the super-Hubble scales, as expected, the contribution to the perturbed energy density can be ignored. (iii) The spatial higher derivative terms appear only in the equation of the motion of the perturbed inflation field ($`\delta \phi `$) while the speed of propagation of the perturbations \[in the equation of motion of the scalar perturbations ($`\mathrm{\Phi }`$)\] is different from that of the standard inflation. (iv) The speed of propagation of the perturbations ($`c_s^2`$) is different from that of the canonical single scalar field inflation. (v) The perturbations are not purely adiabatic. $`\xi `$ act as an extra scalar field during inflation and hence can act as a source. This introduces non-adiabatic (entropic) perturbations. (vi) Since, the trans-Planckian corrections do not change the pressure perturbations, the perturbation equations for the tensor modes do not change. Hence, the tensor perturbation equation remain unchanged. Recently, Lim Lim (2005) had show that general Lorentz violating models (with out taking in-account the higher derivatives of the scalar field) can modify the pressure perturbations and hence the tensor perturbation equations. However, in this model, this is not the case. Since the tensor perturbations remain the same, the well-know consistency relation between the scalar and tensor ratio will also be broken in this model Hui and Kinney (2002).
We combined Eqs. (55,56,57) to obtain a single differential equation of $`\mathrm{\Phi }`$. We have shown that the resultant differential equation of $`\mathrm{\Phi }`$ is different from that of the standard canonical scalar field driven inflation. More importantly, the differential equation of $`\mathrm{\Phi }`$ in our model is fourth order while in the standard canonical scalar field it is second order. We also obtained the solutions of $`\mathrm{\Phi }`$ in the three regimes for the power-law inflation.
We have also obtained the equation of motion of the Mukhanov-Sasaki variable for the perturbed inflaton field with higher derivatives. In all the earlier analyzes, the Mukhanov-Sasaki variable ($`Q`$) was assumed to satisfy the differential equation Eq. (62) in which the source term ($`𝒮_k`$) was assumed to be zero. More importantly, we had shown that the source term in Eq. (62) dominates during the trans-Planckian regime. The Mukhanov-Sasaki variable of the two fields are strongly coupled and hence obtaining the solution analytically is possible only in a particular limit.
In this work, we calculated the power-spectrum corresponding to the inflaton field, during the power-law inflation by assuming that (i) the quantum field $`\mu __S`$ is coupled to an external, classical source field $`𝒮_k(\eta )`$ which is determined by solving the coupled differential equations (55, 57) (ii) the quantum field is initially in a minimum energy state and (iii) $`d_1/M_{_{\mathrm{Pl}}}^2b_{_{11}}M_{_{\mathrm{Pl}}}^2`$. We have shown that in this particular limit, the power-spectrum is same as that obtained in Refs. Martin and Brandenberger (2001); Niemeyer and Parentani (2001).
The work suggests various possible directions for further study:
* We have obtained the power-spectrum analytically in the limit of $`d_1/M_{_{\mathrm{Pl}}}^2b_{_{11}}M_{_{\mathrm{Pl}}}^2`$. The trans-Planckian corrections in this limit are small to be observed in the present or the future CMB experiments. It would be interesting to obtain the power-spectrum by solving the system of differential equations numerically and obtain the leading order trans-Planckian corrections in these models.
* As we have mentioned earlier, this model introduces non-adiabatic perturbations which can lead to the non-Gaussianity in the CMB. Recently in Refs. Martin and Ringeval (2004a, b); Martin and Ringeval (2005); Easther et al. (2005a, b), trans-Planckian constraints from the CMB was studied in detail. It would be interesting to do a similar analysis for this scenario. The non-Gaussian signatures may place stringent and independent constraints on the parameters $`b_{_{11}},d_1`$.
* In this work, we have ignored the solenoidal part of the perturbed $`u`$ field. The solenoidal part contributes to the vector perturbations. It would be interesting to see whether the solenoidal part of the perturbed unit-time like vector field can lead to the growing large-scale vorticity and hence the production of large scale primordial magnetic field.
* In this work, we have ignored the back-reaction of the field excitations on the perturbed FRW background. There have been claims in the literature Tanaka (2000); Starobinsky (2001) that trans-Planckian modes may effect the evolution of cosmological fluctuations in the early stages of cosmological inflation in a non-trivial way. In Ref. Brandenberger and Martin (2005), the authors have discussed in detail the backreaction problem of the trans-Planckian inflation in a toy model and have shown that the back-reaction of the trans-Planckian modes may lead to a renormalization of the cosmological constant driving inflation.. It would be interesting to perform a similar analysis for this model.
We hope to return to study some of these issues in the near future.
## Acknowledgments
The authors wish to thank J. Martin, L. Sriramkumar for comments on the earlier version of the paper. The authors also wish to thank S. Bashinsky, U. Seljak and in particular, N. Bartolo for stimulating discussions. SS thanks D. Mattingly, B. van Tent for useful email correspondences.
## Appendix A The Background.
In this appendix, we give key steps in obtaining the stress-tensor corresponding to the two corrective Lagrangians (29a, 29b), and the equations of motion of the scalar field and the unit vector field. Having obtained these, we discuss their properties in the FRW background.
This question has been addressed in Ref. Lemoine et al. (2002). Our treatment differs with the one followed in that paper by the fact that we do not make the decomposition in time and space-like components. Our condensed formulas will prove very useful when computing the perturbations.
In order to do that, it proves easier to go back to the action. Let us first specialize to the contribution of the non-linear part of the Lagrangian (29a):
$$S_{_{\mathrm{cor}}}=b_{_{11}}d^4x\sqrt{g}(𝒟^2\phi )^2.$$
(122)
Using the definition given in Eq.(31), the variation of $`𝒟^2\phi `$ can be written as
$`\delta (𝒟^2\phi )`$ $`=`$ $`\stackrel{~}{A}_{\mu \nu }\delta g^{\mu \nu }+\stackrel{~}{B}_{\mu \nu }^\sigma _\sigma \delta g^{\mu \nu }+\stackrel{~}{C}^\mu \delta u_\mu +\stackrel{~}{D}^{\nu \mu }_\nu \delta u_\mu .`$ (123)
$`+\stackrel{~}{E}^\mu _\mu \delta \phi +\stackrel{~}{H}^{\mu \nu }_\mu _\nu \delta \phi .`$
The quantities $`\stackrel{~}{A}\mathrm{}\stackrel{~}{H}`$ can be written explicitly. For example,
$$\stackrel{~}{A}_{\mu \nu }=\frac{}{g^{\mu \nu }}𝒟^2\phi ,\stackrel{~}{H}^{\mu \nu }=\frac{}{(_\mu _\nu \phi )}𝒟^2\phi .$$
(124)
These partial derivatives have to be taken keeping in mind the choice of variables made in this work. In order to be consistent, we choose the following set of independent variables
$`\phi ,_\beta \phi ,_\alpha _\beta \phi ,_\sigma _\alpha _\beta \phi ,g^{\mu \nu },_\rho g^{\mu \nu },_\tau _\rho g^{\mu \nu },`$
$`u_\sigma ,_\beta u_\sigma ,_\alpha _\beta u_\sigma ,\lambda .`$ (125)
Using the relations
$$_\epsilon (g_{\varkappa \varrho }g^{\varrho \varsigma })=0;\frac{}{g_{\alpha \beta }}(g_{\varkappa \varrho }g^{\varrho \varsigma })=0,$$
(126)
we get,
$`{\displaystyle \frac{g^{\alpha \beta }}{g^{\rho \sigma }}}`$ $`=`$ $`\delta _\rho ^\alpha \delta _\sigma ^\beta ,{\displaystyle \frac{g_{\alpha \beta }}{g^{\rho \sigma }}}=g_{\rho \alpha }g_{\sigma \beta },`$
$`_ϵg_{\alpha \beta }`$ $`=`$ $`g_{\alpha \rho }g_{\beta \sigma }_ϵg^{\rho \sigma }.`$ (127)
The expressions giving the quantities $`\stackrel{~}{A},\mathrm{},\stackrel{~}{H}`$ are very long and time consuming to obtain and will not be displayed here. One of the most important aspects of the way we will present our results is that for our purposes we only need to know their values on the FRW background. This is the subject of the next appendix.
Substituting Eq. (123) in (122) and integrating by parts the resultant expression, we obtain
$`\delta S_{_{cor}}`$ $`=`$ $`b_{11}{\displaystyle }d^4x[\sqrt{g}({\displaystyle \frac{1}{2}}g_{\mu \nu }\left(𝒟^2\phi \right)^22\stackrel{~}{A}_{\mu \nu }𝒟^2\phi )`$
$`+_\sigma \left(\sqrt{g}𝒟^2\phi \stackrel{~}{B}_{\mu \nu }^\sigma \right)]\delta g^{\mu \nu }`$
$`2b_{11}{\displaystyle d^4x\left[\sqrt{g}𝒟^2\phi \stackrel{~}{C}^\mu _\nu \left(\sqrt{g}𝒟^2\phi \stackrel{~}{D}^{\nu \mu }\right)\right]\delta u_\mu }`$
$`2b_{11}{\displaystyle }d^4x[_\mu \left(\sqrt{g}𝒟^2\phi \stackrel{~}{E}^\mu \right)`$
$`+_\mu _\nu \left(\sqrt{g}𝒟^2\phi \stackrel{~}{H}^{\mu \nu }\right)]\delta \phi .`$
From the above expression, it is easy to infer the contribution of the corrective Lagrangian to the stress-energy tensor as well as to the field equations of the inflaton and the vector field $`u_\mu `$, i. e.,
$`T_{\mu \nu }^{\left(\mathrm{cor}\right)}`$ $`=`$ $`b_{_{11}}g_{\mu \nu }\left(𝒟^2\phi \right)^24b_{_{11}}\stackrel{~}{B}_{\mu \nu }^\rho _\rho \left(𝒟^2\phi \right)`$
$`+`$ $`b_{_{11}}\left(4\stackrel{~}{A}_{\mu \nu }2{\displaystyle \frac{_\rho g}{g}}\stackrel{~}{B}_{\mu \nu }^\rho 4_\rho \stackrel{~}{B}_{\mu \nu }^\rho \right)𝒟^2\phi `$
$``$ $`b_{_{11}}g_{_{\mu \nu }}\left(𝒟^2\phi \right)^2+b_{_{11}}E_{_{\mu \nu }}𝒟^2\phi +4b_{_{11}}C_{_{\mu \nu }}^{\rho \varkappa }_\varkappa \phi _\rho \left[𝒟^2\phi \right],`$
$`eq_{1,\phi }`$ $`=`$ $`{\displaystyle \frac{2b_{_{11}}}{\sqrt{g}}}([_\mu \left(\sqrt{g}\stackrel{~}{E}^\mu \right)_\mu _\nu \left(\sqrt{g}\stackrel{~}{H}^{\mu \nu }\right)]𝒟^2\phi `$ (130)
$`+`$ $`\left[\sqrt{g}\stackrel{~}{E}^\mu _\nu \left(\sqrt{g}\stackrel{~}{H}^{\mu \nu }\right)_\nu \left(\sqrt{g}\stackrel{~}{H}^{\nu \mu }\right)\right]_\mu 𝒟^2\phi `$
$`\sqrt{g}\stackrel{~}{H}^{\mu \nu }_\mu _\nu 𝒟^2\phi ),`$
$`eq_{2,\phi }`$ $`=`$ $`2b_{11}[(\stackrel{~}{C}^\mu {\displaystyle \frac{1}{2}}{\displaystyle \frac{_\nu g}{g}}\stackrel{~}{D}^{\nu \mu }_\nu \stackrel{~}{D}^{\nu \mu })𝒟^2\phi `$ (131)
$`\stackrel{~}{D}^{\nu \mu }_\nu 𝒟^2\phi ],`$
where (130) gives the RHS of Eq. (III.1) and (131) gives the RHS of Eq. (III.1).
Similarly, the variation of the action for the unit vector field (29b) leads to
$`\delta S_u`$ $`=`$ $`{\displaystyle d^4x\sqrt{g}T_{\mu \nu }\delta g^{\mu \nu }}2\lambda {\displaystyle d^4x\sqrt{g}g^{\mu \nu }u_\nu \delta u_\mu }`$ (132)
$`+`$ $`{\displaystyle d^4x\sqrt{g}(g^{\alpha \beta }u_\alpha u_\beta +1)\delta \lambda }`$
$``$ $`d_1{\displaystyle d^4x\sqrt{g}g^{\alpha \rho }g^{\beta \rho }\delta (F_{\alpha \beta }F_{\rho \sigma })}.`$
From the above expression, we obtain the contribution to the stress-energy tensor and the field equation of the vector field, i. e.,
$`T_{\mu \nu }^{(u)}`$ $`=`$ $`\left({\displaystyle \frac{1}{2}}\lambda (g^{\alpha \beta }u_\alpha u_\beta 1)+{\displaystyle \frac{1}{2}}d_1F_{\alpha \beta }F^{\alpha \beta }\right)g_{\mu \nu }`$ (133)
$``$ $`\lambda u_\mu u_\nu 2d_1g^{\alpha \rho }F_{\alpha \mu }F_{\rho \nu },`$
$`eq_{2,u}`$ $`=`$ $`2\lambda g^{\mu \nu }u_\nu 2d_1{\displaystyle \frac{1}{\sqrt{g}}}_\nu \left[\sqrt{g}F^{\mu \nu }\right],.`$ (134)
where (134) gives the LHS of Eq. (III.1).
It is easy to verify that the energy momentum tensor corresponding to the corrective Lagrangian and the vector field vanish while the equation of motion of the vector field and the inflaton are satisfied when the following relation holds
$$\overline{u}_\mu =a(\eta )(1,0,0,0),\overline{\lambda }=0,_\mu ^\mu \overline{\phi }=\frac{\overline{\phi }_{}^{^{}}{}_{}{}^{2}}{a^2(\eta )}.$$
(135)
The intermediary relations needed are
$`\overline{F}_{\mu \nu }=0,\overline{}_{\alpha \beta }=a^2\left(\eta _{\alpha \beta }\delta __\alpha ^0\delta __\beta ^0\right),𝒟^2\left(\overline{\phi }\right)=0.`$ (136)
The evolution of the modified scalar field in the FRW background is same as that of the inflaton. Hence, the trans-Planckian effects, in this model, only affect the perturbations. As we will see in the next appendix, the above equations play a significant role in discarding many terms from the perturbed stress-tensor.
## Appendix B The Perturbed FRW space time.
In this appendix, we obtain the linear order perturbed stress-tensor corresponding to the two corrective Lagrangians (29a, 29b), and the equations of motion of the perturbed scalar field and the unit vector field.
We will use more than once the Leibniz rule:
$`\delta (XY)`$ $`=`$ $`\overline{X}\delta Y+\overline{Y}\delta XX_o\delta Y+Y_o\delta X.`$ (137)
and the vanishing of the three dimensional Laplacian of the scalar field in the FRW background.
For example, the variation of the component of the stress-energy tensor coming from the non-linear part of the Lagrangian reads
$`\delta T_{\mu \nu }^{corr}`$ $`=`$ $`\left[4\stackrel{~}{A}_{\mu \nu }2{\displaystyle \frac{_\rho g}{g}}\stackrel{~}{B}_{\mu \nu }^\rho 4_\rho \stackrel{~}{B}_{\mu \nu }^\rho \right]_o\delta 𝒟^2\phi `$ (138)
$``$ $`4\left(\stackrel{~}{B}_{\mu \nu }^\rho \right)_o_\rho (\delta 𝒟^2\phi ).`$
From the above expression, we see that it is enough to know the quantities $`\stackrel{~}{A},\mathrm{},\stackrel{~}{H}`$ on the background. Going back to Eq. (123), we have to compute the variation of the vector field, the Christoffel symbol and other quantities. Let us begin by the connection coefficient:
$`\mathrm{\Gamma }_{\rho \sigma }^\alpha ={\displaystyle \frac{1}{2}}g^{\alpha \tau }(_\tau g_{\rho \sigma }+_\rho g_{\sigma \tau }+_\sigma g_{\tau \rho }).`$ (139)
Due to our choice of independent variables, we need to express the variations of the covariant components of the metric in terms of the contravariant ones. This is achieved simply:
$$\delta (g_{\alpha \beta }g^{\beta \gamma })=0\delta g_{\alpha \beta }=g_{\alpha \gamma }g_{\beta \sigma }\delta g^{\gamma \sigma }.$$
(140)
One then obtains
$`\delta \mathrm{\Gamma }_{\rho \sigma }^\alpha `$ $`=`$ $`aa^{}(\delta _\mu ^\alpha \delta _\nu ^0\eta _{\rho \sigma }+\delta _\mu ^\alpha \eta _{\sigma \nu }\delta _\rho ^0\delta _\mu ^\alpha \eta _{\rho \nu }\delta _\sigma ^0`$
$`+2\delta _0^\alpha \eta _{\rho \mu }\eta _{\sigma \nu }2\delta _\nu ^\alpha \delta _\rho ^0\eta _{\mu \sigma })\delta g^{\mu \nu }`$
$`+`$ $`{\displaystyle \frac{a^2}{2}}\left(\eta _{\rho \mu }\eta _{\sigma \nu }\eta ^{\alpha ϵ}\eta _{\sigma \mu }\delta _\nu ^\alpha \delta _\rho ^ϵ\eta _{\rho \nu }\delta _\mu ^\alpha \delta _\sigma ^ϵ\right)_ϵ\delta g^{\mu \nu }.`$
Similarly, we get
$`\delta ^{\rho \sigma }`$ $`=`$ $`(\delta _\mu ^\rho \delta _\nu ^\sigma +\delta _0^\sigma \delta _\mu ^\rho \delta _\nu ^0++\delta _0^\rho \delta _\mu ^\sigma \delta _\nu ^0)\delta g^{\mu \nu }`$ (142)
$`+`$ $`(\delta _0^\rho \eta ^{\sigma \mu }+\delta _0^\sigma \eta ^{\rho \mu })\delta u_\mu .`$
After a long but straightforward calculation, one obtains the result given in Eqs. (42,43 ) with the following expressions for the background:
$`(\stackrel{~}{A}_{\mu \nu })_o=(\overline{\phi }^{^{\prime \prime }}+3{\displaystyle \frac{a^{^{}}}{a}}\overline{\phi }{}_{}{}^{})\delta _\mu ^0\delta _\nu ^0`$ , $`(\stackrel{~}{B}_{\mu \nu }^ϵ)_o={\displaystyle \frac{1}{2}}\overline{\phi }^{^{}}\delta _\mu ^0\delta _\nu ^0\delta _0^ϵ,`$
$`(\stackrel{~}{C}^\mu )_o={\displaystyle \frac{1}{a^3}}\left(2\overline{\phi }^{^{\prime \prime }}+3{\displaystyle \frac{a^{^{}}}{a}}\overline{\phi }^{^{}}\right)\delta _0^\mu `$ , $`(\stackrel{~}{D}^{\sigma \rho })_o={\displaystyle \frac{1}{a^3}}\overline{\phi }^{^{}}\eta ^{\rho \sigma },`$
$`(\stackrel{~}{H})_o={\displaystyle \frac{1}{a^3}}(\eta ^{\mu \nu }+\delta _0^\mu \delta _0^\nu )`$ , $`(\stackrel{~}{E})_o=0.`$ (143)
At this point it is worth noticing that $`00`$ is the only non-zero component of the stress-tensor. Since the trans-Planckian corrections do not change the pressure perturbations, the perturbation equation for the tensor modes do not change. Hence, the tensor perturbation equations is given by Eq. (26). Recently, Lim Lim (2005) had show that general Lorentz violating models (with out taking into account higher derivatives of the scalar field) can modify the pressure perturbations and hence the tensor perturbation equations. However, in our specific Lorentz violating model, this is not the case.
Let us now specialize to the scalar perturbations. The covariant components read
$$\delta g_{\mu \nu }=a^2(\eta )\left(\begin{array}{ccc}2\varphi & & _iB\\ _iB& & 2\psi \delta ^{ij}2_i_jE\end{array}\right).$$
(144)
Using the identities given in Eq. (126), we get the contravariant components
$`\delta g^{00}={\displaystyle \frac{2\varphi }{a^2\left(\eta \right)}},\delta g^{0i}={\displaystyle \frac{_iB}{a^2\left(\eta \right)}},\delta g^{ij}={\displaystyle \frac{\left(2_i_jE2\psi \delta _{ij}\right)}{a^2\left(\eta \right)}}.`$ (145)
Using the fact that $`u_\mu `$ is a unit vector, we have
$$\delta (g^{\alpha \beta }u_\alpha u_\beta )=0\delta u_0=a\varphi .$$
(146)
Using the form of the scalar perturbations Eq. (144) and the results obtained in Eq.(143), we obtain the following simple formula
$$\delta (𝒟^2\phi )=\frac{1}{a^2}\left(\frac{1}{a}\overline{\phi }^{^{}}\delta ^{ij}_i\delta u_j+^2\delta \phi \right).$$
(147)
Let us now turn our attention to the equation of motion obeyed by the perturbations. Concerning the inflaton field, the first two contributions to the formula given in Eq. (130) vanish and the result is
$$\delta eq_{1,\phi }=\frac{2b_{11}}{a^4}\left[^4(\delta \phi )+\frac{1}{a}\overline{\phi }^{^{}}^2(\delta ^{ij}_i\delta u_j)\right].$$
(148)
In the same way, we get
$`b_{11}{\displaystyle \frac{1}{a^6}}\delta _0^\mu \left(10{\displaystyle \frac{a^{^{}}}{a}}(\overline{\phi }^{^{}})^2\delta ^{ij}_i\delta u_j(2\overline{\phi }^{^{\prime \prime }}+8{\displaystyle \frac{a^{^{}}}{a}}\overline{\phi }^{^{}})a^2\delta \phi \right)`$ (149)
$`+`$ $`b_{11}{\displaystyle \frac{1}{a^6}}\eta ^{\mu \nu }\overline{\phi }^{^{}}(2a^2_\nu \delta \phi 2\overline{\phi }^{^{}}_\nu \delta ^{ij}_i\delta u_j)2{\displaystyle \frac{1}{a}}\delta _0^\mu \delta \lambda `$
$`+`$ $`4d_1{\displaystyle \frac{1}{a^4}}(\eta ^{\rho \nu }\eta ^{\sigma \mu }\eta ^{\rho \mu }\eta ^{\sigma \nu })_\nu _\rho \delta u_\sigma =0.`$
From the temporal index ($`\mu =0`$), we obtain the variation of the Lagrange multiplier, i. e.,
$`\delta \lambda `$ $`=`$ $`2d_1{\displaystyle \frac{1}{a^3}}\delta ^{ij}_j(\delta u_i)^{^{}}+{\displaystyle \frac{b_{11}}{a^5}}[a\overline{\phi }^{}^2\delta \phi ^{}\overline{\phi }^2\delta ^{ij}_i\delta u_j^{}`$ (150)
$`+5\overline{\phi }^2\delta ^{ij}_j\delta u_ia(\overline{\phi }^{\prime \prime }+4\overline{\phi }^{^{}})^2\delta \phi ].`$
From the spatial indices $`\mu =k`$, we obtain
$`b_{11}\left(2a\overline{\phi }^{^{}}^2_k\delta \phi +2(\overline{\phi }^{^{}})^2\delta ^{ij}_k_i\delta u_j\right)`$
$`+`$ $`4d_1a^2\left[\delta u_k^{\prime \prime }+^2\delta u_k+_k(a\varphi )^{^{}}_k_i\delta u_i\right]=0.`$
Having these results, especially the variation of the Lagrange multiplier, we obtain the expression for the non-vanishing component of the stress-tensor:
$`\delta T_{00}^{(\mathrm{cor})}`$ $`=`$ $`{\displaystyle \frac{2b_{11}}{a^2}}(5{\displaystyle \frac{}{a}}\overline{\phi }^2_i\delta u_i{\displaystyle \frac{1}{a}}\overline{\phi }^2_i\delta u_i^{}`$
$`(\overline{\phi }^{\prime \prime }+4\overline{\phi }^{})^2\delta \phi +\overline{\phi }^{}^2\delta \phi ^{}),`$
$`\delta T_{00}^u`$ $`=`$ $`2a^2\delta \lambda .`$ (153)
We now wish to express all the quantities in terms of gauge invariant quantities. For the vector field, we have
$$\delta u_i=\delta u_i^{(gi)}u_0_i(BE^{^{}}).$$
(154)
The equation of the vector field now becomes
$`b_{_{11}}\left(a\overline{\phi }^{}^2_k\delta \phi ^{\left(gi\right)}\overline{\phi }_{}^{}{}_{}{}^{2}\delta ^{ij}_k_i\delta u_j^{\left(gi\right)}\right)`$ (155)
$`+2d_1a^2\left(\delta u_{k}^{\left(gi\right)}{}_{}{}^{\prime \prime }+^2\delta u_k^{\left(gi\right)}+_k\left(a\mathrm{\Phi }\right)^{}+_k_i\delta u_i^{\left(gi\right)}\right)=0,`$
while for the inflaton one obtains
$`2V_{,\phi }a^2\mathrm{\Phi }4\overline{\phi }^{^{}}\mathrm{\Phi }^{^{}}+a^2V_{,\phi \phi }\delta \phi ^{\left(gi\right)}+2\delta \phi _{}^{\left(gi\right)}{}_{}{}^{}+\delta \phi _{}^{\left(gi\right)}{}_{}{}^{\prime \prime }`$ (156)
$`^2\delta \phi ^{\left(gi\right)}+{\displaystyle \frac{2b_{11}}{a^2}}\left[^4\delta \phi ^{\left(gi\right)}+{\displaystyle \frac{1}{a}}\overline{\phi }^{^{}}^2\delta ^{ij}_i\delta u_j^{\left(gi\right)}\right]=0.`$
Similarly, we can obtain expressions for the stress-energy tensor (53) interms of the gauge-invariant variables.
Until this point we have not assumed any specific form of the spatial part of the perturbed unit-vector field. In general, the time-dependent spatial components of the perturbed $`u`$ field can be expressed as a sum of irrotational and solenoidal parts, i. e.
$$\delta u_i(\xi )_i+(\times \varpi )_i.$$
(157)
The irrotational part of the perturbed $`u`$ field will contribute to the scalar perturbations while the solenoidal part of the perturbed $`u`$ field contributes to the vector perturbations. Since we are interested in the scalar perturbations, for the rest of the calculations, we ignore the solenoidal part. Thus, we get
$$\delta u_i^{(gi)}=_i\xi ^{(gi)}.$$
(158)
Substituting this in the Einstein and the field equations, we obtain Eqs. (53).
## Appendix C The fourth order equation
In this appendix, we provide key steps in obtaining the equation of motion of Bardeen potential ($`\mathrm{\Phi }`$), in a general FRW background, by combining Eqs. (55,56,57).
From the constraint equation we deduce that
$$\delta \phi _k=\frac{2}{3}M_{pl}^2\left(\frac{1}{\overline{\phi }^{^{}}}\mathrm{\Phi }_k^{^{}}+\frac{}{\overline{\phi }^{^{}}}\mathrm{\Phi }_k\right).$$
(159)
Substituting this expression and its derivatives into the equation of motion of $`\delta \phi _k`$, one obtains $`\xi _k`$, i. e.,
$$\xi _k=\frac{M_{pl}}{k^2}\left(Q_3\mathrm{\Phi }_k^{^{\prime \prime \prime }}+Q_2\mathrm{\Phi }_k^{^{\prime \prime }}+Q_1\mathrm{\Phi }_k^{^{}}+Q_0\mathrm{\Phi }_k\right),$$
(160)
where
$`Q_3`$ $`=`$ $`{\displaystyle \frac{M_{pl}^2}{3b_{11}}}{\displaystyle \frac{a^3}{\overline{\phi }^{}_{}{}^{}2}}{\displaystyle \frac{1}{k^2}},Q_2={\displaystyle \frac{M_{pl}^2}{3b_{11}}}{\displaystyle \frac{a^3}{\overline{\phi }^{^{}}}}\left[2\left({\displaystyle \frac{1}{\overline{\phi }^{^{}}}}\right)^{^{}}+{\displaystyle \frac{}{\overline{\phi }^{^{}}}}\right]{\displaystyle \frac{1}{k^2}},`$
$`Q_1`$ $`=`$ $`\left[{\displaystyle \frac{M_{pl}^2}{3b_{11}}}{\displaystyle \frac{a^3}{\overline{\phi }^{^{}}}}\left(\left({\displaystyle \frac{1}{\overline{\phi }^{^{}}}}\right)^{^{\prime \prime }}+2\left({\displaystyle \frac{}{\overline{\phi }^{^{}}}}\right)^{^{}}2{\displaystyle \frac{1}{b_{11}}}a^3\right)\right]{\displaystyle \frac{1}{k^2}},`$
$`+`$ $`{\displaystyle \frac{M_{pl}^2}{3b_{11}}}{\displaystyle \frac{a^3}{\overline{\phi }^{}_{}{}^{}2}}+{\displaystyle \frac{2}{3}}M_{pl}^2{\displaystyle \frac{a}{\overline{\phi }^{}_{}{}^{}2}}k^2,`$
$`Q_0`$ $`=`$ $`\left[{\displaystyle \frac{M_{pl}^2}{3b_{11}}}{\displaystyle \frac{a^3}{\overline{\phi }^{^{}}}}\left({\displaystyle \frac{}{\overline{\phi }^{^{}}}}\right)^{^{\prime \prime }}+{\displaystyle \frac{1}{b_{11}}}{\displaystyle \frac{a^5}{\overline{\phi }^{^{}}}}\overline{V}_{,\phi }\right]{\displaystyle \frac{1}{k^2}}+{\displaystyle \frac{M_{pl}^2}{3b_{11}}}{\displaystyle \frac{a^3}{\overline{\phi }^{}_{}{}^{}2}}`$ (161)
$`+`$ $`{\displaystyle \frac{2}{3}}M_{pl}^2{\displaystyle \frac{a}{\overline{\phi }^{}_{}{}^{}2}}k^2.`$
Substituting the above expressions in the field equation of the Bardeen potential, we get,
$`\mathrm{\Phi }_k^{^{\prime \prime \prime \prime }}+X_3\mathrm{\Phi }_k^{^{\prime \prime \prime }}+(X_2+Y_2k^2+Z_2k^4)\mathrm{\Phi }_k^{^{\prime \prime }}`$
$`+`$ $`(X_1+Y_1k^2+Z_1k^4)\mathrm{\Phi }_k^{^{}}+(X_0+Y_0k^2+Z_0k^4)\mathrm{\Phi }_k=0,`$
where
$`X_3`$ $`=`$ $`{\displaystyle \frac{\overline{\phi }^{}_{}{}^{}2}{a^3}}\left({\displaystyle \frac{a^3}{\overline{\phi }^{}_{}{}^{}2}}\right)^{^{}}+\overline{\phi }^{^{}}\left[2\left({\displaystyle \frac{1}{\overline{\phi }^{^{}}}}\right)^{^{}}+{\displaystyle \frac{}{\overline{\phi }^{^{}}}}\right],`$
$`X_2`$ $`=`$ $`{\displaystyle \frac{\overline{\phi }^{}_{}{}^{}2}{a^3}}\left[{\displaystyle \frac{a^3}{\overline{\phi }^{^{}}}}\left(2\left({\displaystyle \frac{1}{\overline{\phi }^{^{}}}}\right)^{^{}}+{\displaystyle \frac{}{\overline{\phi }^{^{}}}}\right)\right]^{^{}}`$
$`+`$ $`\overline{\phi }^{^{}}\left[\left({\displaystyle \frac{1}{\overline{\phi }^{^{}}}}\right)^{^{\prime \prime }}+2\left({\displaystyle \frac{}{\overline{\phi }^{^{}}}}\right)^{^{}}\right]6{\displaystyle \frac{1}{M_{pl}^2}}\overline{\phi }^{}_{}{}^{}2,`$
$`Y_2`$ $`=`$ $`{\displaystyle \frac{3}{2}}{\displaystyle \frac{b_{11}}{d_1}}{\displaystyle \frac{\overline{\phi }^{}_{}{}^{}2}{a^3}}+1,Z_2=2b_{11}{\displaystyle \frac{1}{a^2}},`$
$`X_1`$ $`=`$ $`{\displaystyle \frac{\overline{\phi }^{}_{}{}^{}2}{a^3}}\left[{\displaystyle \frac{a^3}{\overline{\phi }^{}}}\left(\left({\displaystyle \frac{1}{\overline{\phi }^{^{}}}}\right)^{^{\prime \prime }}+2\left({\displaystyle \frac{}{\overline{\phi }^{^{}}}}\right)^{^{}}\right)6{\displaystyle \frac{1}{M_{pl}^2}}a^3\right]^{^{}},`$
$`Y_1`$ $`=`$ $`{\displaystyle \frac{\overline{\phi }^{}_{}{}^{}2}{a^3}}\left({\displaystyle \frac{a^3}{\overline{\phi }^{}_{}{}^{}2}}\right)^{^{}}++3{\displaystyle \frac{b_{11}}{d_1}}{\displaystyle \frac{\overline{\phi }^{}_{}{}^{}2}{a^3}}\left({\displaystyle \frac{\overline{\phi }^{^{\prime \prime }}}{\overline{\phi }^{^{}}}}\right),`$
$`Z_1`$ $`=`$ $`2b_{11}{\displaystyle \frac{\overline{\phi }^{}_{}{}^{}2}{a^3}}\left[\left({\displaystyle \frac{a}{\overline{\phi }^{}_{}{}^{}2}}\right)^{^{}}+{\displaystyle \frac{a}{\overline{\phi }^{}_{}{}^{}2}}\right],`$ (163)
$`X_0`$ $`=`$ $`{\displaystyle \frac{\overline{\phi }^{}_{}{}^{}2}{a^3}}\left[{\displaystyle \frac{a^3}{\overline{\phi }^{^{}}}}\left({\displaystyle \frac{}{\overline{\phi }^{^{}}}}\right)^{^{\prime \prime }}+3{\displaystyle \frac{1}{M_{pl}^2}}{\displaystyle \frac{a^5}{\overline{\phi }^{^{}}}}\overline{V}_{,\phi }\right]^{^{}},`$
$`Y_0`$ $`=`$ $`{\displaystyle \frac{\overline{\phi }^{}_{}{}^{}2}{a^3}}\left({\displaystyle \frac{a^3}{\overline{\phi }^{}_{}{}^{}2}}\right)^{^{}}+3{\displaystyle \frac{b_{11}}{d_1}}{\displaystyle \frac{\overline{\phi }^{}_{}{}^{}2}{a^3}}\left(^{^{}}{\displaystyle \frac{\overline{\phi }^{^{\prime \prime }}}{\overline{\phi }^{^{}}}}\right),`$
$`Z_0`$ $`=`$ $`2b_{11}{\displaystyle \frac{\overline{\phi }_{}^{}{}_{}{}^{2}}{a^3}}\left({\displaystyle \frac{a^{}}{\overline{\phi }_{}^{}{}_{}{}^{2}}}\right)^{^{}}+{\displaystyle \frac{3}{2}}b_{11}\left({\displaystyle \frac{1}{d_1}}2{\displaystyle \frac{1}{M_{pl}^2}}\right){\displaystyle \frac{\overline{\phi }^{}_{}{}^{}2}{a^3}}.`$ |
warning/0507/cond-mat0507502.html | ar5iv | text | # Low frequency Rabi spectroscopy for a dissipative two-level system
<sup>1</sup><sup>1</sup>institutetext: Novosibirsk State Technical University, 20 K. Marx Ave., 630092 Novosibirsk, Russia
Institute for Physical High Technology, P.O. Box 100239, D-07702 Jena, Germany
Decoherence; open systems; quantum statistical methods Quantum well devices Quantum interference devices
## Abstract
We have analyzed the interaction of a dissipative two level quantum system with high and low frequency excitation. The system is continuously and simultaneously irradiated by these two waves. If the frequency of the first signal is close to the level separation the response of the system exhibits undamped low frequency oscillations whose amplitude has a clear resonance at the Rabi frequency with the width being dependent on the damping rates of the system. The method can be useful for low frequency Rabi spectroscopy in various physical systems which are described by a two level Hamiltonian, such as nuclei spins in NMR, double well quantum dots, superconducting flux and charge qubits, etc. As the examples, the application of the method to a nuclear spin and to the readout of a flux qubit are briefly discussed.
It is well known that under resonant irradiation a quantum two level system (TLS) can undergo coherent (Rabi) oscillations. The frequency of these oscillations is proportional to the amplitude of the resonant field and is much lower than the gap frequency of TLS. The effect is widely used in molecular beam spectroscopy , and in quantum optics .
During the last several years it has been proven experimentally that Rabi spectroscopy can serve as a valuable tool for the determination of relaxation times in solid state quantum mechanical two-level systems, qubits, to be used for quantum information processing . These systems normally are strongly coupled to the environment, which results in the fast damping of Rabi oscillations. It prevents the use of conventional continuous measurements schemes for their detection, though the special schemes for the detection of coherent oscillations through a weak continuous measurement of a TLS were proposed in . That is why Rabi oscillations are measured with the pulse technique through the statistic of switching events of the occupation probability between two energy levels with excitation and read out being taken at the gap frequency of TLS, which normally, lies in GHz range . The main drawback of this technique is that it requires sophisticated high frequency readout electronics.
In this work we propose a new experimental method for a Rabi spectroscopy of a TLS, when readout electronics is continuously swept across the low (compare to the gap) Rabi resonance. The method is rather general and can be applied to a great variety of two level systems.
We consider a TLS which is irradiated continuously by two external sources. The first with a frequency $`\omega _0`$, which is close to the energy gap between the two levels, excites the low frequency Rabi oscillations. Normally, Rabi oscillations are damped out with a rate, which is dependent on how strongly the system is coupled to the environment. However, if a second low frequency source is applied simultaneously to TLS it responds with persistent low frequency oscillations. The amplitude of these low frequency oscillations has a resonance at the Rabi frequency with the width being dependent on the damping rates of the system. Note, that this approach has a well known classical analog. Indeed, a damping rate of a classical oscillator can be easily obtained from its amplitude-frequency characteristics.
We start with a Hamiltonian of a driven TLS , which is subjected to both high and low frequency excitation:
$$H=\frac{\mathrm{\Delta }}{2}\sigma _x+\frac{\epsilon }{2}\sigma _z\sigma _zF\mathrm{cos}\omega _0t\sigma _zG(t).$$
(1)
Here the first two terms describe an isolated TLS, which can model a great variety of situations in physics and chemistry: from a spin 1/2 particle in a magnetic field to superconducting flux and charge qubits , . In order to be exact we consider the first two terms in (1) to describe a double-well system where only the ground states of the two wells are occupied, with $`\mathrm{\Delta }`$ being the energy splitting of a symmetric ($`\epsilon =0`$) TLS due to quantum tunnelling between two wells. The quantity $`\epsilon `$ is the bias, the external energy parameter which makes the system asymmetric. The last two terms in (1) describe the interaction with external time-dependent high frequency, $`F`$, and low frequency, $`G`$ fields which modulate the energy asymmetry between the two wells.
Hamiltonian (1) is written in the localized state basis, i.e., in the basis of states localized in each well. In terms of the eigenstates basis, which we denote by upper case subscripts for the Pauli matrices $`\sigma _X,\sigma _Y,\sigma _Z`$, Hamiltonian (1) reads:
$`H=\left[{\displaystyle \frac{\mathrm{\Delta }}{\mathrm{\Delta }_\epsilon }}F\mathrm{cos}\omega _0t+{\displaystyle \frac{\mathrm{\Delta }}{\mathrm{\Delta }_\epsilon }}G(t)\right]\sigma _X`$
$`+\left[{\displaystyle \frac{\mathrm{\Delta }_\epsilon }{2}}{\displaystyle \frac{\epsilon }{\mathrm{\Delta }_\epsilon }}F\mathrm{cos}\omega _0t{\displaystyle \frac{\epsilon }{\mathrm{\Delta }_\epsilon }}G(t)\right]\sigma _Z,`$ (2)
where $`\mathrm{\Delta }_\epsilon =\sqrt{\mathrm{\Delta }^2+\epsilon ^2}`$ is the gap between two energy states.
The inclusion of the dissipative environment in Hamiltonian (1), results in the Bloch-Redfield equations for the matrix operators $`\sigma _X`$, $`\sigma _Y`$, $`\sigma _Z`$, , . For weak driving ($`F\mathrm{\Delta }`$), and for weak coupling of the TLS to the bath these equations can be approximated by Bloch-type equations :
$$\dot{\sigma }_Z=\left(2f\mathrm{cos}\omega _0t+2g(t)\right)\sigma _Y\mathrm{\Gamma }_Z\left(\sigma _ZZ_0\right),$$
(3)
$`\dot{\sigma }_Y=\left(2f\mathrm{cos}\omega _0t+2g(t)\right)\sigma _Z`$
$`+\left[{\displaystyle \frac{\mathrm{\Delta }_\epsilon }{\mathrm{}}}{\displaystyle \frac{2\epsilon }{\mathrm{\Delta }}}f\mathrm{cos}\omega _0t{\displaystyle \frac{2\epsilon }{\mathrm{\Delta }}}g(t)\right]\sigma _X\mathrm{\Gamma }\sigma _Y,`$ (4)
$$\dot{\sigma }_X=\left[\frac{\mathrm{\Delta }_\epsilon }{\mathrm{}}\frac{2\epsilon }{\mathrm{\Delta }}f\mathrm{cos}\omega _0t\frac{2\epsilon }{\mathrm{\Delta }}g(t)\right]\sigma _Y\mathrm{\Gamma }\sigma _X,$$
(5)
where $`f=\mathrm{\Delta }`$F$`/\mathrm{}\mathrm{\Delta }_\epsilon `$, $`g(t)=\mathrm{\Delta }`$G(t)$`/\mathrm{}\mathrm{\Delta }_\epsilon `$, and $`Z_0=\mathrm{tanh}\left(\mathrm{\Delta }_\epsilon /k_BT\right)`$ is the equilibrium polarization of the system in the absence of external excitation sources ($`f=0`$, $`g=0`$). The angled brackets in Eqs. (3), (Low frequency Rabi spectroscopy for a dissipative two-level system), and (5) denote the trace over reduced density matrix $`\rho (t)`$, which is obtained by tracing out all environment degrees of freedom: $`\sigma _X=Tr(\sigma _X\rho (t))`$, etc. In order to simplify the problem we assume the relaxation, $`\mathrm{\Gamma }_Z`$ and dephasing, $`\mathrm{\Gamma }`$, rates in (3)-(5) are time-independent, i.e. the rates are slowly varying functions on the scale of Rabi period which is of the order of $`1/f`$.
Assuming that a high frequency driving amplitude $`F`$ is sufficiently small we write the desired solution of Eqs. (3), (Low frequency Rabi spectroscopy for a dissipative two-level system), and (5) as:
$$\sigma _Z=Z(t),$$
(6)
$$\sigma _Y=Y(t)+A(t)\mathrm{cos}(\omega _0t)+B(t)\mathrm{sin}(\omega _0t),$$
(7)
$$\sigma _X=X(t)+C(t)\mathrm{cos}(\omega _0t)+D(t)\mathrm{sin}(\omega _0t),$$
(8)
where $`X`$, $`Y`$, $`Z`$, $`A`$, $`B`$, $`C`$, and $`D`$ are slowly varying, as compared to the high frequency $`\omega _0`$, quantities.
As is known, a two-level system resonantly irradiated with a high frequency undergoes a low frequency Rabi oscillations. However, if the external low frequency excitation, $`G(t)`$ is absent $`(g=0)`$, the Rabi oscillations are damped out. For this case we obtain at the degeneracy point ($`\epsilon =0`$) the following solution: $`X=0`$, $`Y=0`$, $`C=B`$, $`D=A`$,
$$Z(t)=z^{}+z_0e^{\mathrm{\Gamma }_1t}+e^{\mathrm{\Gamma }_2t}\left(z_1\mathrm{cos}\omega _Rt+z_2\mathrm{sin}\omega _Rt\right),$$
(9)
$$A(t)=a^{}+a_0e^{\mathrm{\Gamma }_1t}+e^{\mathrm{\Gamma }_2t}\left(a_1\mathrm{cos}\omega _Rt+a_2\mathrm{sin}\omega _Rt\right),$$
(10)
$$B(t)=b^{}+b_0e^{\mathrm{\Gamma }_1t}+e^{\mathrm{\Gamma }_2t}\left(b_1\mathrm{cos}\omega _Rt+b_2\mathrm{sin}\omega _Rt\right),$$
(11)
where the relaxation rates, $`\mathrm{\Gamma }_1`$, $`\mathrm{\Gamma }_2`$, and the Rabi frequency $`\omega _R`$ are determined by the equation
$$\left(\lambda i\mathrm{\Gamma }\right)^2\left(\lambda i\mathrm{\Gamma }_Z\right)\left(\lambda i\mathrm{\Gamma }\right)f^2\left(\lambda i\mathrm{\Gamma }_Z\right)\delta ^2=0.$$
(12)
Here $`\lambda `$ is the eigenvalue for the oscillation mode, e. g., $`Z(t)e^{i\lambda t}`$. This cubic equation has simple analytic solutions only for two cases. Firstly, in the absence of a relaxation in Eqs. (3) - (5) $`(\mathrm{\Gamma }_Z=\mathrm{\Gamma }=0)`$ we obtain $`\mathrm{\Gamma }_1=\mathrm{\Gamma }_2=0`$, $`\omega _R=\sqrt{\delta ^2+f^2}`$, where $`\delta `$ is the high frequency detuning parameter, $`\delta =\omega _0\mathrm{\Delta }_\epsilon /\mathrm{}`$. Secondly, for $`\delta =0`$ we obtain $`\mathrm{\Gamma }_1=\mathrm{\Gamma }`$, $`\mathrm{\Gamma }_2=(\mathrm{\Gamma }+\mathrm{\Gamma }_Z)/2`$, $`\omega _R=\sqrt{f^2(\mathrm{\Gamma }\mathrm{\Gamma }_Z)^2/4}`$. For the subsequent derivation we need only the quantities $`z^{}`$, $`a^{}`$, $`b^{}`$ which can be determined as the steady state solution of the equations (3), (Low frequency Rabi spectroscopy for a dissipative two-level system), and (5): $`z^{}=P_ZZ(0)`$, $`a^{}=z^{}f\mathrm{\Gamma }/(\mathrm{\Gamma }^2+\delta ^2)`$, $`b^{}=z^{}f\delta /(\mathrm{\Gamma }^2+\delta ^2)`$, where $`P_Z=\frac{\mathrm{\Gamma }_Z\left(\mathrm{\Gamma }^2+\delta ^2\right)}{\mathrm{\Gamma }_Z\left(\mathrm{\Gamma }^2+\delta ^2\right)+f^2\mathrm{\Gamma }}`$.
The quantity $`P_ZZ_0`$ is the nonequilibrium polarization, i. e., the steady state difference of occupation probabilities between two energy levels in the case when the high frequency excitation is applied to the TLS. Therefore, in the absence of the low frequency excitation the Rabi oscillations decay with a rate $`\mathrm{\Gamma }_2`$ given by Eq. (12). The main goal of our investigation is to obtain the persistent oscillations of the quantities $`Z(t)`$ (6), $`Y(t)`$ (7), $`X(t)`$ (8) which can be detected with low frequency (compared to the gap) electronic circuitry. In what follows we show that a low frequency signal applied to TLS can sustain the persistent low frequency oscillations of the above mentioned quantities. The amplitude of these oscillations has a clear resonance at the Rabi frequency with the width of the resonance being dependent on the damping rates of the system.
The Eqs. (3), (Low frequency Rabi spectroscopy for a dissipative two-level system), and (5) are analized in the presence of a low frequency excitation $`G(t)`$. We insert Eqs. (6), (7), (8) into (3), (Low frequency Rabi spectroscopy for a dissipative two-level system), and (5), and in accordance with the ideology of rotating wave approximation retain only the low frequency terms. We consider the force $`G`$ to be small, and therefore neglect the terms which are of second and higher order in $`G`$. In addition, we keep only the terms which oscillate within the bandwidth of the Rabi frequency and neglect the terms which are of the order of $`\mathrm{}f/\mathrm{\Delta }_\epsilon `$, $`\mathrm{}\mathrm{\Gamma }/\mathrm{\Delta }_\epsilon `$, $`\mathrm{}\mathrm{\Gamma }_Z/\mathrm{\Delta }_\epsilon `$. Consequently we obtain the following set of equations for low frequency quantities:
$$\dot{Z}=fA\mathrm{\Gamma }_ZZ$$
(13)
$$\dot{Y}=\frac{\mathrm{\Delta }_\epsilon }{\mathrm{}}X\mathrm{\Gamma }Y+2gz^{}\frac{\epsilon }{\mathrm{\Delta }}fB$$
(14)
$$\dot{X}=\frac{\mathrm{\Delta }_\epsilon }{\mathrm{}}Y\mathrm{\Gamma }X+\frac{\epsilon }{\mathrm{\Delta }}fA$$
(15)
$`\ddot{A}+2\mathrm{\Gamma }\dot{A}+(\mathrm{\Omega }_R^2+\mathrm{\Gamma }^2)A=`$
$`f\left(\mathrm{\Gamma }_Z\mathrm{\Gamma }\right)Z\mathrm{\Gamma }{\displaystyle \frac{2\epsilon }{\mathrm{\Delta }}}gb^{}\delta {\displaystyle \frac{2\epsilon }{\mathrm{\Delta }}}ga^{}{\displaystyle \frac{2\epsilon }{\mathrm{\Delta }}}\dot{g}b^{}`$ (16)
$`\ddot{B}+2\mathrm{\Gamma }\dot{B}+\left(\delta ^2+\mathrm{\Gamma }^2\right)B=`$
$`\delta fZ+\mathrm{\Gamma }{\displaystyle \frac{2\epsilon }{\mathrm{\Delta }}}ga^{}\delta {\displaystyle \frac{2\epsilon }{\mathrm{\Delta }}}gb^{}+{\displaystyle \frac{2\epsilon }{\mathrm{\Delta }}}\dot{g}a^{}`$ (17)
From these equations we can obtain the Fourier components for the slowly varying quantities $`A(\omega )`$ and $`B(\omega )`$, and for the low frequency persistent response of the TLS to a small low frequency excitation, $`\stackrel{~}{Z}(\omega )`$, $`\stackrel{~}{Y}(\omega )`$, and $`\stackrel{~}{X}(\omega )`$: $`\stackrel{~}{C}(\omega )=\stackrel{~}{B}(\omega )`$, $`\stackrel{~}{D}(\omega )=\stackrel{~}{A}(\omega )`$,
$$\stackrel{~}{Z}(\omega )=\stackrel{~}{g}(\omega )\frac{2\epsilon }{\mathrm{\Delta }}f^2P_ZZ_0\frac{\delta }{\delta ^2+\mathrm{\Gamma }^2}\frac{2\mathrm{\Gamma }+i\omega }{s(\omega )},$$
(18)
$$\stackrel{~}{A}(\omega )=\stackrel{~}{g}(\omega )\frac{2\epsilon }{\mathrm{\Delta }}fP_ZZ_0\frac{\delta }{\delta ^2+\mathrm{\Gamma }^2}\frac{\left(2\mathrm{\Gamma }+i\omega \right)\left(i\omega +\mathrm{\Gamma }_Z\right)}{s(\omega )},$$
(19)
$`\stackrel{~}{B}(\omega )=\stackrel{~}{g}(\omega ){\displaystyle \frac{2\epsilon }{\mathrm{\Delta }}}fP_ZZ_0{\displaystyle \frac{\delta }{\delta ^2+\mathrm{\Gamma }^2}}{\displaystyle \frac{1}{\left(\delta ^2+(\mathrm{\Gamma }+i\omega )^2\right)}}\left\{{\displaystyle \frac{f^2\delta \left(2\mathrm{\Gamma }+i\omega \right)}{s(\omega )}}{\displaystyle \frac{\delta ^2\mathrm{\Gamma }^2i\omega \mathrm{\Gamma }}{\delta }}\right\},`$ (20)
$$\stackrel{~}{Y}(\omega )=\frac{\mathrm{}}{\mathrm{\Delta }_\epsilon }\frac{\epsilon }{\mathrm{\Delta }}f\stackrel{~}{A}(\omega )+\left(\frac{\mathrm{}}{\mathrm{\Delta }_\epsilon }\right)^2\left(2\stackrel{~}{g}(\omega )P_ZZ_0\frac{\epsilon }{\mathrm{\Delta }}f\stackrel{~}{B}(\omega )\right)\left(\mathrm{\Gamma }+i\omega \right),$$
(21)
$$\stackrel{~}{X}(\omega )=\left(\frac{\mathrm{}}{\mathrm{\Delta }_\epsilon }\right)^2\frac{\epsilon }{\mathrm{\Delta }}f\stackrel{~}{A}(\omega )\left(\mathrm{\Gamma }+i\omega \right)\frac{\mathrm{}}{\mathrm{\Delta }_\epsilon }\left(2\stackrel{~}{g}(\omega )P_ZZ_0\frac{\epsilon }{\mathrm{\Delta }}f\stackrel{~}{B}(\omega )\right),$$
(22)
where $`s(\omega )=(\mathrm{\Omega }_R\omega +i\mathrm{\Gamma })(\mathrm{\Omega }_R+\omega i\mathrm{\Gamma })(i\omega +\mathrm{\Gamma }_Z)f^2\left(\mathrm{\Gamma }_Z\mathrm{\Gamma }\right)`$; $`\mathrm{\Omega }_R=\sqrt{\delta ^2+f^2}`$ is the Rabi frequency in the absence of the damping ($`\mathrm{\Gamma }_Z=\gamma =0`$). <sup>1</sup><sup>1</sup>1In above expressions we disregard the Fourier components of the damping terms (Eqs. (9), (10), (11)) since for the persistent signal $`g(t)`$ they are unimportant.
As can be concluded from these expressions the persistent low frequency oscillations of the spin components $`Z(t)`$, $`X(t)`$, $`Y(t)`$ appear as the response to the low frequency external force, $`g(t)`$, only in the presence of the high frequency excitation ($`f0`$). Exactly at resonance ($`\delta =0`$) $`\stackrel{~}{A}(\omega )=0`$, and $`\stackrel{~}{Z}(\omega )=0`$, but $`\stackrel{~}{B}(\omega )0`$. At this point the population of the two levels are equalized, and spin circularly rotates in the $`XY`$ plane with frequency $`\omega _0`$, with the center of the circle being precessed with the frequency $`\omega `$ of the low frequency external source.
As an example we show below the time evolution of the quantity $`\sigma _Z(t)=Z(t)`$, obtained from the numerical solution of the equations (3), (Low frequency Rabi spectroscopy for a dissipative two-level system), and (5) where we take a low frequency excitations as $`G(t)=Gcos(\omega _Lt)`$. The calculations have been performed with initial conditions $`\sigma _Z(0)=1`$, $`\sigma _X(0)`$=$`\sigma _Y(0)=0`$ for the following set of the parameters: $`F/h=36`$ MHz, $`\mathrm{\Delta }/h=1`$ GHz, $`\mathrm{\Gamma }/2\pi =4`$ MHz, $`\mathrm{\Gamma }_z/2\pi =1`$ MHz, $`ϵ/\mathrm{\Delta }=1`$, $`Z_0=1`$, $`\delta /2\pi =6.366`$ MHz, $`\omega _L/\mathrm{\Omega }_R=1`$. As is seen from Fig.1 in the absence of low frequency signal ($`G=0`$) the oscillations are damped out, while if $`G0`$ the oscillations persist.
The Fourier spectra of these signals are shown on Fig.2 for different amplitudes of low frequency excitation. For $`G=0`$ the Rabi frequency is positioned at approximately 26.2 MHz, which is close to $`\mathrm{\Omega }_R=26.24`$ MHz. With the increase of $`G`$ the peak becomes higher. It is worth noting the appearance of the peak at the second harmonic of Rabi frequency. This peak is due to the contribution of the terms on the order of $`G^2`$ which we omitted in our theoretical analysis.
The comparison of analytical and numerical resonance curves calculated for low frequency amplitude, $`G/h=1`$ MHz and different dephasing rates, $`\mathrm{\Gamma }`$ are shown on Fig.3. The curves at the figure are the peak-to-peak amplitudes of oscillations of $`Z(t)`$ calculated from Eq. (18) with $`\stackrel{~}{g}(\omega )=g(\delta (\omega +\omega _L)+\delta (\omega \omega _L))/2`$, where $`\delta (\omega )`$ is Dirac delta function. The point symbols are found from numerical solution of Eqs. (3),(Low frequency Rabi spectroscopy for a dissipative two-level system),(5). The widths of the curves depend on $`\mathrm{\Gamma }`$ (see the insert) and the positions of the resonances coincide with the Rabi frequency. A good agreement between numerics and Eq. 18, as shown at Fig. 3, is observed only for relative small low frequency amplitude $`G/h`$, for which our linear response theory is valid.
The key point of the method we described above is that it allows for the detection of the high frequency response of a TLS at a frequency which is much less than the gap frequency. The low frequency dynamics of the quantities $`\sigma _X,\sigma _Y`$, $`\sigma _Z`$ bears the information about the relaxation, $`\mathrm{\Gamma }_Z`$, and dephasing, $`\mathrm{\Gamma }`$, rates.
The experimental realization of this method (the detection scheme) depends on the problem under investigation. For example, the method can easily be adapted for NMR. In this case the quantities $`\sigma _Z,\sigma _X`$, $`\sigma _Y`$ are the longitudinal, $`M_Z`$, and the transversal polarizations of the sample, $`M_X`$, and $`M_Y`$. Indeed, from the comparison of the Hamiltonian (Low frequency Rabi spectroscopy for a dissipative two-level system) and the equations of motion (3), (Low frequency Rabi spectroscopy for a dissipative two-level system), (5) with those for nuclear spin ($`H=\stackrel{}{\mu }\stackrel{}{B}`$, $`d\stackrel{}{\mu }/dt=\gamma [\stackrel{}{\mu }\stackrel{}{B}]`$, where $`\gamma `$ is the gyromagnetic ratio), it is clear that the NMR case corresponds to the polarization of a sample with a field $`B_0`$ along the $`z`$ axes with a high frequency excitation $`B_1cos(\omega _0t)`$ and a low frequency probe $`G(t)`$ being applied in the ZX plane at an angle $`\theta `$ to the $`z`$ axis: $`B_Z=B_0+cos\theta (B_1cos(\omega _0t)+G(t))`$, $`B_X=sin\theta (B_1cos(\omega _0t)+G(t))`$, $`B_Y=0`$. This analogy allows for the direct application of Eqs. (18)-(22) to the nuclear spin with the following substitutions: $`\mathrm{\Delta }_\epsilon =\gamma \mathrm{}B_0`$, $`\epsilon /\mathrm{\Delta }_\epsilon =cos\theta `$, $`\mathrm{\Delta }/\mathrm{\Delta }_\epsilon =sin\theta `$, $`f=\gamma B_1sin\theta `$, $`g(\omega )=\gamma G(\omega )sin\theta `$, $`\epsilon /\mathrm{\Delta }=\mathrm{cot}\theta `$, $`Z_0=M_0`$ \- the equilibrium magnetization of a sample. In NMR all low frequency components of the magnetization, $`\stackrel{~}{Z}(\omega )`$ (18), $`\stackrel{~}{Y}(\omega )`$ (21), $`\stackrel{~}{X}(\omega )`$ (22), and their combinations are accessible for the measurements.
Our method can be directly applied to a persistent current qubit, which is a superconducting loop interrupted by three Josephson junctions , . For these qubits the successful experimental implementation of low frequency readout electronics has been demonstrated . The average current in the qubit loop is proportional to the low frequency part of the quantity $`\sigma _Z(t)`$, which is directly connected to the probabilities of occupation of the ground, $`P_{}(t)`$, and the excited, $`P_+(t)`$ states: $`\sigma _Z(t)`$=$`P_{}(t)`$-$`P_+(t)`$. Therefore, this current can be detected through the variation of its magnetic flux either by a DC SQUID or by a high quality resonant tank circuit inductively coupled to the qubit .
In conclusion, we proposed a method to study the Rabi oscillations in a dissipative TLS by irradiating it simultaneously with high, resonant, and low frequency. The low frequency response of a system has a clear resonance at the Rabi frequency with the resonance width being dependent on the damping rates. Therefore, the method allows for the experimental determination, in the low frequency domain, of the relaxation and dephasing rates of dissipative two-level systems.
###### Acknowledgements.
We are grateful to D. Averin, M. Grajcar, A. Korotkov, A. Shnirman, A. Smirnov, A. Zagoskin, A. Maassen van den Brink, W. Krech and V. Shnyrkov for fruitful discussions. The authors acknowledge the support from D-Wave Systems. Ya. G. acknowledges partial support by the INTAS grant 2001-0809. E. I. thanks the EU for support through the RSFQubit project. |
warning/0507/hep-ph0507177.html | ar5iv | text | # Possible Molecular Structure of the Newly Observed 𝑌(4260)
## 1 General discussion
The BaBar collaboration has recently announced that a very intriguing new state/structure $`Y(4260)\pi ^+\pi ^{}J/\psi `$ is observed in $`e^+e^{}ISR\pi ^+\pi ^{}J/\psi `$, where ISR stands for Initial-State Radiation . Their results indicate that $`Y(4260)`$ has spin-parity $`J^{PC}=1^{}`$. Its mass and width are
$`m=4.26\mathrm{GeV}/\mathrm{c}^2\text{and}\mathrm{\Gamma }90\mathrm{MeV}/\mathrm{c}^2.`$
An enhancement near 4.26 $`GeV/c^2`$ is clearly observed in $`e^+e^{}ISR\pi ^+\pi ^{}J/\psi `$ channel, but has not been observed in $`e^+e^{}hadrons`$, especially not in the $`D_{(s)}\overline{D}_{(s)}`$ channel. It may imply that the branching ratio of $`Y(4260)J/\psi \pi ^+\pi ^{}`$ is much larger than that of $`Y(4260)D\overline{D}`$.
Since its remarkable characteristics, this discovery stimulates intensive discussions about the structure of $`Y(4260)`$, especially if there is some new physics involved.
First, it exists in the energy-range of $`\psi `$ family, and one may expect that it involves both $`c`$ and $`\overline{c}`$ since in its strong-decay products there is no single charm(anti-charm). On other aspects, $`Y(4260)J/\psi \pi ^+\pi ^{}`$ is a three-body decay whereas $`Y(4260)D\overline{D}`$ is a two-body decay, and usually the former is about two orders smaller than the later due to a suppression from the phase space of final state. However, the data indicate a reversed pattern. This characteristic challenges our theory and demands a plausible interpretation.
The newly observed resonance is very unlikely to be accommodated in the regular $`c\overline{c}`$ structure even with higher radial and/or orbital excitations. It may be a clear signal for a new structure. In Ref. , the authors analyze the characteristics of $`Y(4260)`$ and proposes that $`Y(4260)`$ is perhaps a hybrid charmonium. Different from this explanation, Maiani et al. consider that the new resonance $`Y(4260)`$ may be the first orbital excitation of a diquark-antidiquark state $`[cs][\overline{c}\overline{s}]`$ . With a different point of view, instead of supposing $`Y(4260)`$ to be an exotic state, Llanes-Estrada proposes that the experimental evidence is not compelling to declare this state an exotic, and can be fitted within a standard quarkonium scenario.
In this work, we propose an alternative possibility that $`Y(4260)`$ is an s-wave molecular state of $`\rho \chi _{c1}(1P)`$, which is an isovector. In this framework, we can naturally explain why the branching ratio of $`Y(4260)J/\psi \pi ^+\pi ^{}`$ is larger than that of $`Y(4260)D\overline{D}`$. Meanwhile, we further predict existence of possible partner resonances of $`Y(4260)`$.
In next section we present our picture in detail and then we will draw our conclusion and make a discussion in the last section.
## 2 The molecular structure for $`Y(4260)`$.
Actually, there has been a long history about the molecular structure of hadrons. To explain some phenomena which are hard to find natural interpretations in the regular valence quark structure, people have tried to look for new structures beyond it. The molecular structure is one of the possible candidates. Okun and Voloshin studied the interaction between charmed mesons and proposed possibilities of the molecular states involving charmed quarks . Rujula, Geogi and Glashow suggested that $`\psi (4040)`$ is a $`D^{}\overline{D}^{}`$ molecular state . Moreover, the measured resonances $`f_0(980)`$, $`a_0(980)`$ may be reasonably interpreted as $`K\overline{K}`$ molecules . It seems that at the energy region of charm, molecular structure might be more favorable than at other energy regions.
Therefore, before invoking some fancy structures, let us study possibility to construct a molecular state for $`Y(4260)`$ and see if it coincides with the observed characteristics.
From the Data-book , we find that three particles $`\chi _{c0}`$, $`\chi _{c1}`$ and $`\chi _{c2}`$ in the $`c\overline{c}`$ meson spectrum may be candidates for the constituents in $`Y(4260)`$. The quantum numbers of $`\chi _{c0}`$, $`\chi _{c1}`$ and $`\chi _{c2}`$ are $`J^{PC}=0^{++},\mathrm{\hspace{0.33em}1}^{++}`$ and $`2^{++}`$ respectively. If combining them with $`\rho `$ meson to construct $`\chi _c\rho `$ systems, one can obtain states with spin-parity $`1^{}`$. Meanwhile, the masses of $`\chi _{c0}`$, $`\chi _{c1}`$ and $`\chi _{c2}`$ are well measured as $`3415.19\pm 0.34`$ MeV, $`3510.59\pm 0.10`$ MeV and $`3556.26\pm 0.11`$ MeV, thus we have that $`M_{\chi _c}+M_\rho `$ is 4185 MeV, 4280 MeV and 4326 MeV respectively for $`\chi _{c0}`$, $`\chi _{c1}`$ and $`\chi _{c2}`$.
For an s-wave molecular state, one should expect that the sum of the constituent masses is closer to the mass of the resonance. The difference is due to the interaction between the constituents which in general results in a negative binding energy for s-wave. For the three-combinations $`\rho \chi _{c0}`$, $`\rho \chi _{c1}`$ and $`\rho \chi _{c2}`$, one can observe that the mass sum of $`\rho `$ and $`\chi _{c1}`$ is mostly close to the mass of $`Y(4260)`$. Based on the above considerations, we propose that $`Y(4260)`$ may be a molecular state of $`\rho `$ and $`\chi _{c1}`$. Namely the mass sum of $`\rho `$ and $`\chi _{c1}`$ is about 20 MeV above 4260 MeV and the difference is paid to the negative binding energy.
The decay pattern of $`Y(4260)`$ is the most important issue to concern, because it may provide us the information about the structure of $`Y(4260)`$.
In the Fig.1, we present the quark diagrams for $`Y(4260)J/\psi \pi ^+\pi ^{}`$ and $`Y(4260)D\overline{D}`$.
For Fig.1 (a) and (b), the transition matrix elements can be expressed as
$`(Y(4260)J/\psi \pi ^+\pi ^{})`$ $`=`$ $`\rho ^0,J/\psi |_{dis}|Y(4260)\times \pi ^+\pi ^{}||\rho ^0,`$ (1)
$`(Y(4260)D\overline{D})`$ $`=`$ $`D\overline{D}|_{cross}|Y(4260)`$ (2)
where $`_{dis}`$ corresponds to the hamiltonian which breaks the bound state $`Y(4260)`$ into free $`J/\psi `$ and $`\rho ^0`$ via exchanging $`\sigma `$ meson (maybe, exchanges of multi-soft-gluons and even glueball of $`0^{++}`$ can also contribute, but definitely $`\sigma `$exchange plays the leading role), $``$ is a strong interaction which causes $`\rho ^0`$ decay into $`\pi ^+\pi ^{}`$. $`_{cross}`$ is an interaction, by which quarks (antiquarks) in $`\chi _c`$ and $`\rho ^0`$ exchange and turn into hadronic $`D`$ and $`\overline{D}`$, in the process quark lines cross with each other (see Fig.1 (b)).
$`\chi _{c1}`$ is a $`1^{++}`$ axial vector, $`J/\psi `$ is a $`1^{}`$ vector and both of them are isosinglet, the couplings of $`\chi _{c1}\pi (\rho )\chi _{c1}`$ and $`\chi _{c1}\pi (\rho )J/\psi `$ are forbidden by the isospin conservation, and only $`\sigma `$ of $`0^{++}`$ can be exchanged and is the main contribution to the potential which holds the constituents in a molecule. The interactions of $`\chi _{c1}\sigma \chi _{c1}`$ and $`\chi _{c1}\sigma J/\psi `$ are obviously OZI suppressed , so cannot be very large. One may write down the effective lagrangians
$$L_1=g_1A_{1\mu }A_1^\mu \sigma ,\text{ for}\chi _{c1}\sigma \chi _{c1},$$
and
$$L_2=g_2\stackrel{~}{F}_{1\mu \nu }F_2^{\mu \nu }\sigma ,\text{ for}\chi _{c1}\sigma J/\psi ,$$
where
$$\stackrel{~}{F}_{1\mu \nu }\frac{1}{2}ϵ_{\mu \nu \alpha \beta }F_1^{\alpha \beta },$$
and $`A_{1\mu },A_{2\mu }`$ correspond to axial vector $`\chi _{c1}`$ and vector $`J/\psi `$ respectively.
It is supposed that the $`\sigma `$ exchange provides an attractive potential $`\frac{e^{m_\sigma r}}{r}`$ between $`\chi _{c1}`$ and $`\rho `$ to construct a bound state. Apparently, the coupling is OZI suppressed, and the binding is relatively loose.
More concretely, in Fig.1 (a), $`\chi _c`$ may convert into $`J/\psi `$ mainly via exchanging $`\sigma `$ particle with the constituent $`\rho `$ meson. It is noted that $`L_2`$, which turns $`\chi _{c1}`$ into $`J/\psi `$ is a p-wave interaction and proportional to the linear momentum to guarantee the parity match. The differentiation may result in an opposite sign to the potential between $`\chi _{c1}`$ and $`\rho ^0`$ and provide an effective repulsion. Then the bound state dissolves into free $`J/\psi `$ and $`\rho `$, and then a strong decay of $`\rho ^0\pi ^+\pi ^{}`$ follows. Here, for a general discussion, we ignore all the dynamical details and make only an estimate on the order of magnitude. Since the branching ratio of $`\rho ^0\pi ^+\pi ^{}`$ is almost 100%, we can suppose that the transition of the constituent of $`Y(4260)`$, i.e. $`\rho ^0`$ to $`\pi ^+\pi ^{}`$ is overwhelming.
The total width is then,
$`\mathrm{\Gamma }(Y(4260)J/\psi \pi ^+\pi ^{})`$ $`=`$ $`{\displaystyle \frac{1}{2M}}{\displaystyle }{\displaystyle \frac{d^3p_{_{J/\psi }}}{(2\pi )^3}}{\displaystyle \frac{1}{2E_{J/\psi }}}{\displaystyle \frac{d^3p__\rho }{(2\pi )^3}}{\displaystyle \frac{1}{2E_\rho }}(2\pi )^4\delta ^4(MP_{J/\psi }P_\rho )`$ (3)
$`|(Y(4260)J/\psi +\rho ^0)|^2\times BR(\rho ^0\pi ^+\pi ^{}),`$
where $`M`$ is the mass of $`Y(4260)`$ and $`P_{J/\psi },p_{_{J/\psi }},P_\rho ,p__\rho `$ are the four- and three-momenta of $`J/\psi `$ and $`\rho `$ respectively.
Comparing with Fig.1. (a), Fig.1 (b) involves an extra color re-combination process which leads to a suppression, this suppression factor is
$`{\displaystyle \frac{|(Y(4260)D\overline{D})|}{|(Y(4260)J/\psi \pi ^+\pi ^{})|}}\alpha ={\displaystyle \frac{1}{3}}.`$ (4)
There may be a numerical factor $`g`$ coming from dynamics and it is completely a non-perturbative QCD factor. For a rough estimate it can be approximated as unity.
$`Y(4260)J/\psi \pi ^+\pi ^{}`$ seems to be a three-body decay, thus there could be a suppression from the phase space of final states. However, in our picture of molecular state, it is not a real three-body decay, instead, it is a two-step process, namely first $`Y(4260)`$ dissolves into $`J/\psi `$ and $`\rho ^0`$ and then $`\rho ^0`$ transits into $`\pi ^+\pi ^{}`$. Since the total width is proportional to a two-body decay rate multiplied by the branching ratio of $`\rho ^0\pi ^+\pi ^{}`$ which is 100% almost, there does not exist the phase space suppression factor at all.
Due to the color re-matching factor, one can expect that the decay rate of $`Y(4260)J/\psi \pi ^+\pi ^{}`$ is about one order larger than that of $`Y(4260)D\overline{D}`$. The concrete dynamics may change this ratio more or less, but here we just take this value from estimate of order of magnitude. This value qualitatively coincides with the experimental results.
## 3 More discussions and conclusion
We suggest that the observed $`Y(4260)`$ is an s-wave molecular state of $`\chi _{c1}`$ and $`\rho ^0`$. It is natural to consider another two partner molecular states, namely $`\chi _{c0}+\rho ^0`$ and $`\chi _{c2}+\rho ^0`$ in s-wave. Their spin-parity can be different, but which one is dominant depends on the concrete dynamics. For the simplest case, supposing they are also $`1^{}`$, we may expect that the molecular state of $`\chi _{c2}+\rho ^0`$ is only 40 MeV above 4260 MeV (supposing it has the same binding energy as that for $`\chi _{c1}+\rho ^0`$), on other side, the total width of $`Y(4260)`$ is 90 MeV, thus this molecular state might be hidden in the observed peak of $`Y(4260)`$, in other words, the experimentally observed peak $`Y(4260)`$ may cover two close states. Meanwhile, the molecular state of $`\chi _{c0}+\rho ^0`$ could be 100 MeV below the central value of the peak and thus corresponds to a new state which can be used as a test of the model. Namely, if this partner resonance is observed in the future experiments, one can claim that the molecular structure postulation may be correct, otherwise, we need to consider other possible mechanisms to suppress its production rate from dynamics or abandon the molecular state interpretation. Moreover, the molecule of $`\chi _{c1}\rho ^0`$ is a component of an isovector, so there may exist another two components of the isotriplet, i.e. $`\chi _{c1}\rho ^\pm `$ which may decay into $`J/\psi \pi ^\pm \pi ^0`$ with comparable rates of $`Y(4260)J/\psi \pi ^+\pi ^{}`$. They may be experimentally observable. For the molecular structure, $`\sigma `$ exchange between $`\chi _{c1}`$ and $`\rho ^0`$ may result in an attractive potential which binds them into a molecule. Since the coupling is OZI suppressed, the binding is relatively loose.
In our scenario, the favorable decay mode of $`Y(4260)`$ is $`Y(4260)J/\psi \pi ^+\pi ^{}`$. The molecular structure of $`Y(4260)`$ results in different decay pattern from $`\psi (3770)`$ which is supposed to be a pure $`c\overline{c}`$ charmonium. Namely if we take the $`D\overline{D}`$ mode as a standard, the rate of $`Y(4260)J/\psi \pi ^+\pi ^{}`$ is larger than that of $`Y(4260)D\overline{D}`$ by an order.
By contraries, in the hybrid charmonium structure , where a color-octet $`c\overline{c}`$ system is bound with an octet valence gluon, since gluon is flavor-blind, it has the same coupling to $`q\overline{q}(q=u,d)`$ and $`s\overline{s}`$, thus besides a small suppression from the phase space, $`Y(4260)J/\psi \pi ^+\pi ^{}`$ and $`Y(4260)J/\psi K\overline{K}`$ should be comparable unless there exist certain mechanisms to suppress $`K\overline{K}`$ production. In the diquark-anti-diquark picture of $`[cs][\overline{c}\overline{s}]`$, the mode $`Y(4260)J/\psi K\overline{K}`$ overwhelms $`Y(4260)J/\psi \pi ^+\pi ^{}`$. In our picture of molecular state, the mode of $`Y(4260)J/\psi K\overline{K}`$ can only be realized via final state interaction $`\pi ^+\pi ^{}K\overline{K}`$, so that the rate of $`Y(4260)J/\psi K\overline{K}`$ is much smaller than that of $`Y(4260)J/\psi \pi ^+\pi ^{}`$.
Let us turn to a subtle and difficult subject, the production of $`Y(4260)`$ in $`e^+e^{}`$ collisions. The production may occur via the so-called hairpin mechanism which does not suffer from the suppression due to color matching. It seems that it has a larger production rate than the direct (non-resonant) production of $`D\overline{D}`$ at first glimpse. However, a detailed analysis indicates that unless the energy $`\sqrt{s}`$ of $`e^+e^{}`$ collisions can be precisely tuned to 4260 MeV, the energy conservation demands production of other hadrons such as pions in company with $`Y(4260)`$, and the constraint from the final product phase space would greatly suppress its production rate. To achieve concrete values one must carry out model-dependent calculations and it is beyond the scope of this work.
One more observation is that $`\rho ^0`$ only decays into $`\pi ^+\pi ^{}`$, but not $`\pi ^0\pi ^0`$, therefore, if the molecular picture is right, the mode of $`Y(4260)J/\psi \pi ^0\pi ^0`$ must be very suppressed. Moreover, since in our picture $`\pi ^+\pi ^{}`$ are produced from the real $`\rho ^0`$meson, the measured invariant-mass spectrum of $`\pi ^+\pi ^{}`$ should peak up at $`m_\rho `$. Looking at the figure (Fig. 3 of ), the dipion mass distribution of $`Y(4260)J/\psi \pi ^+\pi ^{}`$ seems to show some fine structures. Since the spectrum is due to $`\rho ^0\pi ^+\pi ^{}`$ decay, it is a p-wave structure. The authors of ref. indicate that the observed $`\pi \pi `$ spectrum is somehow rather an s-wave comparing with the Mote-Carlo results, thus the molecular interpretation offers a non-standard interpretation for the bump. More precise experiments in the future may give a decisive conclusion. Thus the possibility that $`Y(4260)`$ is a molecular state, is indeed worth careful studies.
Now let us draw a brief conclusion. We propose that the newly observed $`Y(4260)`$ is a molecular state of $`\chi _{c1}`$ and $`\rho ^0`$, and our analysis indicates that this picture qualitatively coincides with the experimental data. We naturally explain why the rate of $`Y(4260)\pi ^+\pi ^{}J/\psi `$ is larger than that of $`Y(4260)D\overline{D}`$, namely why $`Y(4260)`$ is only observed in $`e^+e^{}ISR\pi ^+\pi ^{}J/\psi `$, but not in $`Y(4260)D\overline{D}`$. We have also made predictions on existence of two other components of the isotriplet, $`\chi _{c1}\rho ^\pm `$ which may be observed in channels $`J/\psi \pi ^\pm \pi ^0`$, and the extra partner resonances $`\chi _{c0(c2)}\rho ^0`$ along with their isotriplet components. It is suggested that the state of $`\chi _{c0}+\rho ^0`$ may be distinguished from $`Y(4260)`$ and can be experimentally measured, so should serve as a test of the model. The future experiments will collect more data and confirm or negate the various theoretical models as well as ours. For such experiments besides the B-factories, BES and CLEO are also ideal places.
In our model, we only discuss the qualitative characteristics and make estimate of order of magnitude, but ignore all the dynamics. Definitely all the details of dynamics may change the numbers quite much, but we hope that the qualitative conclusion and analysis would remain unchanged, because they are independent of the dynamical details.
Acknowledgment:
This work is supported by the National Natural Science Foundation of China. We are grateful to Prof. K.T. Chao, Dr. S.W. Ye and S.L. Zhu for helpful discussions. |
warning/0507/gr-qc0507038.html | ar5iv | text | # Loop Quantum Geometry: A primer
## I Introduction
Loop Quantum Gravity (LQG) has become in the past years a mayor player as a candidate for a quantum theory of gravity. On the one hand it has matured into a serious contender together with other approaches such as String/M Theory, but on the other is it not as well understood, neither properly credited as a real physical theory of quantum gravity. The purpose of this contribution is to provide a starting point for those interested in learning the basics of the theory and to provide at the same time an introduction to the rich literature on the subject which includes very well written reviews and monographs. Given the space constraint we shall not attempt to write a comprehensive review of LQG, but to provide, we hope, and useful guide to the subject.
Let us start by providing a list of references that will be useful in the various stages. Firstly, there are several primer introductions to the subject, written for different purposes. For instance, there was for many years the canonical primer by Pullin pullin . Unfortunately, it is now somewhat dated. Good introductions to spin networks and recoupling theory needed in LQG are given by the primers by Rovelli Rovelli:1998gg and Major seth . There are recent up-to-date accounts written for non-experts that give nice motivation, historical perspective and an account of recent and in progress work from two different perspectives AA:NJP and Smolin:2004sx . There are also technical reviews that give many details and are certainly a good read AL:review , Perez:2004hj , Thiemann:2002nj , and (from an outside perspective) Nicolai:2005mc .
The subject has matured enough so that several monographs have been written, including some recent and updated. These monographs approach and present the subject from different perspectives depending, of course, on the authors own taste. From these, it is worth mentioning two. The first one by Rovelli is physically motivated but not so heavy in its mathematical treatment, and can be found in Rovelli:2004tv . A mathematically precise treatment, but not for the faint of heart is given by the monograph by Thiemann Thiemann:2001yy . There have been also several nice reviews that motivate and give a birdseye view of the subject such as Rovelli:1999hz , Rovelli:1997yv and pullin2 . Finally, there are several accounts on comparisons between loop quantum gravity and other approaches, such as string theory. On chronological order, we have a review by Rovelli Rovelli:1997qj , an entertaining dialog Rovelli:2003wd and a critical assessment by Smolin Smolin:2003rk .
The second purpose of this paper is to present an introduction to the formalism known as (loop) quantum geometry. The difference between quantum geometry and loop quantum gravity is that the former is to be thought of as the (new) formalism dealing with background-free quantum theories based on connections, whereas the later is a particular implementation where gravity, as defined by general relativity, is the theory under consideration. For instance, one could think of applying the same formalism to more general theories such as supergravity and/or higher derivative theories. In the remainder of this section we shall give a motivation for why one should study loop quantum gravity, when one is interested in the basic problem of uniting quantum mechanics and the theory of gravitation.
Why should one study loop quantum gravity? For one thing, it is based on two basic principles, namely the general principles of quantum theory and one of the main lessons from general relativity: that physics is diffeomorphism invariant. This means that the field describing the gravitational interaction, and the geometry of spacetime is fully dynamical and interacting with the rest of the fields present. When one is to consider its quantum description, this better be background independent. The fact that LQG is based in general principles of quantum mechanics means only that one is looking for a description based on the standard language of quantum mechanics: states are elements on a Hilbert space (well defined, of course), observables will be Hermitian operators thereon, etc. This does not mean that one should use all that is already known about quantizing fields. Quite on the contrary, the tools needed to construct a background independent quantization (certainly not like the quantization we know), are rather new.
Another reason for studying LQG is that this is the most serious attempt to perform a full non-perturbative quantization of the gravitational field. It is an attempt to answer the following question: can we quantize the gravitational degrees of freedom without considering matter on the first place? Since LQG aims at being a physical theory, which means it better be falsifiable, one expects to answer that question unambiguously, whenever one has the theory fully developed. This is one of the main present challenges of the theory, namely to produce predictions that can be tested experimentally. Since the theory does not suffer from extra dimensions nor extended symmetries, one expects that the task will be feasible to complete, without the burden of getting rid on those extra features. Will this be the final theory describing the quantum degrees of freedom of the gravitational field? Only experiments will tell, but for the time being, LQG remains an intriguing possibility very well worth the trial.
This contribution is organized as follows: In Sec. II we provide some of material needed in order to be ready to fully grasp the details of LQG. In particular, we recall the canonical description go general relativity in the geometrodynamics language and perform the change of variables to go to the description in terms of connections. In Sec. III we consider within the classical description of the gravitational field, the observables that will be regarded as the basic objects for the quantization. We shall see that both background independence and diffeomorphism invariance lead us to select holonomies and electric fluxes as the basic objects. Section IV will be devoted to the discussion of the Hilbert space of the theory, its multiply characterizations and a particular basis that is very convenient, namely, the spin network basis. In Section V we provide a discussion of the new (quantum) geometry that the theory presents for us, focusing on basic geometrical operators. Section VI will be devoted to a list of accomplishments and open issues.
Finally, we should note that our list of references is minimalist, trying to concentrate on review articles or monographs rather than in the original articles. A more complete reference list can be found in the books Thiemann:2001yy and Rovelli:2004tv and in the bibliography compilation biblio . An online guide for learning LQG with references can be found in seth2 .
## II Preliminaries
What are the pre-requisites for learning LQG? First of all, a reasonable knowledge of General Relativity (GR), specially its Hamiltonian formulation as in wald1 and poisson , acquaintance with QFT on flat space-time and preferably some notions of “quantization”, namely the passage from a classical theory to a quantum one as in wald2 . Finally, the language of gauge field theories is essential, including connections and holonomies. A good introduction to the subject, including the necessary geometry, is given in baez .
The first step is to introduce the basic classical variables of the theory. Since the theory is described by a Hamiltonian formalism, this means that the 4-dim spacetime $`M`$ is of the form $`M=\mathrm{\Sigma }\times `$, where $`\mathrm{\Sigma }`$ is a 3-dimensional manifold. The first thing to do is to start with the geometrodynamical phase space $`\mathrm{\Gamma }_\mathrm{g}`$ of Riemannian metrics $`q_{ab}`$ an $`\mathrm{\Sigma }`$ and their canonical momenta $`\stackrel{~}{\pi }^{ab}`$ (related to the extrinsic curvature $`K_{ab}`$ of $`\mathrm{\Sigma }`$ into $`M`$ by $`\stackrel{~}{\pi }^{ab}=\sqrt{q}(K^{ab}\frac{1}{2}q^{ab}K)`$, with $`q=\mathrm{det}(q_{ab})`$ and $`K=q^{ab}K_{ab}`$). Recall that they satisfy,
$$\{\stackrel{~}{\pi }^{ab}(x),q_{cd}(y)\}=2\kappa \delta _{(c}^a\delta _{d)}^b\delta ^3(x,y);\{q_{ab}(x),q_{cd}(y)\}=\{\stackrel{~}{\pi }^{ab}(x),\stackrel{~}{\pi }^{cd}(y)\}=0$$
(1)
General Relativity in these geometrodynamical variables is a theory with constraints, which means that the canonical variables $`(q_{ab},\stackrel{~}{\pi }^{ab})`$ do not take arbitrary values but must satisfy four constraints:
$$^b=D_a(\stackrel{~}{\pi }^{ab})=0\mathrm{and},=\sqrt{q}[R^{(3)}+q^1(\frac{1}{2}\stackrel{~}{\pi }^2\stackrel{~}{\pi }^{ab}\stackrel{~}{\pi }_{ab})]0$$
(2)
The first set of constraints are known as the vector constraint and what they generate (its gauge orbit) are spatial diffeomorphisms on $`\mathrm{\Sigma }`$. The other constraint, the scalar constraint (or super-Hamiltonian) generates “time reparametrizations”. We start with 12 degrees of freedom, minus 4 constraints means that the constraint surface has 8 dimensions (per point) minus the four gauge orbits generated by the constraints giving the four phase space degrees of freedom, which corresponds to the two polarizations of the gravitational field.
In order to arrive at the connection formulation, we need first to enlarge the phase space $`\mathrm{\Gamma }_\mathrm{g}`$ by considering not metrics $`q_{ab}`$ but the co-triads $`e_a^i`$ that define the metric by,
$$q_{ab}=e_a^ie_b^j\delta _{ij}$$
(3)
where $`i,j=1,2,3`$ are internal labels for the frames. These represent 9 variables instead of the 6 defining the metric $`q_{ab}`$, so we have introduced more variables, but at the same time a new symmetry in the theory, namely the $`SO(3)`$ rotations in the triads. Recall that a triad $`e_a^i`$ and a rotated triad $`e_a^i(x)=U_{}^{i}{}_{j}{}^{}(x)e_a^j(x)`$ define the same metric $`q_{ab}(x)`$, with $`U_{}^{i}{}_{j}{}^{}(x)SO(3)`$ a local rotation. In order to account for the extra symmetry, there will be extra constraints (first class) that will get rid of the extra degrees of freedom introduced. Let us now introduce the densitized triad as follows:
$$\stackrel{~}{E}_i^a=\frac{1}{2}ϵ_{ijk}\stackrel{~}{\eta }^{abc}e_b^je_c^k$$
(4)
where $`\stackrel{~}{\eta }^{abc}`$ is the naturally defined levi-civita density one antisymmetric object. Note that $`\stackrel{~}{E}_i^a\stackrel{~}{E}_j^b\delta ^{ij}=qq^{ab}`$.
Let us now consider the canonical variables. It turns out that the canonical momenta to the densitized triad $`\stackrel{~}{E}_i^a`$ is closely related to the extrinsic curvature of the metric,
$$K_a^i=\frac{1}{\sqrt{\mathrm{det}(\stackrel{~}{E})}}\delta ^{ij}\stackrel{~}{E}_j^bK_{ab}$$
(5)
For details see Perez:2004hj . Once one has enlarged the phase space from the pairs $`(q_{ab},\stackrel{~}{\pi }^{ab})`$ to $`(\stackrel{~}{E}_i^a,K_b^j)`$, the next step is to perform the canonical transformation to go to the $`AshtekarBarbero`$ variables. First we need to introduce the so called spin connection $`\mathrm{\Gamma }_a^i`$, the one defined by the derivative operator that annihilates the triad $`e_a^i`$ (in complete analogy to the Christoffel symbol that defined the covariant derivative $`D_a`$ killing the metric). It can be inverted from the form,
$$_{[a}e_{b]}^i+ϵ_{}^{i}{}_{jk}{}^{}\mathrm{\Gamma }_a^je_b^k=0$$
(6)
This can be seen as an extension of the covariant derivative to objects with mixed indices. The key to the definition of the new variables is to combine these two objects, namely the spin connection $`\mathrm{\Gamma }`$ with the object $`K_a^i`$ (a tensorial object), to produce a new connection
$${}_{}{}^{\gamma }A_{a}^{i}:=\mathrm{\Gamma }_a^i+\gamma K_a^i$$
(7)
This is the Ashtekar-Barbero Connection. Similarly, the other conjugate variable will be the rescaled triad,
$${}_{}{}^{\gamma }\stackrel{~}{E}_{i}^{a}=\stackrel{~}{E}_i^a/\gamma $$
(8)
Now, the pair $`({}_{}{}^{\gamma }A_{a}^{i},{}_{}{}^{\gamma }\stackrel{~}{E}_{i}^{a})`$ will coordinatize the new phase space $`\mathrm{\Gamma }_\gamma `$. We have emphasized the parameter $`\gamma `$ since this labels a one parameter family of different classically equivalent theories, one for each value of $`\gamma `$. The real and positive parameter $`\gamma `$ is known as the Barbero-Immirzi parameter barbero ; Immirzi . In terms of these new variables, the canonical Poisson brackets are given by,
$$\{{}_{}{}^{\gamma }A_{a}^{i}(x),{}_{}{}^{\gamma }\stackrel{~}{E}_{j}^{b}(y)\}=\kappa \delta _a^b\delta _j^i\delta ^3(x,y).$$
(9)
and,
$$\{{}_{}{}^{\gamma }A_{a}^{i}(x),{}_{}{}^{\gamma }A_{b}^{j}(y)\}=\{{}_{}{}^{\gamma }\stackrel{~}{E}_{i}^{a}(x),{}_{}{}^{\gamma }\stackrel{~}{E}_{j}^{b}(y)\}=0$$
(10)
Let us summarize. i) We started with the geometrodynamical phase space where the configuration space was taken to be the space of 3-Riemannian metrics, ii) enlarged it by considering triads instead of metrics, and iii) performed a canonical transformation, and changed the role of the configuration variable; the connection $`{}_{}{}^{\gamma }A_{a}^{i}`$ (that has the information about the extrinsic curvature) is now the configuration variable and the (densitized) triad is regarded as the canonical momenta. The connection $`{}_{}{}^{\gamma }A_{a}^{i}`$ is a true connection since it was constructed by adding to the spin connection $`\mathrm{\Gamma }_a^i`$ a tensor field, that yields a new connection. We started by considering connection taking values in $`so(3)`$, the Lie algebra of $`SO(3)`$, but since as Lie algebras it is equivalent to $`su(2)`$, we will take as the gauge group the simply connected group $`SU(2)`$ (which will allow to couple fermions when needed). The phase space $`\mathrm{\Gamma }_\gamma `$ is nothing but the phase space of a $`SU(2)`$ Yang-Mills theory.
Let us now write down the constraints in terms of the new phase space variables. The new constraint that arises because of the introduction of new degrees of freedom takes a very simple form,
$$G_i=𝒟_a\stackrel{~}{E}_i^a0$$
(11)
that is, it has the structure of Gauss’ law in Yang-Mills theory and that is the name that has been adopted for it. We have denoted by $`𝒟`$ the covariant defined by the connection $`{}_{}{}^{\gamma }A_{a}^{i}`$, such that $`𝒟_a\stackrel{~}{E}_i^a=_a\stackrel{~}{E}_i^a+ϵ_{ij}^{}{}_{}{}^{k}{}_{}{}^{\gamma }A_{a}^{j}\stackrel{~}{E}_k^a`$. The vector and scalar constraints now take the form,
$$V_a=F_{ab}^i\stackrel{~}{E}_i^b(1+\gamma ^2)K_a^iG_i0$$
(12)
where $`F_{ab}^i=_a{}_{}{}^{\gamma }A_{b}^{i}_b{}_{}{}^{\gamma }A_{a}^{i}+ϵ_{}^{i}{}_{jk}{}^{}{}_{}{}^{\gamma }A_{a}^{j}{}_{}{}^{\gamma }A_{b}^{k}`$ is the curvature of the connection $`{}_{}{}^{\gamma }A_{b}^{j}`$. The other constraint is,
$$𝒮=\frac{\stackrel{~}{E}_i^a\stackrel{~}{E}_j^b}{\sqrt{\mathrm{det}(\stackrel{~}{E})}}\left[ϵ_{}^{ij}{}_{k}{}^{}F_{ab}^k2(1+\gamma ^2)K_{[a}^iK_{b]}^j\right]0$$
(13)
Note incidently that if $`(1+\gamma ^2)=0`$, the constraints would simplify considerably. That was the original choice of Ashtekar, that rendered the connection complex (and thus the corresponding gauge group non-compact), and it had a nice geometrical interpretation in terms of self-dual fields. Historically the emphasis from complex to real connections (and a compact group as a consequence) was due to the fact that in this case the mathematical well defined construction of the Hilbert space has been completed, whereas the non-compact case remains open. For more details see AL:review and specially Sec 3.2 of Perez:2004hj .
The next step is to consider the right choice of variables, now seen as functions of the phase space $`\mathrm{\Gamma }_\gamma `$ that are preferred for the non-perturbative quantization we are seeking. As we shall see, the guiding principle will be that the functions (defined by an appropriate choice of smearing functions) will be those that can be defined without the need of a background structure, i.e. a metric on $`\mathrm{\Sigma }`$.
## III Holonomies and Fluxes
We have seen that the classical setting for the formulation of the theory is the phase space of a Yang-Mill theory, with extra constraints. Since the theory possesses these constraints, the strategy to be followed is to quantize first and then to impose the set of constraints as operators on a Hilbert space. This Hilbert space is very important since it will be the home where the imposition of the constraints will be implemented and its structure will have some physical relevance. This is known as the Kinematical Hilbert Space $`_{\mathrm{kin}}`$. One of the main achievements of LQG is that this space has been rigourously defined, something that was never done in the old geometrodynamics program.
The main objective of this section is to motivate and construct the classical algebra of observables that will be the building blocks for the construction of the space $`_{\mathrm{kin}}`$. In QFT one could say that there are two important choices when quantizing a classical system. The first one is the choice of the algebra of observables, consisting of two parts, the choice of variables, and of functions thereof. The second choice is a representation of this chosen classical algebra into a Hilbert space. Both steps normally involve ambiguities. This is also true of the phase space we are starting with. Remarkably for our case, the physical requirements of background independence and diffeomorphism invariance will yield a choice of variables and of a representation that is in a sense, unique. The infinite freedom one is accustomed to in ordinary QFT is here severely reduced. Background independence and diffeomorphism invariance impose very strong conditions on the possible quantum theories.
Let us start by considering the connection $`A_a^i`$ (from now on we shall omit the $`\gamma `$ label). The most natural object one can construct from a connection is a holonomy $`h_\alpha (A)`$ along a loop $`\alpha `$. This is an element of the gauge group $`G=SU(2)`$ and is denoted by,
$$h_\alpha (A)=𝒫\mathrm{exp}\left(_\alpha A_ads^a\right)$$
(14)
The path-order exponential of the connection. Note that for notational simplicity we have omitted the ‘lie-algebra indices’. The connection as an element of the lie algebra, in the fundamental representation, should be written as $`A_a^i(\tau _i)_B^A`$, with $`A,B=1,2`$ the $`2\times 2`$-matrix indices of the Pauli matrices $`\tau _i`$. From the holonomy, it is immediate to construct a gauge invariant function by taking the trace arriving then at the Wilson loop $`T[\alpha ]:=\frac{1}{2}\mathrm{Tr}𝒫\mathrm{exp}(_\alpha A_ads^a)`$.
Several remarks are in order. i) When seen as a smearing of the connection, it is clear that the holonomy represents a one dimensional smearing (as compared with, say, the three dimensional smearings $`A_a^i(x)g_i^a(x)\mathrm{d}^3x`$ used in ordinary QFT); ii) The loop $`\alpha `$ can be seen as a label, but the holonomy is a function of the connection $`A_a^i`$; iii) Even when the holonomies are functions of the connection $`A`$, they only “prove” the connection along the loop $`\alpha `$. In order to have a useful set of functions that can separate points of the space of smooth connections $`𝒜`$, one needs to consider for instance Wilson loop functions along all possible loops on $`\mathrm{\Sigma }`$. The algebra generate by such functions is called the holonomy algebra $`𝒜`$.
In order to implement the idea that one should look for generalized notions that will replace the loops, let us consider the most obvious extension. In recent years the emphasis has shifted from loops to consider instead closed graphs $`\mathrm{{\rm Y}}`$, that consist of $`N`$ edges $`e_I`$ ($`I=1,2,\mathrm{},N`$), and $`M`$ vertices $`v_\mu `$, with the restriction that there are no edges with ‘loose ends’. As a aside remark one should mention that every graph $`\mathrm{{\rm Y}}`$ can be decomposed in independent loops $`\alpha _i`$ based at a chosen vertex $`v`$. Given a graph $`\mathrm{{\rm Y}}`$, one can consider the parallel transport along the edges $`e_I`$, the end result is an element of the gauge group $`g_I=h(e_I)G`$ for each such edge. One can then think of the connection $`A_a^i`$ as a map from graphs to $`N`$-copies of the gauge group: $`A_a^i:\mathrm{{\rm Y}}G^N`$. Furthermore, one can think of $`𝒜_\mathrm{{\rm Y}}`$ as the configuration space for the graph $`\mathrm{{\rm Y}}`$, that is homeomorphic to $`G^N`$ (one can regard $`𝒜_\mathrm{{\rm Y}}`$ as the configuration space of the ‘floating lattice’ gauge theory over $`\mathrm{{\rm Y}}`$).
Once we have recognized that one can associate a configuration space for all graphs, one should not loose perspective that the relevant classical configuration space is still the space $`𝒜`$ of all (smooth) connections $`A_a^i`$. What we are doing at the moment is to construct relevant configuration functions, making use of the graphs and the space $`𝒜_\mathrm{{\rm Y}}`$. In particular, what we need is to consider generalizations of the Wilson loops $`T[\alpha ]`$ defined previously. As we have mentioned before, every graph $`\mathrm{{\rm Y}}`$ can be decomposed into independent loops $`\alpha _i`$ and the corresponding Wilson loops $`T[\alpha _i]`$ are a particular example of functions defined over $`𝒜_\mathrm{{\rm Y}}`$. What we shall consider as a generalization of the Wilson loop are all possible functions defined over $`𝒜_\mathrm{{\rm Y}}`$ (in a sense the Wilson loops generate these functions, but are overcomplete). Thus, a function $`c:G^N`$ defines a cylindrical function $`C_\mathrm{{\rm Y}}`$ of the connection $`A`$ as,
$$C_\mathrm{{\rm Y}}:=c(h(e_1),h(e_2),\mathrm{},h(e_N))$$
(15)
By considering all possible functions $`c`$ and all possible embedded graphs $`\mathrm{{\rm Y}}`$, we generate the algebra of functions known as Cyl (it is closely related to the holonomy algebra, and it can be converted into a $`C^{}`$-algebra $`\overline{\mathrm{Cyl}}`$, by suitable completion).
Even when we shall consider in what follows the full algebra Cyl, one should keep in mind that the basic objects that build it are precisely the holonomies, functions of the connection smeared along one dimensional objects. Let us now discuss why this choice of configuration functions is compatible with the basic guiding principles for the quantization we are building up, namely diffeomorphism invariance and background independence. Background independence is clear since there is no need for a background metric to define the holonomies. Diffeomorphism invariance is a bit more subtle. Clearly, when one applies a diffeomorphism $`\varphi :\mathrm{\Sigma }\mathrm{\Sigma }`$, the holonomies transform in a covariant way
$$\varphi _{}h(e_I)=h(\varphi ^1e_I),$$
(16)
that is, the diffeomorphism acts by moving the edge (or loop). How can we then end up with a diffeo-invariant quantum theory? The strategy in LQG is to look for a diffeomorphism invariant representation of the diffeo-covariant configuration functions. As we shall see later, this has indeed been possible and in a sense represents the present ‘success’ of the approach.
Let us now consider the functions depending of the momenta that will be fundamental in the (loop) quantization. The basic idea is again to look for functions that are defined in a background independent way, that are natural from the view point of the geometric character of the object (1-form, 2-form, etc), and that transform covariantly with respect to the gauge invariances of the theory. Just as the the connection $`A_a^i`$ can be identified with a one form that could be integrated along a one-dimensional object, one would like analyze the geometric character of the densitised triad $`\stackrel{~}{E}`$ in order to naturally define a smeared object. Recall that the momentum is a density-one vector field on $`\mathrm{\Sigma }`$, $`\stackrel{~}{E}_i^a`$ with values in the dual of the lie-algebra $`su(2)`$. In terms of its tensorial character, it is naturally dual to a (lie-algebra valued) two form,
$$E_{abi}:=\frac{1}{2}\stackrel{~}{}\eta {}_{abc}{}^{}\stackrel{~}{E}_{i}^{c}$$
(17)
where $`\stackrel{~}{}\eta _{abc}`$ is the naturally defined Levi-Civita symbol. It is now obvious that the momenta is crying to be integrated over a two-surface $`S`$. It is now easy to define the objects
$$E[S,f]:=_SE_{abi}f^idS^{ab},$$
(18)
where $`f^i`$ is a lie-algebra valued smearing function on $`S`$. This ‘Electric flux’ variable does not need a background metric to be defined, and it transforms again covariantly as was the case of the holonomies. The algebra generated by holonomies and flux variables is known as the Holonomy-Flux algebra $``$.
Perhaps the main reason why this Holonomy-Flux algebra $``$ is interesting, is the way in which the basic generators interact, when considering the classical (Poisson) lie-bracket. First, given that the configuration functions depend only on the connection and the connections Poisson-commute, one expects that $`\{T[\alpha ],T[\beta ]\}=0`$ for any loops $`\alpha `$ and $`\beta `$. The most interesting poisson bracket one is interested in is the one between a configuration and a momenta variable,
$$\{T[\alpha ],E[S,f]\}=\kappa \underset{\mu }{}f^i(v)\iota (\alpha ,S|_v)\mathrm{Tr}(\tau _ih(\alpha ))$$
(19)
where the sum is over the vertices $`v`$ and $`\iota (\alpha ,S|_v)=\pm 1`$ is something like the intersection number between the loops $`\alpha `$ and the surface $`S`$ at point $`v`$. The sum is over all intersection of the loop $`\alpha `$ and the surface $`S`$. The most important property of the Poisson Bracket is that it is completely topological. This has to be so if we want to have a fully background independent classical algebra for the quantization.
A remark is in order. The value of the constant $`\iota |_v`$ depends not only on the relative orientation of the tangent vector of the loop $`\alpha `$ with respect to the orientation of $`\mathrm{\Sigma }`$ and $`S`$, but also on a further decomposition of the loop into edges, and whether they are ‘incoming’ or ‘outgoing’ to the vertex $`v`$. The end result is that is we have, for simple intersections, that the number $`\iota |_v`$ becomes insensitive to the ‘orientation’. This is different to the $`U(1)`$ case where the final result is the intersection number. For details see ACZ (and note the difference with the claims in Nicolai:2005mc ).
Let us now consider the slightly more involved case of a cylindrical function $`C_\mathrm{{\rm Y}}`$ that is defined over a graph $`\mathrm{{\rm Y}}`$ with edges $`e_I`$, intersecting the surface $`S`$ at points $`p`$. We have then,
$$\{C_\gamma ,E[S,f]\}=\frac{\kappa }{2}\underset{p}{}\underset{I_p}{}\iota (I_p)f^i(p)X_{I_p}^ic$$
(20)
where the sum is over the vertices $`p`$ of the graph that lie on the surface $`S`$, $`I_p`$ are the edges starting or finishing in $`p`$ and where $`X_{I_P}^ic`$ is the result of the action of the $`i`$-th left (resp. right) invariant vector field on the $`I_p`$-th copy of the group if the $`I_p`$-th edge is pointing away from (resp. towards) the surface $`S`$. Note the structure of the right hand side. The result is non-zero only if the graph $`\mathrm{{\rm Y}}`$ used in the definition of the configuration variable $`C_\mathrm{{\rm Y}}`$ intersects the surface $`S`$ used to smear the triad. If the two intersect, the contributions arise from the action of right/left invariant vector fields on the arguments of $`c`$ associated with the edges at the intersection.
Finally, the next bracket we should consider is between two momentum functions, namely $`\{E[S,f],E[S^{},g]\}`$. Just as in the case of holonomies, these functions depend only on one of the canonical variables, namely the triad $`\stackrel{~}{E}`$. One should then expect that their Poisson bracket vanishes. Surprisingly, this is not the case and one has to appropriately define the correct algebraic structure<sup>1</sup><sup>1</sup>1The end result is that one should not regard $`E[S,f]`$ as phase space functions subject to the ordinary Poisson bracket relations, but rather should be viewed as arising from vector fields $`X^\alpha `$ on $`𝒜`$. The non-trivial bracket is then due to the non-commutative nature of the corresponding vector fields. This was shown in ACZ where details can be found.
We have arrived then to the basic variables that will be used in the quantization in order to arrive at LQG. They are given by,
$$h(e_I)\mathrm{Configuration}\mathrm{function}$$
(21)
and
$$E[S,f]\mathrm{Momentum}\mathrm{function},$$
(22)
subject to the basic Poisson bracket relations given by Eqs. (19) and (20). In the next section we shall take the Holonomy-Flux algebra $``$ as the starting point for the quantization.
## IV The Hilbert Space
In this section we shall first outline the quantization strategy to arrive at the kinematical Hilbert space $`_{\mathrm{kin}}`$. Later on, we shall give certain details of the several parts that arise in the construction. As we have emphasized, the Kinematical Hilbert space is the starting point for the program of implementing the constraints as quantum operators. We shall consider briefly that issue later on.
### IV.1 General Considerations
The strategy is to build the Hilbert space as in the ordinary Schrödinger representation, where states are to be represented by wave-functions of the configuration space. Recall that in the present approach, one has decided that the space $`𝒜`$ is the space of (classical) configurations. The natural strategy is then to consider wave functions $`\mathrm{\Psi }`$ of the form,
$$\mathrm{\Psi }=\mathrm{\Psi }(A)$$
(23)
Following the analogy, one should expect that the Hilbert space will consist of wave-functions that are square integrable. This means that one has to introduce a measure $`\mathrm{d}\mu `$ on the space of connections, in order to define the inner product heuristically as,
$$\mathrm{\Phi }|\mathrm{\Psi }=d\mu \overline{\mathrm{\Phi }}\mathrm{\Psi }$$
(24)
The Hilbert space would be denoted then as $`_{\mathrm{kin}}=L^2(𝒜,\mathrm{d}\mu )`$. In this configuration representation one expects that the connection and any function of it will act as a multiplication operator, and that the momenta will be represented as a derivation, with possible a correction term<sup>2</sup><sup>2</sup>2For a discussion of the corresponding Schrödinger representation for a scalar field see CCQ ..
There are several different technical issues that need to be properly addressed in order to complete the construction as we have outlined it. Basically what needs to be done is to give mathematical meaning to the different aspects of the kinematical Hilbert space. Let us then make some general remarks for each of them.
i) Configuration space. From the experience that has been gained from the scalar field case, we know that the space where the wave functions have support is a much larger (functional) space as the space of classical configurations. In the scalar case, the classical configuration space involves smooth fields, whereas the quantum configuration space is made of tempered distributions (dual to the Schwarz space). In our case we expect that the quantum configuration space $`\overline{𝒜}`$ will also be an extension of the classical space. This is indeed the case as was shown in the mid-nineties by Ashtekar and collaborators. There are several characterizations of this space (that can be found in AL:review ; Thiemann:2001yy ; Velhinho ), but here we shall mention only two of them (for a third one see below). The first one is to recall that one could think of a connection as an operator that acts on 1-dimensional objects (edges) and gives a group element. Certainly, a smooth connection has this property. It turns out that a generalized connection, an element of the quantum configuration space $`\overline{𝒜}`$, is precisely any such map (satisfying some tame conditions such as the composition: $`h(e_1e_2)=h(e_1)h(e_2);,e_I`$). In particular, what is dropped is any notion of continuity of the connection. The other characterization is more algebraic. It is based on the observation that the (Abelian) algebra of cylindrical functions Cyl can be extended to a $`C^{}`$-algebra $`\overline{\mathrm{Cyl}}`$ with unit. Standard theorems on representations of such algebras tell us that those algebras can be viewed as the space of continuous functions over a compact space $`\mathrm{\Delta }`$, called the spectrum of the algebra. It turns out that $`\overline{𝒜}`$ is precisely the spectrum of the algebra $`\overline{\mathrm{Cyl}}`$.
ii) Measure. A extremely important issue in the choice of quantum theory is the inner product on the Hilbert space. In the case of wave-functions this implies defining a measure on the configuration space $`\overline{𝒜}`$. At his point one should require that the measure –responsible in a sense for the resulting quantum representation– be diffeomorphism invariant. This would implement naturally the intuitive notion that the theory should be spatially diffeo-invariant. The construction of a measure with this property is one of the main results of the program since it was generally thought that no such measures existed for gauge theories. The diffeo-invariant measure is known as the Ashtekar-Lewandowski measure $`\mu _{\mathrm{AL}}`$.
iii) Momentum operators. The other important part of the quantum representation is of course, the way in which the functions $`E[S,f]`$ are promoted to operators on the Hilbert space $`_{\mathrm{kin}}`$. Intuitively, whenever the wave function $`\mathrm{\Psi }`$ is a function of the configuration variable (the connection), the conjugate variable acts as derivation $`\widehat{E^a}=\delta /\delta A_a`$. However, one should also consider the possibility of adding a multiplicative term to the derivative. This term would commute with the (multiplicative) configuration operator $`\widehat{h}(e_I)`$, so the commutation relations would be satisfied. This extra term would also account, for instance, for the existence of a non-trivial measure (See CCQ for a discussion in the scalar case). Therefore, the choice of measure and the representation of the momentum operator are intertwined, and one should ensure that the measure not only be diffeomorphism invariant, but should also “support” the triad. A systematic study of these issues initiated by Sahlmann has yielded a uniqueness result: The only diffeo-covariant representation of the Holonomy-Flux algebra $``$ is given by the Ashtekar-Lewandowski representation LOST .
iv) Hybrid variables. Recall that the basic functions we have chosen for the quantization are the pairs $`(h(e_I),E[S,f])`$. These generators of the classical algebra to be quantized (the Holonomy-Flux algebra), are not canonical. One of them, namely the holonomy, is a exponentiated version of the connection, whereas the electric flux is linear in the triad. This would be equivalent to a choice $`U(\lambda )=\mathrm{exp}[i\lambda q]`$ and $`\pi (\nu )=\nu p`$, for a quantum mechanical system with phase space $`\mathrm{\Gamma }=(q,p)`$. This means that the basic functions of phase space are neither of the canonical type $`(q,p)`$, nor of the Weyl form (both exponentiated); the quantization based on this functions is in a strict sense non-canonical. In finite dimensions the Stone von Neumann theorem assures us that the Weyl relations are equivalent to the Canonical Commutation Relations, provided the quantum representations are regular<sup>3</sup><sup>3</sup>3For a nice discussion see for instance wald2 .. For a representation to be regular means that the exponentiated form of the operators –which are unitary– are continuous with respect to the parameters ($`\lambda `$ in the $`U(\lambda )`$ above). One of the main features of the loop quantization is that the exponentiated variable, the holonomy, has a well defined associated operator $`\widehat{h}(e_I)`$ on $`_{\mathrm{kin}}`$ that is however, discontinuous. Thus, the operator $`\widehat{A}_i^a`$ does not exist! In the finite-dimensional system example, the equivalent statement would be to say that $`\widehat{U}(\lambda )`$ is well defined but $`\widehat{q}`$ is not. Such non-regular representations in quantum mechanics do exist poly and form the basic starting point for Loop Quantum Cosmology LQC ; bojo .
### IV.2 Ashtekar-Lewandowski Hilbert Space
Let us now consider the particular representation that defines LQG. As we have discussed before, the basic observables are represented as operators acting on wave functions $`\mathrm{\Psi }_\mathrm{{\rm Y}}(\overline{A})_{\mathrm{kin}}`$ as follows:
$$\widehat{h}(e_I)\mathrm{\Psi }(\overline{A})=\left(h(e_I)\mathrm{\Psi }\right)(\overline{A})$$
(25)
and
$$\widehat{E}[S,f]\mathrm{\Psi }_\mathrm{{\rm Y}}(\overline{A})=i\mathrm{}\{\mathrm{\Psi }_\mathrm{{\rm Y}},{}_{}{}^{2}E[S,f]\}=i\frac{\mathrm{}_\mathrm{P}^2}{2}\underset{p}{}\underset{I_p}{}\kappa (I_p)f^i(p)X_{I_p}^i\psi $$
(26)
where $`\mathrm{}_\mathrm{P}^2=\kappa \mathrm{}=8\pi G\mathrm{}`$, the Planck area is giving us the scale of the theory (recall that the Immirzi parameter $`\gamma `$ does not appear in the basic Poisson bracket, and should therefore not play any role in the quantum representation). Here we have assumed that a cylindrical function $`\mathrm{\Psi }_\mathrm{{\rm Y}}`$ on a graph $`\mathrm{{\rm Y}}`$ is an element of the Kinematical Hilbert space (which we haven’t defined yet!). This implies one of the most important assumption in the loop quantization prescription, namely, that objects such as holonomies and Wilson loops that are smeared in one dimension are well defined operators on the quantum theory<sup>4</sup><sup>4</sup>4this has to be contrasted with the ordinary Fock representation where such objects do not give raise to well defined operators on Fock space. This implies that the loop quantum theory is qualitatively different from the standard quantization of gauge fields..
Let us now introduce the third characterization of the quantum configuration space which will be useful for constructing the Hilbert space $`_{\mathrm{kin}}`$. The basic idea for the construction of both the Hilbert space $`_{\mathrm{kin}}`$ (with its measure) and the quantum configuration space $`\overline{𝒜}`$, is to consider the projective family of all possible graphs on $`\mathrm{\Sigma }`$. For any given graph $`\mathrm{{\rm Y}}`$, we have a configuration space $`𝒜_\mathrm{{\rm Y}}=(SU(2))^N`$, which is $`n`$-copies of the (compact) gauge group $`SU(2)`$. Now, it turns out that there is a preferred (normalized) measure on any compact semi-simple Lie group that is left and right invariant. It is known as the Haar measure $`\mu _\mathrm{H}`$ on the group. We can thus endow $`𝒜_\mathrm{{\rm Y}}`$ with a measure $`\mu _\mathrm{{\rm Y}}`$ that is defined by using the Haar measure on all copies of the group. Given this measure on $`𝒜_\mathrm{{\rm Y}}`$, we can consider square integrable functions thereon and with them the graph-$`\mathrm{{\rm Y}}`$ Hilbert space $`_\mathrm{{\rm Y}}`$, which is of the form:
$$_\mathrm{{\rm Y}}=L^2(𝒜_\mathrm{{\rm Y}},\mathrm{d}\mu _\mathrm{{\rm Y}})$$
(27)
If we were working with a unique, fixed graph $`\mathrm{{\rm Y}}_0`$ (which we are not), we would be in the case of a lattice gauge theory on an irregular lattice. The main difference between that situation and LQG is that, in the latter case, one is considering all graphs on $`\mathrm{\Sigma }`$, and one has a family of configurations spaces $`\{𝒜_\mathrm{{\rm Y}}/\mathrm{{\rm Y}}\mathrm{a}\mathrm{graph}\mathrm{in}\mathrm{\Sigma }\}`$, and a family of Hilbert spaces $`\{_\mathrm{{\rm Y}}/\mathrm{{\rm Y}}\mathrm{a}\mathrm{graph}\mathrm{in}\mathrm{\Sigma }\}`$. In order to have consistent families of configuration and Hilbert spaces one needs some conditions. In particular, there is a notion of when a graph $`\mathrm{{\rm Y}}`$ is ‘larger’ than $`\mathrm{{\rm Y}}^{}`$. We say that if $`\mathrm{{\rm Y}}`$ contains $`\mathrm{{\rm Y}}^{}`$ then $`\mathrm{{\rm Y}}>\mathrm{{\rm Y}}^{}`$.<sup>5</sup><sup>5</sup>5by containing we mean that the larger graph can be obtained from the smaller by adding some edges or by artificially dividing the existing edges (by proclaiming that an interior point $`pe_I`$ is now a vertex of the graph).
Given this (partial) relation “$`>`$”, we have a corresponding projection $`P_{\mathrm{{\rm Y}},\mathrm{{\rm Y}}^{}}:\mathrm{{\rm Y}}\mathrm{{\rm Y}}^{}`$, which in turn induces a projection operator for configuration spaces $`P:𝒜_\mathrm{{\rm Y}}𝒜_\mathrm{{\rm Y}}^{}`$ and an inclusion operator for Hilbert spaces $`\iota :_\mathrm{{\rm Y}}^{}_\mathrm{{\rm Y}}`$. The consistency conditions that need to be imposed are rather simple. The intuitive idea is that, if one has a graph $`\mathrm{{\rm Y}}_o`$ and a function thereon $`\psi _\mathrm{{\rm Y}}[A]`$ one should be able to consider larger graphs $`\mathrm{{\rm Y}}_i`$ where the function $`\psi `$ be well defined; one should be able to talk of the ‘same function’, even when defined on a larger graph, and more importantly, its integral (and inner product with other functions) should be independent of the graph $`\mathrm{{\rm Y}}_i`$ we have decided to work on.
We are now in the position of giving a heuristic definition of the configuration space $`\overline{𝒜}`$ and $`_{\mathrm{kin}}`$: The quantum configuration space $`\overline{𝒜}`$ is the configuration space for the “largest graph”; and similarly, the kinematical Hilbert space $`_{\mathrm{kin}}`$ is the largest space containing all Hilbert spaces in $`\{_\mathrm{{\rm Y}}/\mathrm{{\rm Y}}\mathrm{a}\mathrm{graph}\mathrm{in}\mathrm{\Sigma }\}`$. Of course, this can be made precise mathematically, where the relevant limits are called projective. For details see alm2t ; AL:review and Thiemann:2001yy . The Ashtekar-Lewandowski measure $`\mu _{\mathrm{AL}}`$ on $`_{\mathrm{kin}}`$ is then the measure whose projection to any $`𝒜_\mathrm{{\rm Y}}`$ yields the corresponding Haar measure $`\mu _\mathrm{{\rm Y}}`$. The resulting Hilbert space can thus be written as
$$_{\mathrm{kin}}=L^2(\overline{𝒜},\mathrm{d}\mu _{\mathrm{AL}})$$
Let us now see how it is that the cylindrical functions $`\mathrm{\Psi }_\mathrm{{\rm Y}}`$ Cyl belong to the Hilbert space of the theory. Let us consider a cylindrical function $`\mathrm{\Psi }_\mathrm{{\rm Y}}`$ defined on the space $`𝒜_\mathrm{{\rm Y}}`$. If it is continuous, then it is bounded (since $`𝒜_\mathrm{{\rm Y}}`$ is compact), and thus it is square integrable with respect to the measure $`\mu _\mathrm{{\rm Y}}`$. Therefore, $`\mathrm{\Psi }_\mathrm{{\rm Y}}_\mathrm{{\rm Y}}`$. Finally, we have the inclusion between Hilbert spaces given by $`_\mathrm{{\rm Y}}_{\mathrm{kin}}`$ which implies that the cylinder function $`\mathrm{\Psi }_\mathrm{{\rm Y}}`$ indeed belongs to the kinematical Hilbert space.
Let us now see how this inner product works for known functions such as Wilson loops. Consider $`T[\alpha ]`$ and $`T[\beta ]`$ two Wilson loops with $`\alpha \beta `$ (and nonintersecting, for simplicity). Each loop can be regarded as a graph on itself, with only an edge and no vertex (or many bi-valent vertices by artificially declaring them to be there). In order to take the inner product between these two function which looks like
$$T[\alpha ]|T[\beta ]_{\mathrm{kin}}=_{\overline{𝒜}}d\mu _{\mathrm{AL}}\overline{T[\alpha ]}T[\beta ]$$
(28)
we have to construct a graph $`\mathrm{{\rm Y}}`$ that contains both $`\alpha `$ and $`\beta `$. This is rather easy to do and in fact there is large freedom in doing that. The end result should be independent of the particular choice. The simplest possibility is to take an edge $`e`$ that connects any point of $`\alpha `$ with any point of $`\beta `$ (and for simplicity it does not intersect neither loop at another point). The resulting graph has now three edges $`(\alpha ,\beta ,e)`$, so we can construct its configuration space to be $`𝒜_\mathrm{{\rm Y}}=(SU(2))^3=`$. The Wilson loops are now very simple cylinder functions on $`𝒜_\mathrm{{\rm Y}}`$ (or rather, the induced function by the inclusion of $`\alpha `$ into $`\mathrm{{\rm Y}}`$). $`T[\alpha ]`$ has associated a function $`t_\alpha :(SU(2))^3`$ that only depends on the first argument (corresponding to $`\alpha `$), whereas $`T[\beta ]`$ has a function $`t_\beta `$ that depends only on the second argument .Its functional dependence (the same for $`T[\alpha ]`$ and $`T[\beta ]`$), as a function of the holonomy in each edge, is very simple:
$$T[\alpha ]=t_\alpha (h(\alpha ),h(\beta ),h(e))=\frac{1}{2}\mathrm{Tr}(h(\alpha )).$$
Given that we are integrating a function ($`\overline{T[\alpha ]}T[\beta ]`$) that is defined on the graph $`\mathrm{{\rm Y}}=\alpha \beta e`$, and that the functional integral over the full space $`\overline{𝒜}`$ reduces to an integral over the minimal graph were the function is defined, the inner product (28) can then be rewritten as,
$$T[\alpha ]|T[\beta ]_{\mathrm{kin}}=_{G_1}d\mu _\mathrm{H}\overline{T[\alpha ]}_{G_2}d\mu _\mathrm{H}T[\beta ]_{G_3}d\mu _\mathrm{H}$$
(29)
Each integral is performed over each copy of the gauge group, where $`G_1`$ denotes the first entry in the functions $`t_i`$, namely the holonomy along $`\alpha `$, and $`G_2`$ the holonomy along $`\beta `$. The third copy of the group, that associated to $`e`$ has the property that the function we are integrating has no dependence on it, thus one is only integrating the unit function. We are assuming that $`\mu _H`$ is normalized so this integral is equal to one. The question now reduces to that of calculating the integral:
$$_{G_1}d\mu _\mathrm{H}\overline{T[\alpha ]}=\frac{1}{2}_{G_1}d\mu _\mathrm{H}\overline{\mathrm{Tr}(h(\alpha ))}$$
(30)
It turns out that the Haar measure is such that this integral vanishes. Thus, the inner product (28) is zero, for any two (different) loops $`\alpha `$ and $`\beta `$ and for any edge $`e`$ connecting them. Note that if we had assumed from the very beginning that the loops $`\alpha `$ and $`\beta `$ are equal, then there would be no need to define a connecting edge and the inner product would be given by only one integral $`_{G_1}d\mu _\mathrm{H}\overline{T[\alpha ]}T[\alpha ]=\frac{1}{2}_{G_1}d\mu _\mathrm{H}|\mathrm{Tr}(h(\alpha ))|^20`$
Several remarks are in order:
i) The fact that the inner product between any two Wilson loop functions vanishes is a clear signature that the measure $`\mu _{\mathrm{AL}}`$ is diffeomorphism invariant. More precisely, a diffeomorphism $`\varphi `$ induces a transformation $`U(\varphi ):`$ given by: $`U(\varphi )\mathrm{\Psi }_\mathrm{{\rm Y}}=\mathrm{\Psi }_{(\varphi ^1\mathrm{{\rm Y}})}`$. The measure $`\mu _{\mathrm{AL}}`$ is such that the operator $`U(\lambda )`$ is unitary. Thus, diffeomorphism are unitarily implemented in this quantum theory.
ii) The integral (30) can be interpreted as the vacuum expectation value of the operator $`\widehat{T}[\alpha ]`$; The “vacuum” $`\mathrm{\Psi }_0`$ in this representation is simply the unit function, so we have:
$$\widehat{T}[\alpha ]_0=_GT[\alpha ]d\mu _\mathrm{H}=0.$$
We can then conclude that in the LQG representation, the vacuum expectation value of all Wilson loops vanishes exactly.
iii) As we mentioned before, the kinematical Hilbert space $`_{\mathrm{kin}}`$ can be regarded as the “largest” Hilbert space by combining the Hilbert spaces over all possible graphs. We shall use the following symbol to denote that idea. We write then:
$$_{\mathrm{kin}}=_\mathrm{{\rm Y}}_\mathrm{{\rm Y}}$$
Note that since the space of graphs $`\mathrm{{\rm Y}}`$ is uncountable, the Hilbert space $`_{\mathrm{kin}}`$ is non-separable. In the next part we shall se how a convenient choice of basis for each Hilbert space $`_\mathrm{{\rm Y}}`$ will allow us to have a full decomposition of the Hilbert space.
Let us then recall what is the structure of simple states in the theory. The vacuum or ‘ground state’ $`|0`$ is given by the unit function. One can then create excitations by acting via multiplication with holonomies or Wilson loops. The resulting state $`|\alpha =\widehat{T}[\alpha ]|0`$ is an excitation of the geometry but only along the one dimensional loop $`\alpha `$. Since the excitations are one dimensional, the geometry is sometimes said to be polymer like. In order to obtain a geometry that resembles a three dimensional continuum one needs aa huge number of edges ($`10^68`$) and vertices.
### IV.3 A choice of basis: Spin Networks
The purpose of this part is to provide a useful decomposition of the Hilbert space $`_\mathrm{{\rm Y}}`$, for all graphs. In practice one just takes one particular graph $`\mathrm{{\rm Y}}_o`$ and works in that graph. This would be like restricting oneself to a fixed lattice and one would be working on the Hilbert space of a Lattice Gauge Theory. Thus, all the results of this subsection are also relevant for a LGT, but one should keep in mind that one is always restricting the attention to a little part of the total Hilbert space $`_{\mathrm{kin}}`$.
Let us begin by sketching the basic idea of what we want to do. For simplicity, let us consider just one edge, say $`e_i`$. What we need to do is to be able to decompose any function $`F`$ on $`G`$ (in this case we only have one copy of the group), in a suitable basis. In the simplest case, when the group is just $`G=U(1)`$ we just have a circle. In this case, the canonical decomposition for any function $`F`$ is given by the Fourier series: $`F=_nA_n\mathrm{e}^{\mathrm{i}n\theta }`$, with $`A_n=\frac{1}{2\pi }F(\theta )\mathrm{e}^{\mathrm{i}n\theta }d\theta `$. As is well known, the functions $`\mathrm{e}^{\mathrm{i}n\theta }`$ represent a basis for any function on the group $`G=U(1)`$, that is also orthonormal with respect to the measure $`\mathrm{d}\mu _\mathrm{H}=\frac{1}{2\pi }\mathrm{d}\theta `$, which is nothing but the Haar measure on this group. For notational simplicity one can denote by $`f_n(\theta )=\mathrm{e}^{\mathrm{i}n\theta }`$, and they represent irreducible representations of the group $`U(1)`$, labelled by $`n`$. Then the series looks like $`F(\theta )=_nA_nf_n(\theta )`$. Here we should think of assigning, to each of the basis functions $`f_n`$ the trivial label $`n`$ (for all integers), and the general function is a linear combination of the $`n`$-labelled basis.
In the case of the group $`G=SU(2)`$, there is en equivalent decomposition of a function $`f(g)`$ of the group ($`gG`$). It reads,
$$f(g)=\underset{j}{}\sqrt{j(j+1)}f_j^{mm^{}}\underset{mm^{}}{\overset{j}{\mathrm{\Pi }}}(g)$$
(31)
where,
$$f_j^{mm^{}}=\sqrt{j(j+1)}_Gd\mu _H\underset{mm^{}}{\overset{j}{\mathrm{\Pi }}}(g^1)f(g)$$
(32)
is the equivalent of the Fourier component. When doing harmonic analysis on groups the generalization of the Fourier decomposition is known as the Peter-Weyl decomposition. The functions $`\underset{mm^{}}{\overset{j}{\mathrm{\Pi }}}(g)`$ play the role of the Fourier basis. In this case these are unitary representation of the group, and the label $`j`$ labels the irreducible representations. In the $`SU(2)`$ case with the interpretation of spin, these represent the spin-$`j`$ representations of the group. In our case, we will continue to use that terminology (spin) even when the interpretation is somewhat different.
Given a cylindrical function $`\mathrm{\Psi }_\mathrm{{\rm Y}}[A]=\psi (h(e_1),h(e_2),\mathrm{},h(e_N))`$, we can then write an expansion for it as,
$`\mathrm{\Psi }_\mathrm{{\rm Y}}[A]`$ $`=`$ $`\psi (h(e_1),h(e_2),\mathrm{},h(e_N))`$ (33)
$`=`$ $`{\displaystyle \underset{j_1\mathrm{}j_N}{}}f_{j_1\mathrm{}j_N}^{m_1\mathrm{}m_N,n_1\mathrm{}n_N}\varphi _{m_1n_1}^{j_1}(h(e_1))\mathrm{}\varphi _{m_Nn_N}^{j_N}(h(e_N)),`$
where $`\varphi _{mn}^j(g)=\sqrt{j(j+1)}\underset{mn}{\overset{j}{\mathrm{\Pi }}}(g)`$ is the normalized function satisfying
$$_Gd\mu _\mathrm{H}\overline{\varphi _{mn}^j(g)}\varphi _{m^{}n^{}}^j^{}(g)=\delta _{j,j^{}}\delta _{m,m^{}}\delta _{n,n^{}}.$$
The expansion coefficients can be obtained by projecting the state $`|\mathrm{\Psi }_\mathrm{{\rm Y}}`$,
$$f_{j_1\mathrm{}j_N}^{m_1\mathrm{}m_N,n_1\mathrm{}n_N}=\varphi _{m_1n_1}^{j_1}\mathrm{}\varphi _{m_Nn_N}^{j_N}|\mathrm{\Psi }_\mathrm{{\rm Y}}$$
(34)
This implies that the products of components of irreducible representations $`_{i=1}^N\varphi _{m_in_i}^{j_i}[h(e_i)]`$ associated with the $`N`$ edges $`e_I\mathrm{{\rm Y}}`$, for all values of spins $`j`$ and for $`jm,nj`$ and for any graph $`\mathrm{{\rm Y}}`$, is a complete orthonormal basis for $`_{\mathrm{kin}}`$. We can the write,
$$_\mathrm{{\rm Y}}=_j_{\mathrm{{\rm Y}},j}$$
(35)
where the Hilbert space $`_{\alpha ,j}`$ for a single loop $`\alpha `$, and a label $`j`$ is the familiar $`(2j+1)`$ dimensional Hilbert space of a particle of ‘spin $`j`$’. For a complete treatment see Perez:2004hj .
We have been able to decompose the Hilbert space of a given graph, using the Peter-Weyl decomposition theorem, but how can we make contact with the so called spin networks? Let us for a moment focus our attention to the $`U(1)`$ case. Then, for each edge one could decompose any function in a Fourier series where the basis was labelled by an integer, that correspond to all possible irreducible unitary representations of the group. As we saw later, one could think of a basis of the graph Hilbert by considering the product over functions labelled by these representations. What one can do it to associate a label $`n_I`$ for each edge, to denote the function made out of those basis functions corresponding to the chosen labels. A graph with the extra labelling, known as dressing, can be given different names. In the $`U(1)`$ theory the labelling denotes the possible electric flux, so the graph $`\mathrm{{\rm Y}}_{(n_1,n_2,\mathrm{},n_N)}`$ with the labellings for each edge is called a flux network. In the case of geometry with group $`SU(2)`$, the graphs with labelling $`j_I=𝐣`$ are known as spin networks. As the reader might have noticed, in the geometry case there are more labels than the spins for the edges. Normally these are associated to vertices and are known as intertwiners. This means that the Hilbert spaces $`_{\mathrm{{\rm Y}},𝐣}`$ is finite dimensional. Its dimension being a measure of the extra freedom contained in the intertwiners. One could then introduce further labelling $`𝐥`$ for the graph, so we can decompose the Hilbert space as
$$_\mathrm{{\rm Y}}=_j_{\mathrm{{\rm Y}},𝐣}=_{𝐣,𝐥}_{\mathrm{{\rm Y}},𝐣,𝐥}$$
(36)
where now the spaces $`_{\mathrm{{\rm Y}},𝐣,𝐥}`$ are one-dimensional. For more details see baez2 , Sec. 4.1.5 of Perez:2004hj and Sec. 4.2 of AL:review .
Let us see the simplest possible example, the Wilson loop $`T[\alpha ]`$. In this case there is only one edge, and the function only involves the lowest representation $`j=1/2`$. Furthermore, the example is very easy since it is only the trace of the object $`\varphi _{mn}^{1/2}`$. In terms of labelling it corresponds to a $`(j=1/2)`$ label on the loop. Since there are no vertices, there is no choice for ‘intertwiner’. The next simplest example of a spin network is for the simplest graph, namely a loops $`\alpha `$ but for higher spin $`j`$ labellings. These functions correspond to taking the trace of the holonomy around $`\alpha `$ but in a higher representation $`j`$ of the group. In general a spin-$`j`$ Wilson loop $`T^j[\alpha ]`$ can be written in terms of the fundamental one $`T[\alpha ]`$ as a polynomial expansion of degree $`j`$: $`T^j[\alpha ]=_{k=1,\mathrm{},j}A_k(T[\alpha ])^k`$, for some coefficients $`A_k`$.
Two final remarks are in order. While it is true that spin networks are naturally defined in terms of the harmonic analysis on $`G`$, and due to its properties as an eigen-basis of important operators (the subject of next chapter) have a special place in the theory, they are not, by themselves the theory. Dirac’s transformation theory assures us that there are equivalent descriptions for the quantum theory that are not given by these functions. Another related issue is the existence of the so-called loop representation of the theory. In its early stages, non-perturbative quantum gravity (today known as LQG) was though to have two equivalent representations, the connection representation and the loop representation (precisely in the sense of Dirac’s transformation theory). When the structure of the theory was properly understood and the Hilbert space was characterized, it was clear that the loop representation as originally envisioned was not well defined. The basic idea was to have loops $`\alpha `$ as the argument for wave-functions. Thus, if one had (in Dirac notation) a state $`|\mathrm{\Psi }`$, the wave function of the connection would be $`\mathrm{\Psi }[A]=A|\mathrm{\Psi }`$. If one had a ‘basis of loops’ $`\alpha |`$, one could imagine having $`\alpha |\mathrm{\Psi }`$, the state in the loop representation. Furthermore, there was a proposal for defining this state via the ‘loop transform’:
$$\psi [\alpha ]=\alpha |\mathrm{\Psi }=𝒟A\alpha |AA|\mathrm{\Psi }:=𝒟A\overline{T}[\alpha ,A]\mathrm{\Psi }[A]$$
(37)
It was then implied that the Wilson loops were the kernel of the transformation and were in a sense complete. From the discussion of this section, it is clear that those expectations were superseded. Namely, what we have seen is that loops by themselves are not enough; we would loose information by only considering loops<sup>6</sup><sup>6</sup>6More precisely, when seen as graphs, loops are not enough. What was previously done was to consider integer powers of loops $`\alpha ^n=\alpha \alpha \mathrm{}\alpha `$ in the ‘basis’ of loops. Graphs (that only care about the image of the embedding, and where $`\alpha ^2=\alpha `$) together with spin and intertwiners represent a more efficient (no redundancies) way of characterizing the discrete basis of the theory.. We need the extra information contained in the spin networks, namely we need to consider all possible spins and intertwiners. Instead, what one has to consider is the ‘spin network transform’,
$$\psi [\mathrm{{\rm Y}}_{𝐣,𝐥}]=\mathrm{{\rm Y}}_{𝐣,𝐥}|\mathrm{\Psi }=d\mu _{\mathrm{AL}}\mathrm{{\rm Y}}_{𝐣,𝐥}|AA|\mathrm{\Psi }:=_{\overline{𝒜}}d\mu _{\mathrm{AL}}\overline{𝒩}[\mathrm{{\rm Y}}_{𝐣,𝐥},A]\mathrm{\Psi }[A]$$
(38)
Let us end this part with a two remarks regarding spin networks. So far we have only considered cylindrical functions that satisfy no further requirements. We know, on the other hand that in order to have physical states we need to impose the quantum constraints on the states. In particular Gauss’ law implies that the states be gauge invariant, namely, invariant under the action of (finite) gauge transformations of the connection $`Ag^1Ag+g^1\mathrm{d}g`$. This condition, when translated to the language of holonomies, graphs and so on imposes some restrictions on the cylinder functions, that are conditions imposed only at the vertices of the graph:
$$\underset{v}{}\underset{e_v}{}X_{e_v}^i𝒩[\mathrm{{\rm Y}}_{𝐣,𝐥},A]=0$$
(39)
where the first sum is over vertices $`v`$ of the graph and the second over edges $`e_v`$ for each vertex $`v`$. Spin networks that satisfy these conditions are called gauge invariant spin networks (for a detailed treatment see Sec. 4.1.5 of Perez:2004hj ). Finally, there is a very convenient language for performing calculations involving (gauge invariant) spin networks, that employs graphical manipulations over the vertices of the graphs. This is known as recoupling theory, and (due to space restrictions here) it can be learned from the primers by Major seth and Rovelli Rovelli:1998gg and in Rovelli’s Book Rovelli:2004tv .
In next section we shall explore the picture of the quantum geometry that the mathematical apparatus provides for us.
## V Quantum Geometry
So far, we have constructed the kinematic Hilbert space of the quantum theory. This is the quantum analog of the classical phase space of the theory, the space of space-time metrics. Just as the spacetime, together with its geometrical constructs such as vectors, tensor fields, derivatives, etc. is the arena for doing classical geometry, the space $`_{\mathrm{kin}}`$ is the arena for quantum geometry. What we need is to define the quantum object, i.e. the operators, that will correspond to the geometric objects defined on this arena. Note that the background manifold $`\mathrm{\Sigma }`$ has many of the features of the classical geometry, namely all the structure that comes with $`\mathrm{\Sigma }`$ being a differentiable manifold: tangent spaces, vectors, forms, etc. What is not available are the notions arising from a metric on $`\mathrm{\Sigma }`$: distance, angles, area, volume, etc. These are the properties of the geometry that are quantized. The purpose of this section is to explore this quantum geometry.
This section has three parts. In the first one we describe the nature of the operators associated to electric fluxes. In the second part, we present the simplest of the geometric operators, namely the area of surfaces and in the last part we discuss the non-commutative character of the quantum geometry.
### V.1 Flux operators
The operators $`\widehat{E}[S,f]`$ corresponding to the electric flux observables, are in a sense the basic building blocks for constructing the quantum geometry. We have seen in Sec.III the action of this operators on cylindrical functions,
$$\widehat{E}[S,f]\mathrm{\Psi }_\mathrm{{\rm Y}}(\overline{A})=i\mathrm{}\{\mathrm{\Psi }_\mathrm{{\rm Y}},{}_{}{}^{2}E[S,f]\}=i\frac{\mathrm{}_\mathrm{P}^2}{2}\underset{p}{}\underset{I_p}{}\kappa (I_p)f^i(p)X_{I_p}^i\psi $$
(40)
Here the first sum is over the intersections of the surface $`S`$ with the graph $`\mathrm{{\rm Y}}`$, and the second sum is over all possible edges $`I_p`$ that have the vertex $`v_p`$ (in the intersection of $`S`$ and the graph) as initial of final point. In the simplest case of a loop $`\alpha `$, there are only simple intersections (meaning that there are two edges for each vertex), and in the simplest case of only one intersection between $`S`$ and $`\alpha `$ we have one term in the first sum and two terms in the second (due to the fact that the loop $`\alpha `$ is seen as having a vertex at the intersection point). In this simplified case we have
$$\widehat{E}[S,f]\mathrm{\Psi }_\alpha (\overline{A})=i\mathrm{}_\mathrm{P}^2f^i(p)X_{I_p}^i\psi $$
(41)
Note that the action of the operator is to ‘project’ the angular momentum in the direction given by $`f^i`$ (in the internal space associated with the Lie algebra). As we shall see, this operator is in a sense fundamental the fundamental entity for constructing (gauge invariant) geometrical operators. For this, let us rewrite the action of the flux operator (40), dividing the edges that are above the surface $`S`$, as ‘up’ edges, and those that lie under the surface as ‘down’ edges.
$$\widehat{E}[S,f]=\frac{\mathrm{}_P^2}{2}\underset{p}{}f^i(p)(\widehat{J}_{i(u)}^p\widehat{J}_{i(d)}^p)$$
(42)
where the sum is over the vertices at the intersection of the graph and the surface, and where the ‘up’ operator $`\widehat{J}_{i(u)}^p=\widehat{J}_i^{p,e_1}+\widehat{J}_i^{p,e_2}+\mathrm{}+\widehat{J}_i^{p,e_u}`$ is the sum over all the up edges and the down operator $`\widehat{J}_{i(d)}^p`$ is the corresponding one for the down edges.
### V.2 Area operator
The simplest operator that can be constructed representing geometrical quantities of interest is the area operator, associated to surfaces $`S`$. The reason behind this is again the fact that the densitized triad is dual to a two form that is naturally integrated along a surface. The difference between the classical expression for the area and the flux variable is the fact that the area is a gauge invariant quantity. Let us first recall what the classical expression for the area function is, and then we will outline the regularization procedure to arrive at a well defined operator on the Hilbert space. The area $`A[S]`$ of a surface $`S`$ is given by $`A[S]=_S\mathrm{d}^2x\sqrt{h}`$, where $`h`$ is the determinant of the induced metric $`h_{ab}`$ on $`S`$. When the surface $`S`$ can be parametrized by setting, say, $`x^3=0`$, then the expression for the area in terms of the densitized triad takes a simple form:
$$A[S]=\gamma _S\mathrm{d}^2x\sqrt{\stackrel{~}{E}_i^3\stackrel{~}{E}_j^3k^{ij}}$$
(43)
where $`k^{ij}=\delta ^{ij}`$ is the Killing-Cartan metric on the Lie algebra, and $`\gamma `$ is the Barbero-Immirzi parameter (recall that the canonical conjugate to $`A`$ is $`{}_{}{}^{\gamma }\stackrel{~}{E}_{i}^{a}=\stackrel{~}{E}_i^a/\gamma `$). Note that the functions is again smeared in two dimensions and that the quantity inside the square root is very much a square of the (local) flux. One expects from the experience with the flux operator, that the resulting operator will be a sum over the intersecting points $`p`$, so one should focus the attention to the vertex operator
$$\mathrm{\Delta }_{S,\mathrm{{\rm Y}},p}=\left[(\widehat{J}_{i(u)}^p\widehat{J}_{i(d)}^p)(\widehat{J}_{j(u)}^p\widehat{J}_{j(d)}^p)\right]k^{ij}$$
with this, the area operator takes the form,
$$\widehat{A}[S]=\gamma \mathrm{}_\mathrm{P}^2\underset{p}{}\widehat{\sqrt{\mathrm{\Delta }_{S,\mathrm{{\rm Y}},p}}}$$
(44)
We can now combine both the form of the vertex operator with Gauss’ law $`(\widehat{J}_{i(u)}^p+\widehat{J}_{i(d)}^p)\mathrm{\Psi }=0`$ to arrive at,
$$|(\widehat{J}_{i(u)}^p\widehat{J}_{i(d)}^p)|^2=|2(\widehat{J}_{i(u)}^p)|^2$$
(45)
where we are assuming that there are no tangential edges. The operator $`\widehat{J}_{i(u)}^p`$ is an angular momentum operator, and therefore its square has eigenvalues equal to $`j^u(j^u+1)`$ where $`j^u`$ is the label for the total ‘up’ angular momentum. We can then write the form of the operator
$$\widehat{A}[S]𝒩(\mathrm{{\rm Y}},\stackrel{}{j})=\gamma \mathrm{}_\mathrm{P}^2\underset{vV}{}\sqrt{|\widehat{J}_{i(u)}^p|^2}𝒩(\mathrm{{\rm Y}},\stackrel{}{j})$$
(46)
With these conventions, in the case of simple intersections between the graph $`\mathrm{{\rm Y}}`$ and the surface $`S`$, the area operator takes the well known form:
$$\widehat{A}[S]𝒩(\mathrm{{\rm Y}},\stackrel{}{j})=\gamma \mathrm{}_\mathrm{P}^2\underset{vV}{}\sqrt{j_v(j_v+1)}𝒩(\mathrm{{\rm Y}},\stackrel{}{j})$$
when acting on a spin network $`𝒩(\gamma ,\stackrel{}{j})`$ defined over $`\mathrm{{\rm Y}}`$ and with labels $`\stackrel{}{j}`$ on the edges (we have not used a label for the intertwiners).
Let us now interpret these results in view of the new geometry that the loop quantization gives us. The one dimensional excitations of the geometry carry flux of area: whenever the graph pierces a surface it endows $`S`$ with a quanta of area depending on the value of $`j`$. Furthermore, the eigenvalues of the operator are discrete, giving a precise meaning to the statement that the geometry is quantized: there a minimum (non-zero) value for the area given by taking $`j=1/2`$ in the previous formula. Thus the area gap $`a_\mathrm{o}`$ is given by
$$a_\mathrm{o}=\gamma \mathrm{}_\mathrm{P}^2\frac{\sqrt{3}}{2}$$
(47)
If the value of $`\gamma `$ is of order unit, then we see that the minimum area is of the order of the Planck area. In order to get a macroscopic value for area we would need a very large number of intersections. The Immirzi parameter has to be fixed to select the physical sector of the theory. The current viewpoint is that the black hole entropy calculation can be used for that purpose.
There are operators corresponding to other geometrical objects such as volume, length, angles, etc. The main feature that they exhibit is that their spectrum is always discrete.
### V.3 Non-commutative quantum geometry
One of the somewhat unexpected results coming from loop quantum geometry is the fact that the quantum geometry is inherently non-commutative. This means that the operators associated to geometrical quantities do not, in general, commute.
Let us for simplicity consider area operators. Given a surface $`S`$, any spin network state defined on graphs that do not have vertices on the surface are eigenstates of the area operator (bi-valent vertices are equivalent to no-vertices). There is no degeneracy. If the graph has trivalent vertices on the surface, spin networks are still eigenstates and there is again no degeneracy. It is for a four valent vertex that things start to get interesting. It is not true that any spin network state is an area eigenstate. There is, however, a choice of basis of eigenstates for any given area operator. This is true for any valence, namely that given one surface and any vertex thereon there is always a choice of basis that is an eigenbasis for that surface. There will be, however, degeneracy for higher valence vertices. For a four valent one, we know that for a fixed coloring of the edges, there will be some degeneracy. That is, the (finite dimensional) vector space spanned by the intertwiners can have dimension grater that one. The dimension will be given by the number of possible representations of $`SU(2)`$ compatible with the external colorings (via the triangle inequality). In order to break the degeneracy, it is enough to have one surface operator $`A[S]`$ acting on the vertex, provided that it has two vertices on each side of the surface. That choice of surface will then ‘select’ the decomposition of the vector space into eigenstates of $`A[S]`$, the eigenvalues given by the $`\sqrt{j_{\mathrm{int}}(j_{\mathrm{int}}+1)}`$ kind of formula, where the $`j_{\mathrm{int}}`$ means the ‘internal’ coloring of the representation used to write the intertwiner.
Let us now consider an explicit case of non-commutativity for a graph $`\mathrm{{\rm Y}}`$ having a four valent vertex $`v`$ at the surface $`S`$. We assume that the edges $`e_1`$ and $`e_2`$ are up edges for $`S_1`$ (and $`(e_3,e_4)`$ are down). We denote by $`j_5`$ the internal edge, and we choose as the state a linear combination of spin networks (with coefficients $`(\alpha _0,\alpha _1,\alpha _2)`$) of spin networks having $`j_5=0,j_5=1`$ and $`j_5=2`$ respectively as the internal index: $`\mathrm{\Psi }(\mathrm{{\rm Y}},𝐣_\mathrm{𝟎})=\alpha _0𝒩(\gamma ,j_5=0)+\alpha _1𝒩(\gamma ,j_5=1)+\alpha _2𝒩(\gamma ,j_5=2)=_i\alpha _i𝒩_i`$. The external edges $`(j_1,j_2,j_3,j_4)`$ are all taken to be equal to one. For simplicity let us denote by $`𝒩_i`$ the eigenstates of the area operator: $`\widehat{A}[S_1]𝒩_i=a_i𝒩_i`$, where $`a_0=0,a_1=\gamma \mathrm{}_\mathrm{P}^2\sqrt{2}`$ and $`a_3=\gamma \mathrm{}_\mathrm{P}^2\sqrt{6}`$.
First we need to compute the action of the first surface $`\widehat{A}[S_1]`$,
$$\widehat{A}[S_1]\mathrm{\Psi }(\mathrm{{\rm Y}},𝐣_\mathrm{𝟎})=\underset{i}{}\alpha _ia_i𝒩_i$$
(48)
We know that in this case the vector space of intertwiners is three dimensional (spanned by $`𝒩_i`$), so we can think of the area operator acting on this space as a matrix $`A_j^i=a_i\delta _j^i`$ (no summation over $`i`$).
Before acting with $`\widehat{A}[S_2]`$ we need to change basis from the basis diagonal on $`j_5`$ to the one diagonal to $`j_6`$. Thus we need to expand $`𝒩_i`$ in terms of the other basis. The $`3\times 3`$ matrix $`U_{}^{j}{}_{k}{}^{}`$ that changes basis has to be an element of $`SO(3,)`$. Its precise form is not relevant in what follows but it is known that its components are given by $`6j`$ symbols Rovelli:2004tv . We can then write,
$$𝒩_j=U_{}^{k}{}_{j}{}^{}\stackrel{~}{𝒩}_k$$
(49)
We can now compute
$`\widehat{A}[S_2]\widehat{A}[S_1]\mathrm{\Psi }(\mathrm{{\rm Y}},𝐣_\mathrm{𝟎})`$ $`=`$ $`\widehat{A}[S_2]{\displaystyle \underset{i,k}{}}\alpha _ia_iU_{}^{k}{}_{i}{}^{}\stackrel{~}{𝒩}_k`$ (50)
$`=`$ $`{\displaystyle \underset{i,k}{}}\alpha _ia_iU_{}^{k}{}_{i}{}^{}a_k\stackrel{~}{𝒩}_k`$ (51)
Let us now act with the operators in the reverse order, namely let us act with $`\widehat{A}(S_2)`$ first,
$$\widehat{A}[S_2]\mathrm{\Psi }(\mathrm{{\rm Y}},𝐣_\mathrm{𝟎})=\underset{i,k}{}\alpha _iU_{}^{k}{}_{i}{}^{}a_k\stackrel{~}{𝒩}_k$$
(52)
Now, before acting with $`\widehat{A}[S_1]`$ we need to change basis from the basis diagonal on $`j_6`$ to the basis diagonal to $`j_5`$. Thus we need to expand $`\stackrel{~}{𝒩}_i`$ in terms of the other basis. The matrix that defines this change of basis will be the inverse $`U_{}^{(1)k}{}_{j}{}^{}`$ such that $`\stackrel{~}{𝒩}_j=U_{}^{(1)k}{}_{j}{}^{}𝒩_k`$. We can now compute
$`\widehat{A}[S_1]\widehat{A}[S_2]\mathrm{\Psi }(\mathrm{{\rm Y}},𝐣_\mathrm{𝟎})`$ $`=`$ $`\widehat{A}[S_1]{\displaystyle \underset{i,k,m}{}}\alpha _iU_{}^{k}{}_{i}{}^{}a_iU_{}^{(1)m}{}_{k}{}^{}𝒩_m`$ (53)
$`=`$ $`{\displaystyle \underset{i,k,m}{}}\alpha _iU_{}^{k}{}_{i}{}^{}a_iU_{}^{(1)m}{}_{k}{}^{}a_m𝒩_m`$ (54)
We can then write the commutator as,
$$[\widehat{A}[S_1],\widehat{A}[S_2]]\mathrm{\Psi }(\mathrm{{\rm Y}},𝐣_\mathrm{𝟎})=\underset{i,k,m}{}\alpha _i\left(U_{}^{k}{}_{i}{}^{}a_iU_{}^{(1)m}{}_{k}{}^{}a_ma_iU_{}^{k}{}_{i}{}^{}a_kU_{}^{(1)m}{}_{k}{}^{}\right)𝒩_m$$
(55)
Which is, in general, different from zero (for instance, the area matrix $`A_j^i=a_i\delta _j^i`$ does not commute with $`U_{}^{i}{}_{k}{}^{}`$). This implies clearly that there is an uncertainty relation for intersecting areas
$$(\mathrm{\Delta }\widehat{A}[S_1])^2(\mathrm{\Delta }\widehat{A}[S_2])^20$$
For more details on non-commutativity see non-comm .
## VI Results and open issues
So far we have dealt with the kinematical aspects of the theory, in which the states are (gauge invariant) states labelled by embedded graphs on the three manifold $`\mathrm{\Sigma }`$. However, the fact that general relativity is diffeomorphism invariant as a classical field theory provides a good motivation for trying to implement diffeomorphism invariance at the quantum level. In the $`3+1`$ setting this means implementing the constraints as quantum conditions on the quantum states.
Let us first consider the diffeomorphism constraint, that is, the generator of spatial diffeomorphisms. In the canonical setting, this means that we should ask for diffeomorphism invariant states. The answer is rather easy: the only diffeomorphism invariant state in $`_{\mathrm{kin}}`$ is the vacuum $`\mathrm{\Psi }_0`$! This means that the intuitive idea that diff-invariant states as obtained by taking the quotient by the diffeomorphisms in the Hilbert space does not work. One needs to find a subtler prescription. Fortunately, it is now well understood that, in general, solutions to the quantum constraints do not lie on the original kinematical Hilbert space where the constraints were originally posed. Instead, one has to look for solutions in a larger space. In the case of LQG, diffeo-invariant states live in $`(\mathrm{Cyl})^{}`$, the dual space to the cylindrical functions, where a solution is found by summing over all possible diffeomorphism related states (recall that the action of the diffeomorphisms is a unitary map in the kinematical Hilbert space). Even when the sum is over an uncountable number of states, as a distribution its action on cylindrical functions is well defined. The procedure also includes specifying a diffeo-invariant inner product on solutions, and by completing, one arrives at $`_{\mathrm{diff}}`$, the Hilbert space of diffeomorphism invariant states. Roughly speaking, instead of considering states labelled by embedded graphs as in $`_{\mathrm{kin}}`$, diffeo-invariant states are labelled by diffeomorphism classes of graphs (apart from the internal labellings $`(𝐣,𝐢)`$). Having solved the diffeomorphism invariance, in a rigorous manner and without anomalies, is one of the main achievements of the formalism. For details see AL:review .
Due to lack of space, we will refrain from discussing in detail two of the main applications of LQG, namely Black Holes and Cosmology, so we will only give a brief summary. For black holes, a detailed treatment of the boundary conditions of the theory in the presence of an isolated horizon has lead to a detailed description of the quantum geometry associated with the horizon, to an identification of the quantum degrees of freedom responsible for the entropy, and to a counting of states that yields the entropy proportional to the area. The proportionality factor is set, for large black holes, to $`1/4`$ by adjusting the only free parameter of the theory, that is, the Barbero-Immirzi parameter. The same value can account for the entropy of very general BH, including Einstein-Maxwell, Einstein-Dilaton and non-minimally coupled scalar fields. For details and references see AL:review .
In the case of cosmology, a detailed program known as Loop Quantum Cosmology (LQC), has been developed in the past years (see for instance Bojowald’s contribution to this volume and bojo ). It represents a dimensional reduction, at the classical level, of the gravitational degrees of freedom to instances of great symmetry, as is expected to be the case in the cosmological setting. The quantization is performed via non-standard representations of the classical observable algebra, much in the spirit of LQG, with certain input coming from the full theory. As a result, the theory possesses very different behavior as the standard mini-superspace quantization: the space-time curvature has been shown to be bounded above and thus, the classical singularity is avoided. Furthermore, there is some evidence that the solutions may exhibit an inflationary period.
More recently, a new physical picture for the evaporation of black holes has emerged both from a detailed treatment of dynamical horizons and the end result of gravitational collapse (borrowing from results in LQC). These new developments suggests that one should reconsider the traditional viewpoint towards Hawking evaporation and the information loss problem AA:MB .
Loop quantum gravity is a theory in the making. There are still unsolved issues and things that remain to be done. The most notorious issue that has eluded a complete solution is the implementation of the Hamiltonian constraint. There is no proposal that is fully satisfactory. There is a consistent, free of anomalies quantization by Thiemann, but is far from unique and there is no control (yet) on the physical properties of its solutions (see Thiemann:2001yy ). More recently Thiemann has proposed the ‘master constraint program’ to implement the constraints. There is a completely different approach to the subject, where a projector onto physical states is sought by summing over 4-dim discrete structures known as spin-foams (see for instance Perez:2004hj ). Finally, a new program that tries to solve the constraint in a similar fashion as the diffeomorphism invariance is implemented, namely by considering finite actions of the constraint rather than defining and implementing its generator, is being pursued nowadays.
Another open issue is the semi-classical limit of the theory. That is, one needs to find states within the theory that approximate, in a low-energy/macroscopic limit, classical configurations such as Minkowski spacetime. To deal with this challenge there have been several proposals for semi-classical states, from the so-called weave states, to coherent states Thiemann:2001yy , WKB-like states Smolin:2004sx , and more recently Fock-motivated states AL:review . In most of these cases, the states belong to the kinematical space of states. A systematic approach to the definition of semiclassical states for constrained systems and to the question of whether one can approximate physical semi-classical states with kinematical ones has only begun ABC .
In previous sections we have defined LQG in terms of what it is and does. To end this section we will part from this tradition and make some remarks about what LQG is not. This can be seen as a complement to the FAQ section of Smolin:2004sx .
* Renormalizability. It is sometimes asserted that since LQG is a reformulation of General Relativity, then it is non-renormalizable. The first obvious observation is that LQG is a non-perturbative quantum theory, and as such, it has been shown to be finite (see for instance AL:review ; Thiemann:2001yy ). There is no inconsistency with the fact that Einstein gravity is perturbatively non-renormalizable. The level of rigor with which LQG is defined, as a mathematically well defined theory, is superior to any interacting quantum field theory in 4 dimensions. As such, the theory has also shown to be quite independent of “infrared” cutoffs.
* What is LQG not? One sometimes sees statements like: “various alternative formulations of LQG such as random lattice or spin foam theory…” This claim is incorrect. Spin foams and random lattices are not alternative formulations of LQG. In particular, random lattices are not directly related to LQG, and spin foams are tools (under development) to compute projection operators onto physical states of the theory, i.e. as a tool for solving the quantum scalar constraint of LQG (see Perez:2004hj for details).
* Is LQG a discrete theory? In this regard, it is important to stress, again, that LQG is a well defined theory on the continuum. It is not a discretized gravity model such as Regge calculus or dynamical triangulations, where the discrete structure is assumed from the beginning. It is true that a convenient basis of states in the Hilbert space of LQG is given by graphs and spins (the spin networks basis), but the theory is defined in the continuum, and the quantized character of the geometry (the discrete eigenvalues for the operators) are found as predictions of the theory. There is no continuum limit to be taken. All details and mathematical proofs of these assertions can be found in Refs. AL:review ; Thiemann:2001yy .
To end this primer, let us point out to further references. Once the reader is somewhat familiar with the basic assumptions of the theory, she should move on to move detailed description of the formalism. In this regard, the book by Rovelli Rovelli:2004tv is a good read for the conceptual background and spirit of LQG. For a good degree of detail, rigor and current status of the program see AL:review . For a completely detailed treatment of the formalism see the monograph by Thiemann Thiemann:2001yy .
## Acknowledgments
I would like to thank the organizers of the VI Mexican School for the invitation to write this contribution. This work was in part supported by grants DGAPA-UNAM IN-108103 and CONACyT 36581-E. |
warning/0507/astro-ph0507218.html | ar5iv | text | # Fitting Photometry of Blended Microlensing Events
## 1 Introduction
Gravitational microlensing has become a useful tool in measuring the amount of matter along the line-of-sight to distant stars. Since gravitational lensing depends only on mass, microlensing is sensitive to all compact forms of matter independent of their luminosity. Thus measurements out of the plane of the Galaxy towards the LMC, SMC, and M31 have given important limits/detections of dark matter in the halo ( Aubourg, et al. 1993; Lasserre et al. 2000; Alcock, et al. 1993; 1997a; 2000; Paulin-Henriksson et al. 2003; de Jong et al. 2004), and measurements towards the Galactic bulge give important constraints on the mass and distribution of Galactic stars, including those too faint to observe directly (Griest et al. 1991; Paczyński 1991; Han & Gould 2003; Udalski, et al. 1994; Alcock, et al. 1997b; Afonso, et al. 2003; Popowski, et al. 2005; Sumi, et al. 2005).
The signal of microlensing is a specific transient magnification of a background source star as the lens object passes in front of it, and thus microlensing experiments repeatedly monitor many ordinary stars to find microlensing lightcurves. The probability of microlensing occurring to a given star is called the optical depth, $`\tau `$ and is of order $`10^6`$ or less for many Milky Way lines-of-sight. The smallness of $`\tau `$ means that microlensing experiments concentrate on very crowded star fields where many hundreds of thousands of stars can be simultaneously imaged. This allows many lightcurves to be created simultaneously but also results in blending of the source stars together. This blending causes two problems in using the detected microlensing events to infer the optical depth. First, since each “source object” may contain the light from many stars, the number of stars being monitored is not just the number of objects being photometered. Second, the magnification profile of a microlensing event is changed when unlensed light is blended with the lensed light of the source star.
In this paper we revisit the problem of blending in microlensing lightcurves. There are several methods of dealing with the blending problem. Among these are 1) obtaining high resolution images from space, which will usually allow separationof the source object into its different components, giving a direct measurement of the fraction of light from the lensed source (Alcock, et al. 2001), 2) if the unlensed light is not exactly centered on the lensed source, then the centroid of the light will shift during the microlensing event allowing limits on blending to be placed (Alard, Mao, & Guibert 1995), 3) if the lensed source is a different color than the unlensed light then a color shift will occur as the event proceeds, allowing limits on lensed-light fraction to be made (Alard, Mao, & Guibert 1995), and 4) for image subtraction lightcurves, the source can in principle be removed and this can help break the degeneracy in some cases (Gould & An 2002).
However, we will not discuss the above methods in this paper but will focus on the fitting and interpretation of the photometric data alone; that is, we include the lensed-light fraction as a parameter in the microlensing fit and hope to use the shape of the lightcurve to recover this information. In principle this allows recovery of the actual event duration, and a measurement of the amount of blending in the sample of events, allowing corrections to be made in estimating the optical depth. A related and popular method is to calculate lensing optical depth using only a subsample of very bright source stars (e.g. clump giants). The idea is that very bright source stars are less likely to be blended, and when they are blended, should be blended only by a small amount. In this case one would like to use the blend fits only to determine whether or not a given event is blended.
Unfortunately, as pointed out previously (e.g. Han 1999; DiStefano & Esin 1995; Woźniak & Paczyński 1997; Alard 1997, etc.) blended fits tend to be quite degenerate. A lightcurve with a small lensed-light fraction looks very much like an unblended lightcurve with a smaller maximum magnification and a smaller event duration. As pointed out previously, this means that this fitting method will be of limited use in many cases. Our study adds strength to the conclusions of previous workers, points out several new problems with blend fits, and makes recommendations on how best to proceed with blend fits for those who choose to do them. We will discuss what happens when the microlensing event contains signal from other physical effects such as weak parallax or binary effects. These effects are not rare, and since the difference between blended and unblended lightcurves is small, even an almost undetectable real deviation from the standard point-source-point-lens lightcurve can render blend fit results meaningless.
The plan of the paper is as follows: In § 2 we define our notation and discuss the similarities and differences between unblended microlensing and blended microlensing We also give an analytic approximation that gives the underlying event duration and peak magnification from the lensed-light fraction and the easily measured apparent event duration and maximum magnification.
In § 3 we discuss the usefulness of blend fits and compare with earlier work. In § 4 we discuss the optimal times to take follow-up data in order to improve recovery of parameters from the blend fit. In § 5 we discuss the problem of the baseline magnitude, and In § 6 we discuss the problem of non-Gaussian data and whether the errors returned by fitting programs are reliable.
## 2 Degeneracies in blended lightcurves: analytic approximations for event duration and $`A_{\mathrm{max}}`$
When an isolated lens object crosses close to the line-of-sight of an isolated background source star, the source is magnified and a microlensing lightcurve is generated with magnification
$$A(u)=\frac{u^2+2}{u(u^2+4)^{1/2}},u^2(t)=u_{\mathrm{min}}^2+\left(\frac{tt_{\mathrm{max}}}{t_E}\right)^2,$$
(1)
where $`u`$ is the projected distance between the lens and source in units of the Einstein ring radius, $`t_E`$ is the time to cross the Einstein radius, and $`t`$ is time, with maximum magnification, $`A_{\mathrm{max}}`$, occurring at $`t_{\mathrm{max}}`$.
The most important parameter is the event duration $`t_E`$ since the optical depth depends upon the sum of efficiency weighted event durations:
$$\tau =\frac{\pi }{2E}\underset{\mathrm{events}}{}\frac{t_E}{ϵ(t_E)},$$
(2)
where the exposure $`E`$ is the product of the length in days of the observing program and the number of observed stars, and $`ϵ`$ is the efficiency of detecting an event of duration $`t_E`$.
When other sources of light are contained in the same seeing element as the lensed source star, the microlensing lightcurve is altered since only a fraction of the light is actually lensed:
$$A^{}(u)=f_{\mathrm{ll}}A(u)f_{\mathrm{ll}}+1,$$
(3)
where $`f_{\mathrm{ll}}`$ is the fraction of light that is lensed (a.k.a. the blend fraction, i.e. coming from the source star) before the lensing event begins.<sup>1</sup><sup>1</sup>1Several terms have been used in different ways in the literature for blend fraction, most commonly $`f_b`$ which either means the fraction of light coming from the lens or the fraction of the light coming from non-lens sources. We introduce the new symbol $`f_{\mathrm{ll}}`$ to avoid the extant confusion of nomenclature. Compared with unblended events, blended photometric microlensing lightcurves suffer from a smaller maximum magnification, $`A_{\mathrm{max}}^{}`$ and shorter event duration $`t_E^{}`$, as well as from potential color shifts if the blended light has a different spectrum.
In fact, a blended lightcurve looks remarkably like an unblended lightcurve with different values of $`t_E`$ and $`A_{\mathrm{max}}`$ (e.g. Han 1999; DiStefano & Esin 1995; Woźniak & Paczyński 1997; Alard 1997). However, as illustrated in Figure 1, this similarity is not perfect and there are differences in the shapes of blended and unblended lightcurves. It is these differences which give rise to the hope that information about blending can be extracted by fitting lightcurves with blending parameters. If this similarity were perfect, then there would be no use in fitting blended lightcurves to photometric data. Figure 1 shows that the shape differences are typically small, meaning that extracting blending information will be difficult. Woźniak and Paczyński 1997 (WP) studied this in detail and gave regions of the $`f_{\mathrm{ll}}`$, $`A_{\mathrm{max}}`$ plane where blended fits were useful and where they were not. We return to this subject in § 3, but the qualitative results of WP can be seen from Figure 1, where low magnification events show a maximum difference between the blended lightcurve and the best fit unblended lightcurve of only 1% or so, while the higher magnification events show more substantial differences.
Previous workers have also given analytic formulas relating the measured (apparent) maximum magnification $`A_{\mathrm{max}}^{}`$, and the apparent event duration $`t_E^{}`$ to $`f_{\mathrm{ll}}`$, $`A_{\mathrm{max}}`$, and $`t_E`$. For example, Woźniak and Paczyński (WP) studied the degeneracies by performing expansions of the above equations in the limits of small $`u_{\mathrm{min}}`$ and large $`u_{\mathrm{min}}`$ and in these limits give the formulas relating the actual values of $`t_E`$ and $`A_{\mathrm{max}}`$ to $`f_{\mathrm{ll}}`$ and the measured $`A_{\mathrm{max}}^{}`$ and $`t_E^{}`$, For small $`u_{\mathrm{min}}`$ (large $`A_{\mathrm{max}}`$) they found $`u_{\mathrm{min}}^{}u_{\mathrm{min}}/f_{\mathrm{ll}}`$, and $`t_E^{}f_{\mathrm{ll}}t_E`$, while in the limit of large $`u_{\mathrm{min}}`$ ($`A_{\mathrm{max}}1`$) they found $`u_{\mathrm{min}}^{}u_{\mathrm{min}}/f_{\mathrm{ll}}^{1/4}`$, and $`t_E^{}f_{\mathrm{ll}}^{1/4}t_E`$.
DiStefano & Esin (1995), and Han (1999), and Alard (1997) took a different approach, solving equation 3 for $`A_{\mathrm{max}}`$ and giving the actual $`t_E`$ in terms of $`t_E^{}`$ by requiring that the two different parameterizations give the same amount of time with $`A>1.3416`$. They found:
$$A_{\mathrm{max}}^{}{}_{(\mathrm{HDE})}{}^{}=(A_{\mathrm{max}}^{}1+f_{\mathrm{ll}})/f_{\mathrm{ll}},t_{E}^{}{}_{(\mathrm{HDE})}{}^{}=t_E^{}\left(\frac{u_1^2u_{\mathrm{min}}^2}{1u_{\mathrm{min}}^{}{}_{}{}^{2}}\right)^{1/2},$$
(4)
where $`u_{\mathrm{min}}^{}=u(A_{\mathrm{max}}^{})`$, and $`u_1=u(A(A^{}=1.3416))`$, can be found from the inverse of equation 1:
$$u(A)=(2/\sqrt{11/A^2}2)^{1/2}$$
(5)
Noting in Figure 1 that the differences between the blended and unblended lightcurves tend to be large in the peak, and that the values of $`t_E^{}`$ and $`A_{\mathrm{max}}^{}`$ are found by fitting, we worried that the HDE formula, which assumes equality in the peak, might not be accurate. We also wondered about the range of applicability of the WP formulas and so decided to test these formulas. We did this by fitting artificial blended lightcurves with an unblended source model and finding the best fit values of $`t_E^{}`$ and $`A_{\mathrm{max}}^{}`$. We also fit these lightcurves with blended source models and correctly extracted the input blend parameters.
As shown in Figure 2, we found that the WP formulas are not very useful over most of the parameter range, and that the HDE equations work well only over a restricted range of parameters. For the WP formulas this is not surprising since they were created only to show that the degeneracies exist in certain limits. For relatively large lensed-light fraction and for relatively low values of $`A_{\mathrm{max}}`$ the HDE equations give a good estimation of the best fit $`A_{\mathrm{max}}^{}`$ and $`t_E^{}`$, but for small lensed-light fraction or high $`A_{\mathrm{max}}`$ the estimates of these equation can be far off.
As expected, it is just where the blended lightcurve shape differs the most from an unblended fit that the HDE approximations do not work well. The reason can be seen in Figure 1, where for high magnification events and low lensed-light fraction the blended lightcurve differs strongly in the peak area, but not so much in the lightcurve middle rising and falling regions. Thus, the best fit unblended lightcurve will allow the actual peak magnification to overshoot and compensate for these points by undershooting in the middle regions. Since the HDE formula forces the lightcurves to match at the peak and when $`A_{\mathrm{max}}^{}=1.34`$, it will overestimate the best fit peak magnitude and underestimate the event duration.
By studying many such examples, one can come up with a formula that does a better job of relating the best fit $`A_{\mathrm{max}}^{}`$ and $`t_E^{}`$ to $`A_{\mathrm{max}}`$, $`t_E`$, and $`f_{\mathrm{ll}}`$ in the parameter ranges where the HDE formula does not work well. The points in Figure 2 show the best fit values of $`A_{\mathrm{max}}^{}`$ and $`t_E^{}`$ vs $`f_{\mathrm{ll}}`$ found by fitting artificial blended lightcurves. The dashed lines show the HDE estimates and long dash lines the WP estimates for $`t_E^{}/t_E`$. At small values of $`A_{\mathrm{max}}`$ ($`<3`$) the HDE formulas do work very well (better than the new formula) and they should be used. However for $`A_{\mathrm{max}}>3`$ the HDE formulas do not give accurate estimates. To find a better approximation, one can fit a straight line to the data for a given $`u_{\mathrm{min}}`$ and get a formula which fits well except for very low lensed-light fraction. Repeating this procedure for different values of $`u_{\mathrm{min}}`$, one discovers that the slopes and zero points of the linear fits are also quite linear in $`u_{\mathrm{min}}`$. Thus a simple fitting linear formula that covers much of the parameter space can be found. However, if one fits a quadratic for the low $`f_{\mathrm{ll}}`$ events one can get an even better formula which works very well for $`f_{\mathrm{ll}}<0.3`$. Thus we find an approximation:
$`A_{\mathrm{max}}`$ $``$ $`\{\begin{array}{cc}A_{\mathrm{HDE}},\hfill & \text{if}A_{\mathrm{max}}<3;\hfill \\ \frac{A_{\mathrm{max}}^{}0.9785+0.4150f_{\mathrm{ll}}}{0.8153f_{\mathrm{ll}}+0.00021},\hfill & \text{if}A_{\mathrm{max}}>3\mathrm{and}A_{\mathrm{max}}^{}<10;\hfill \\ \frac{A_{\mathrm{max}}^{}0.3618+0.2106f_{\mathrm{ll}}}{1.0282f_{\mathrm{ll}}0.04433},\hfill & \text{if}A_{\mathrm{max}}>3\mathrm{and}A_{\mathrm{max}}^{}>10,;\hfill \end{array}`$
$`t_E^{}/t_E`$ $``$ $`\{\begin{array}{cc}t_{E}^{}{}_{(\mathrm{HDE})}{}^{}/t_E,\hfill & \mathrm{if}A_{\mathrm{max}}<3;\hfill \\ (1.0946u_{\mathrm{min}}+0.9418)f_{\mathrm{ll}}+1.141u_{\mathrm{min}}+0.0564,\hfill & \mathrm{if}A_{\mathrm{max}}>3\mathrm{and}f_{\mathrm{ll}}>0.3;\hfill \\ 𝑭𝑪𝑼,\hfill & \mathrm{if}A_{\mathrm{max}}>3\mathrm{and}f_{\mathrm{ll}}<0.3,\hfill \end{array}`$ (6)
where
$`𝑭=(1,f_{\mathrm{ll}},f_{\mathrm{ll}}^2),𝑼=\left(\begin{array}{c}1\\ u_{min}\\ u_{min}^2\end{array}\right),`$ (7)
$$𝑪=\left(\begin{array}{ccc}0.02548& 1.0626& 1.6504\\ 1.1914& 7.284& 11.50\\ 0.8824& 26.58& 44.28\end{array}\right),$$
(8)
and $`u_{\mathrm{min}}`$ is found from $`A_{\mathrm{max}}`$ and equation 5.
In using this formula, one typically starts with measured values of $`t_E^{}`$, $`A_{\mathrm{max}}^{}`$, and an initial guess of $`A_{\mathrm{max}}`$ and uses different (unknown) values of $`f_{\mathrm{ll}}`$, to find the corresponding underlying $`A_{\mathrm{max}}`$ and $`t_E`$. If the value of $`A_{\mathrm{max}}`$ found using the new fitting formula is smaller than 3, then one should use the HDE formula instead.
The new fitting formula is shown as the solid line in Figure 2 and does better than HDE or WP for $`A_{\mathrm{max}}>3`$. Over the range $`0.01<f_{\mathrm{ll}}<1.1`$, and $`3<A_{\mathrm{max}}<70`$ the new fit formula gives a typical error in $`t_E`$ (compared with actually fitting the microlensing lightcurve with a blend fit model) of around 3% and a maximum error of 9%. For $`A_{\mathrm{max}}`$ the typical error is 4% and the maximum error is 12%. The HDE formula can be off by more than 50% in $`t_E`$ and 24% in $`A_{\mathrm{max}}`$ in this region of parameter space.
In summary, we tested the HDE formula, Woźniak and Paczyński (WP) formulas and Equation 6 over a wide range of parameters and found the new fitting function works better than HDE for all values of $`f_{\mathrm{ll}}`$ when $`A_{\mathrm{max}}>3`$ and $`A_{\mathrm{max}}^{}>1.34`$, while the old HDE formula works better for low values of $`A_{\mathrm{max}}`$ and $`A_{\mathrm{max}}^{}`$. The WP large $`A_{\mathrm{max}}`$ formula gives $`t_E`$ within 10% only for large $`f_{\mathrm{ll}}(>0.5`$), and large $`A_{\mathrm{max}}`$, while the other WP formula is not useful except for $`A_{\mathrm{max}}1.34`$.
Since in microlensing experiments the event durations are found by photometric fitting and since the optical depth is proportional to the sum of the fit $`t_E`$’s, when making corrections for blending it is important to properly relate the lens-light fraction of each event to the underlying event duration.
## 3 Usefulness of blend fits
Woźniak and Paczyński (1997) (WP) studied the degeneracy of blend fits and concluded that in many cases blended and unblended lightcurves cannot be distinguished by photometric fitting. They described areas of parameter space where blend fits would be useful and areas where they would not. While we think that WP did an accurate and very useful calculation, and we agree with their conclusion that blend fits are usually not very useful, we wanted to repeat their analysis for several reasons. First, WP did not include the baseline magnitude in their fits, reasoning that since many measurements are taken before and after the event, the error in baseline magnitude was not significant. In fact, we find that error in the baseline magnitude is one of the most severe problems in blend fits. We find that errors even at the few percent level can drastically alter the parameter values extracted from the fit. Second, WP considered only evenly spaced observations and we wanted to consider whether different follow-up strategies could improve the ability to extract the parameters.
In our studies, we find the error in fit parameters three ways. First we create artificial lightcurves using the theoretical formula and add Gaussian random noise to each measurement. We perform blended and unblended fits on these lightcurves using Minuit (CERN Lib. 2003). Second we calculate the error matrix by inverting the Hessian matrix as discussed in Gould (2003). Finally to understand the effect of the non-Gaussianity of the errors in real microlensing experiments we create artificial lensing lightcurves by adding microlensing signal into actual non-microlensing lightcurves obtained by the MACHO collaboration, and then fit these.
Since the method of calculating the error matrix is closest to what WP did, we first give these results. Briefly, we calculate the Hessian matrix (the matrix of second derivatives of the light curve residuals with respect to each parameter) then invert it. The square root of the diagonal elements of the resulting matrix are then the one sigma errorbars of the parameters. This accounts for correlations in the parameters, but not any nonlinearities. WP used a very similar method, but used it to calculate the $`\mathrm{\Delta }\chi ^2`$ instead of the error bars. In figure 3 we show that our method brackets WP’s. We show limits calculated as both the one sigma lower limit on $`f_{ll}`$ for an unblended lightcurve and the value of $`f_{ll}`$ that gives $`f_{ll}+\sigma _{f_{ll}}=1`$.
We note, as WP found, that parameter errors scale linearly with
$$a=\frac{\sigma }{\sqrt{(}N)}$$
(9)
for $`N`$ points taken during the peak (defined as lasting $`4t_E`$)<sup>2</sup><sup>2</sup>2This is true for large enough value of $`N`$, for small values of $`N16`$ parameter errors increase faster than $`a`$.. Thus our results can be scaled for other numbers of observations with different values of $`\sigma `$. Thus, we find that our results agree with those of WP if we assume the baseline magnitude is known and take a uniform sampling.
## 4 Follow-up observations
Figure 1 shows that the difference between blended and unblended lightcurves is not always uniform across the lightcurve. So if one wanted to plan follow-up observations to improve the accuracy of the blend fit, one should concentrate on the regions of the lightcurve where the differences are largest. Thus it may be possible to do better than WP suggested with their equal spaced observation calculations. To test this hypothesis we calculated the error matrix for blended fits adding in follow-up observations at different points on the lightcurve. As seen in the Figure 1 examples, for any choice of parameters there are five places where the difference lightcurves are maximum, and therefore where follow-up data is more useful than average: at the peak, in the rising/falling portion of the curve, and in the wings. The precise locations change with the choice of parameters but for Figure 1a they are found to be localized near the peak at $`(|t/t_E|<0.1)`$, in the falling (or rising region) at $`(0.3<|t/t_E|<0.6)`$, and near the baseline at $`(1.0<|t/t_E|<1.5)`$. Observations taken between these regions do little in constraining the parameters. In addition points $`t`$ greater than $`2t_E`$ are very helpful because they fix the baseline in our simulated lightcurve. We discuss the baseline separately in § 5.
In Figure 4 we compare the relative value of added points as a function of the time they are added. We find that, in this case, with 40 observations, 4 extra focused observations can reduce the error on $`f_{ll}`$ by 7.7%. To get the same reduction of error on $`f_{ll}`$ with evenly distributed observations we would need 7 observations, in other words, each added focused observation is equivalent to increasing the sampling by 1.75 points over 4 $`t_E`$. The numerical value of the extra effectiveness obtained using focused versus evenly spaced observations varies with underlying parameters and the total number of added points.
Precise follow-up measurements at multiple focused locations can improve the determination of $`f_{ll}`$ even more as they further constrain the shape of the lightcurve.
In order to see the effect of adding multiple follow-up observations at two distinct times we compare this with adding evenly spaced observations. In Figure 5 we plot the increase (or decrease) in effectiveness of extra focused observations as a function of the two times at which they are taken. The contours around the light areas show regions of increased effectiveness, while dark areas show areas where the focused observations are less valuable than evenly spaced observations. In this example the effectiveness is increased by up to a factor two.
It is important to note that with more observations or higher accuracy in each follow-up region the advantage per added observation is reduced and the relative values of the various minimums vary, though they stay in roughly the same place. For practical use it is important to note that the time of the optimum second follow-up observation(s) varies with the time of the first follow-up observations(s). In practice one would need to calculate optimum observing times for an event in progress as a function of all the previous measurements.
One problem with the above approach is that without knowing the underlying parameters, particularly $`t_E`$, it is difficult to predict the best times to take follow-up data. To test if a practical experiment could be designed to take advantage of focused follow-up data we simulated an experiment. First we generated lightcurves with 80 points over $`8t_E`$ with .05 Gaussian errors at the baseline drawing $`f_{ll}`$ randomly from the interval $`[.01,1)`$ and $`u_{min}`$ randomly from the interval $`(0,1)`$ requiring $`A_{max}^{}>1.34`$ in the HDE approximation. We also adjusted $`t_E`$ to keep $`t_E^{}10\mathrm{d}\mathrm{a}\mathrm{y}\mathrm{s}`$ also using the HDE approximation. We then generated 9 follow-up observations over 3 days at the peak and fit the first half of the light curve plus the follow-up data. From this first fit we calculated the optimum times for two more bouts of follow-up. We generated these, both with 9 observations over 3 days, and then fit the entire lightcurve with the added 27 points. We also generated 27 points of follow-up uniformly distributed over the 20 days starting at the peak, added it to the initial lightcurve and fit the resulting data. To see the relative improvement for the two methods we calculate a parameter $`\zeta =(f_{ll_{focused}}^{}f_{ll})/(f_{ll_{unfocused}}^{}f_{ll})`$, which is the ratio of the error in blend fraction given by focused observations to the error in blend fraction given by uniform follow-up sampling. We plot the distribution of $`\zeta `$ in Figure 6 finding that our strategy gives an improvement $`(\zeta <1)`$ for 71% of the events and a worsening in 29% of the events. We find a substantial improvement ($`\zeta <.5`$) for 45% of the events, and and even larger improvement ($`\zeta <0.1`$) 18% of the time. Thus we conclude that for the the same amount of observing time we can make a more accurate measurement of $`f_{ll}`$ by focussing the follow-up observations.
In summary we find that observations concentrated at a few times can constrain the microlensing parameters as well as many measurements distributed throughout an event. The best place for these measurements are at the peak, in the falling/rising portion, and in the wings with regions between where added observations do no good. In most cases it is possible to constrain the event parameters well enough with the first half of the data and some follow-up observations near the peak to predict the last two optimum observing times.
## 5 Baseline Magnitude
The baseline magnitude of a lightcurve can in principle be very well determined since many measurements can be taken before or after the microlensing event. WP assumed that this was the case and so did not include the baseline magnitude as one of their fit parameters. In real microlensing surveys, however, it may be that the error in average magnitude is not entirely statistical, and may not average down as expected. There may be a systematic limit to the accuracy with which the baseline magnitude can be determined. In fact, detectors and telescope systems drift over time and so measurements made much later may actually reduce the accuracy of the baseline magnitude. To investigate the importance of the baseline magnitude, we created artificial lightcurves without any errors and fit them with a model with a fixed value of baseline magnitude that differed from the actual baseline magnitude by various amounts. Our results are shown in Figure 7. We find that the dependence on baseline is very strong for low amplification events and not as strong for higher amplifications events, but in any case even a 2% error in baseline magnitude determination can strongly bias the recovered blend fraction.
Next, to see how well baseline magnitudes converge in real data, we used the MACHO collaboration database of random stars (Alcock, et al. 2000). We looked at the $`\chi ^2/n_{dof}`$ of a fit to a constant lightcurve for our real lightcurves and compared this to simulated ideal lightcurves with the same number of points and Gaussian errors. For the simulated Gaussian lightcurves we find the $`\chi ^2/n_{dof}`$ distribution peaked near unity and distributed as expected, but for the real data the distribution of $`\chi ^2/n_{dof}`$ is much broader. These two distributions are shown in Figure 8.
As an estimate of the error in the baseline which arises due to the systematic drift and non-Gaussian nature of the magnitude errors, we calculated mean and median for the points in each of 146 lightcurves in MACHO field 119, one of the most frequently observed fields. We found that the distribution of mean minus median had a dispersion of 1.3% indicating that the error in the baseline flux is $`1.3\%`$. Referring to Figure 7 we see that for a typical event with $`u_{min}=0.5`$, this implies a typical spread in $`f_{ll}`$ of 0.18 due to baseline alone. Since half of all events have $`u_{min}<0.5`$, half of all events will have an even larger bias. For more sparsely sampled fields this dispersion due to error in baseline fit would be even larger.
## 6 Errors in Fit Parameters
From Macho Project data (Table 6 of Popowski et al. 2005) it seems blend fits return biased parameters. For the set of Macho clump giant events, which are believed to be minimally blended from their positions on the color magnitude diagram, many are best fit with blending. If the events are not blended then a systematic bias in the fits must make them appear to be blended. A systematic bias in recovered lensed-light fraction would lead to a bias in the optical depth as well. The MACHO collaboration investigated the blending of their clump giant sample and decided to use the parameters from the unblended fits. They also used a subsample of events that were less likely to be blended to check for a bias due to blending and found no such bias.
To test for a systematic bias we generate 1000 lightcurves with Gaussian errors for each of three different values for the error on each point: $`\sigma =0.01`$, $`\sigma =0.05`$, and $`\sigma =0.15`$. The recovered lensed-light fractions for these events are shown in Figure 9. As the error on individual datum increases the distribution of $`f_{ll}^{}`$ becomes increasingly skewed. We find that while the mean $`f_{ll}^{}`$ may not decrease, the most probable value does decrease. This reduction in the mode is at least partially compensated by the large tail of the distribution with $`f_{ll}^{}>1`$, but for the small number of events a microlensing experiment observes it is unlikely that many of the few events with $`f_{ll}^{}1`$ will be observed. Even if one event with $`f_{ll}^{}1`$ is observed it may be ignored as it is an unphysical value of the parameter, thus leading to an underestimate of the average value of $`f_{ll}`$. Thus we find that as the errors in measurement increase blend fitting becomes more and more likely to return biased results. The direction of the bias is more often toward small values of $`f_{ll}`$. Thus events that are in reality unblended become more and more likely to return fit values implying that they are heavily blended.
## 7 Conclusions
We find agreement with previous workers that blend fits are problematic, but can be useful especially for high magnification events. When performing blend fits it is helpful to get extra measurements near the peak and at other specific points along the lightcurve. We find that if care is not taken in the treatment of the lightcurve baseline magnitude the fit results can be severely biased and in real data the errors returned on fit parameters should be treated with caution. We find that blend fits return a biased, skewed distribution of the underlying parameters tending to indicate more blending than actually exists. Finally, note that when the microlensing event contains signal from other physical effects such as weak parallax or binary effects blend fits can yield unreliable results. These effects are not rare, and since the difference between blended and unblended lightcurves is small, even an almost undetectable real deviation from the standard point-source-point-lens lightcurve can render blend fit results meaningless.
We thank David P. Bennett and Piotr Popowski for many useful discussions on the topic of blending. This work is supported in part by the DoE under grant DEFG0390ER40546. |
warning/0507/math0507600.html | ar5iv | text | # 𝐂̄_{𝐧,𝐧-𝟏} revisited
## 1. Introduction
In spite of its importance, the proof of $`\overline{C}_{n,n1}`$ is not so easy to access for the younger generation, including myself. After \[Ka1\] was published, the birational geometry has drastically developed. When Kawamata wrote \[Ka1\], the following techniques and results are not known nor fully matured.
* Kawamata’s covering trick,
* moduli theory of curves, especially, the notion of level structures and the existence of tautological families,
* various notions of singularities such as rational singularities, canonical singularities, and so on.
See \[Ka2, §2\], \[AK, Section 5\], \[AO, Part II\], \[vGO\], \[V2\], and \[KM\]. In the mid 1990s, de Jong gave us fantastic results: \[dJ1\] and \[dJ2\]. The alteration paradigm generated the weak semistable reduction theorem \[AK\]. This paper shows how to simplify the proof of the main theorem of \[Ka1\] by using the weak semistable reduction. The proof may look much simpler than Kawamata’s original proof (note that we have to read \[V1\] and \[V2\] to understand \[Ka1\]). However, the alteration theorem grew out from the deep investigation of the moduli space of stable pointed curves (see \[dJ1\] and \[dJ2\]). So, don’t misunderstand the real value of this paper. We note that we do not enforce Kawamata’s arguments. We only recover his main result. Of course, this paper is not self-contained.
The following result is the main theorem of \[Ka1\]. We call this $`\overline{C}_{n,n1}`$ in this paper. Here, $`n`$ means the dimension of $`X`$.
###### Theorem 1.1 (\[Ka1, Theorem 1\]).
Let $`f:XY`$ be a dominant morphism of algebraic varieties defined over the complex number field $``$. Assume that the general fiber $`X_y=f^1(y)`$ is an irreducible curve. Then we have the following inequality for logarithmic Kodaira dimensions$`:`$
$$\overline{\kappa }(X)\overline{\kappa }(Y)+\overline{\kappa }(X_y).$$
In Section 2, we will give a proof to \[Ka1, Theorem 2\], which is stronger than $`\overline{C}_{n,n1}`$. See the inequality $`(\overline{C}_{n,n1}^{})`$ in the first paragraph of the proof below.
Note that our reference list does not cover all the papers treating the related topics. We apologize in advance to the colleagues whose works were not appropriately mentioned in this paper. In the proof of the main theorem, we do not refer to the original results since they are scattered in various papers. Mori collected them nicely in \[M, §6, §7\].
###### Acknowledgments.
I was grateful to the Institute for Advanced Study for its hospitality. I was partially supported by a grant from the National Science Foundation: DMS-0111298. I would like to thank Professor Noboru Nakayama for comments and Professor Kalle Karu for giving me \[Kr\].
###### Notation.
We will work over $``$ throughout this paper. For the basic properties of the logarithmic Kodaira dimension, see \[I1\], \[I2\], \[I3\], and \[Ka1, §1\].
1. Let $`X`$ be a (not necessarily complete) variety. Then $`\overline{\kappa }(X)`$ denotes the logarithmic Kodaira dimension of $`X`$.
2. Let $`f:XY`$ be a dominant morphism between varieties and $`D`$ a $``$-divisor on $`X`$. We can write $`D=D_{\mathrm{hor}}+D_{\mathrm{ver}}`$ such that every irreducible component of $`D_{\mathrm{hor}}`$ (resp. $`D_{\mathrm{ver}}`$) is mapped (resp. not mapped) onto $`Y`$. If $`D=D_{\mathrm{hor}}`$ (resp.$`D=D_{\mathrm{ver}}`$), $`D`$ is said to be horizontal (resp. vertical).
3. Let $`f:XY`$ be a birational morphism. Then $`\mathrm{Exc}(f)`$ denotes the exceptional locus of $`f`$.
## 2. $`\overline{C}_{n,n1}`$
Here, we prove the following theorem. It is easy to see that this statement is equivalent to Theorem 1.1 by the basic properties of the logarithmic Kodaira dimension.
###### Theorem 2.1 ($`\overline{C}_{n,n1}`$).
Let $`f:XY`$ be a surjective morphism with connected fibers between non-singular projective varieties $`X`$ and $`Y`$. Let $`C`$ and $`D`$ be simple normal crossing divisors on $`X`$ and $`Y`$. We put $`X_0:=XC`$ and $`Y_0:=YD`$. Assume that $`f(X_0)Y_0`$. Then
$$\overline{\kappa }(X_0)\overline{\kappa }(Y_0)+\overline{\kappa }(F_0),$$
where $`F_0`$ is a sufficiently general fiber of $`f_0:=f|_{X_0}:X_0Y_0`$.
Before we start the proof, let us recall the following trivial lemma. We will frequently use it without mentioning it.
###### Lemma 2.2.
Let $`X`$ be a complete normal variety. Let $`D_1`$ and $`D_2`$ be $``$-Cartier $``$-divisors on $`X`$. Assume that $`D_1D_2`$. Then $`\kappa (D_1)\kappa (D_2)`$.
###### Proof of Theorem 2.1.
By \[Ka1\], it is sufficient to prove
($`\overline{C}_{n.n1}^{}`$)
$$\kappa (K_X+Cf^{}(K_Y+D))\overline{\kappa }(F_0).$$
###### Step 1.
By Theorem 2.1 in \[AK\] (see also \[Kr, Chapter 2, Remark 4.5 and Section 9\]), we have the following commutative diagram:
$$\begin{array}{ccccc}X& & X^{}& & U_X^{}\\ & & & & \\ Y& & Y^{}& & U_Y^{}\end{array}$$
such that $`p:X^{}X`$ and $`q:Y^{}Y`$ are projective birational morphisms, $`X^{}`$ is quasi-smooth (in particular, $``$-factorial) and $`Y^{}`$ is non-singular, the inclusion on the right are toroidal embeddings, and such that
1. $`f^{}:(U_X^{}X^{})(U_Y^{}Y^{})`$ is toroidal and equi-dimensional,
2. Let $`C^{}:=(p^{}C)_{\mathrm{red}}`$ and $`D^{}:=(q^{}D)_{\mathrm{red}}`$. Then $`C^{}X^{}U_X^{}`$ and $`D^{}Y^{}U_Y^{}`$.
Since
$$\overline{\kappa }(X_0)=\kappa (K_X+C)=\kappa (K_X^{}+C^{})$$
and
$$\overline{\kappa }(Y_0)=\kappa (K_Y+D)=\kappa (K_Y^{}+D^{}),$$
we can replace $`f:XY`$ with $`f^{}:X^{}Y^{}`$. For the simplicity of the notation, we omit the superscript . So, we can assume that $`f:XY`$ is toroidal with the above extra assumptions.
###### Step 2.
By taking a Kawamata’s Kummer cover $`q:Y^{}Y`$, we obtain the following commutative diagram:
$$\begin{array}{ccc}X& \stackrel{p}{}& X^{}\\ f& & f^{}& & \\ Y& \underset{q}{}& Y^{}\end{array}$$
such that $`f^{}:X^{}Y^{}`$ is weakly semistable, where $`X^{}`$ is the normalization of $`X\times _YY^{}`$ (see \[AK, Section 5\]). Note that $`X^{}`$ is Gorenstein by \[AK, Lemma 6.1\]. We put $`G:=XU_X`$ and $`H:=YU_Y`$. Then we have
$$K_X+Cf^{}(K_Y+D)K_X+C_{\mathrm{hor}}+G_{\mathrm{ver}}f^{}(K_Y+H).$$
Therefore, we can check that
$$p^{}(K_X+Cf^{}(K_Y+D))K_{X^{}/Y^{}}+(p^{}C)_{\mathrm{hor}}.$$
We note that $`(p^{}C)_{\mathrm{hor}}=p^{}(C_{\mathrm{hor}})`$. So, it is sufficient to prove that $`\kappa (K_{X^{}/Y^{}}+(p^{}C)_{\mathrm{hor}})\overline{\kappa }(F_0)`$.
###### Step 3.
Let $`F`$ be a general fiber of $`f:XY`$. We put $`g:=g(F)`$: the genus of $`F`$.
###### Case ($`g2`$).
In this case,
$$\kappa (K_{X^{}/Y^{}}+(p^{}C)_{\mathrm{hor}})\kappa (K_{X^{}/Y^{}})1=\overline{\kappa }(F_0).$$
The last inequality is well-known. See, for example, \[M, (7.5) Theorem\] and \[F1, Theorem 5.3, Remark 5.4\]. So, we stop the proof in this case.
###### Case ($`g=1`$).
It is well-known that
$$\kappa (K_{X^{}/Y^{}})\mathrm{Var}(f^{})=\mathrm{Var}(f)0.$$
See \[M, (7.5) Theorem\] and \[F1, Theorem 5.3, Remark 5.4\]. For the definition of the variation $`\mathrm{Var}(f)`$, see, for instance, \[V3, p.329\] and \[M, (7.1)\]. So, if $`C`$ is vertical or $`\mathrm{Var}(f)1`$, then we obtain
$$\kappa (K_{X^{}/Y^{}}+(p^{}C)_{\mathrm{hor}})\overline{\kappa }(F_0).$$
Therefore, we can assume that $`\mathrm{Var}(f)=0`$ and $`C`$ is not vertical. By Kawamata’s covering trick, we obtain the following commutative diagram:
$$\begin{array}{ccc}X^{}& \stackrel{\pi }{}& X^{\prime \prime }\\ f^{}& & f^{\prime \prime }& & \\ Y^{}& \underset{\eta }{}& Y^{\prime \prime },\end{array}$$
where $`\eta :Y^{\prime \prime }Y^{}`$ is a finite cover from a non-singular projective variety $`Y^{\prime \prime }`$, $`f^{\prime \prime }:X^{\prime \prime }:=X^{}\times _Y^{}Y^{\prime \prime }Y^{\prime \prime }`$ is weakly semistable, and $`f^{\prime \prime }`$ is birationally equivalent to $`Y^{\prime \prime }\times EY^{\prime \prime }`$. Here, $`E`$ is an elliptic curve. Note that, if we need, we can blow-up $`Y^{}`$ and replace $`X^{}`$ with its base change before taking the cover. It is because the property of a morphism being weakly semistable is preserved by a base change under some mild conditions (cf. \[AK, Lemma 6.2\]). For details, see \[AK, Lemma 6.2\] and the proof of \[Ka2, Corollary 19\]. Since
$$\pi ^{}(K_{X^{}/Y^{}}+(p^{}C)_{\mathrm{hor}})=K_{X^{\prime \prime }/Y^{\prime \prime }}+\pi ^{}((p^{}C)_{\mathrm{hor}}),$$
it is sufficient to prove $`\kappa (K_{X^{\prime \prime }/Y^{\prime \prime }}+\pi ^{}((p^{}C)_{\mathrm{hor}}))1`$. Let $`\alpha :\stackrel{~}{X}Y^{\prime \prime }\times E`$, $`\beta :\stackrel{~}{X}X^{\prime \prime }`$ be a common resolution. Since $`X^{\prime \prime }`$ has only rational Gorenstein singularities, $`X^{\prime \prime }`$ has at worst canonical singularities. Thus, we obtain
$$\kappa (K_{X^{\prime \prime }/Y^{\prime \prime }}+\pi ^{}((p^{}C)_{\mathrm{hor}}))=\kappa (K_{\stackrel{~}{X}/Y^{\prime \prime }}+\beta ^{}\pi ^{}((p^{}C)_{\mathrm{hor}})).$$
On the other hand,
$$K_{\stackrel{~}{X}/Y^{\prime \prime }}=K_{\stackrel{~}{X}/Y^{\prime \prime }\times E}+K_{Y^{\prime \prime }\times E/Y^{\prime \prime }}=:A$$
is an effective $`\alpha `$-exceptional divisor such that $`\mathrm{Supp}A=\mathrm{Exc}(\alpha )`$. Let $`B`$ be an irreducible component of $`\beta ^{}\pi ^{}((p^{}C)_{\mathrm{hor}})`$ such that $`B`$ is dominant onto $`Y^{\prime \prime }`$. Then
$$m(A+\beta ^{}\pi ^{}((p^{}C)_{\mathrm{hor}}))\alpha ^{}\alpha _{}B,$$
for a sufficiently large integer $`m`$. Therefore, if is sufficient to prove $`\kappa (Y^{\prime \prime }\times E,\alpha _{}B)1`$. It is true by \[F2, Corollary 5.4\]. Thus, we finish the proof when $`g=1`$.
###### Case ($`g=0`$).
As in the above case, we can take a finite cover and obtain the following commutative diagram:
$$\begin{array}{ccc}X^{}& \stackrel{\pi }{}& X^{\prime \prime }\\ f^{}& & f^{\prime \prime }& & \\ Y^{}& \underset{\eta }{}& Y^{\prime \prime },\end{array}$$
where $`f^{\prime \prime }`$ is birationally equivalent to $`Y^{\prime \prime }\times ^1Y^{\prime \prime }`$. We can further assume that all the horizontal components of $`\pi ^{}((p^{}C)_{\mathrm{hor}})`$ are mapped onto $`Y^{\prime \prime }`$ birationally.
###### Lemma 2.3 (cf. \[F1, Section 7\]).
Let $`f:VW`$ be a surjective morphism between non-singular projective varieties with connected fibers. Assume that $`f`$ is birationally equivalent to $`W\times ^1W`$. Let $`\{C_k\}`$ be a set of distinct irreducible divisors such that $`f:C_kW`$ is birational for every $`k`$ $`(k3)`$. Then
$$\kappa (K_{V/W}+C_1+C_2)0$$
and
$$\kappa (K_{V/W}+C_1+C_2+C_3)1.$$
###### Proof.
By modifying $`V`$ and $`W`$ birationally (see also \[F1, Lemma 7.8\]) and replacing $`C_k`$ with its strict transform, we can assume that there exists a simple normal crossing divisor $`\mathrm{\Sigma }`$ on $`W`$ such that
$$\phi _{ij}:V_0:=f^1(W_0)W_0\times ^1$$
with $`\phi _{ij}(C_i|_{V_0})=W_0\times \{0\}`$ and $`\phi _{ij}(C_j|_{V_0})=W_0\times \{\mathrm{}\}`$ for $`ij`$, where $`W_0:=W\mathrm{\Sigma }`$. We can further assume that there exists $`\psi _{ij}:V^1`$ such that $`\psi _{ij}|_{V_0}=p_2\phi _{ij}`$, where $`p_2`$ is the second projection $`W_0\times ^1^1`$. We also assume that $`_kC_k(f^{}\mathrm{\Sigma })_{\mathrm{red}}`$ is a simple normal crossing divisor. we obtain
$`\psi _{ij}^{}{}_{}{}^{}\left({\displaystyle \frac{dz}{z}}\right)`$ $``$ $`\mathrm{Hom}_{𝒪_V}(f^{}(K_W+\mathrm{\Sigma }),K_V+C_i+C_j+(f^{}\mathrm{\Sigma })_{\mathrm{red}})`$
$``$ $`H^0(V,K_{V/W}+C_i+C_j+(f^{}\mathrm{\Sigma })_{\mathrm{red}}f^{}\mathrm{\Sigma })`$
$``$ $`H^0(V,K_{V/W}+C_i+C_j)`$
for $`ij`$, where $`z`$ denotes a suitable inhomogeneous coordinate of $`^1`$ (see \[F1, Lemma 7.12\]). Therefore,
$$dim_{}H^0(V,K_{V/W}+C_1+C_2)1$$
and
$$dim_{}H^0(V,K_{V/W}+C_1+C_2+C_3)2.$$
Thus, we obtain the required result. ∎
Apply Lemma 2.3 to $`\stackrel{~}{X}Y^{\prime \prime }`$, where $`\beta :\stackrel{~}{X}X^{\prime \prime }`$ is a resolution of $`X^{\prime \prime }`$. Then we obtain
$$\kappa (K_{\stackrel{~}{X}/Y^{\prime \prime }}+\beta ^{}\pi ^{}((p^{}C)_{\mathrm{hor}}))\overline{\kappa }(F_0).$$
Thus, we complete the proof. ∎ |
warning/0507/math0507281.html | ar5iv | text | # On the symplectic volume of the moduli space of Spherical and Euclidean polygons
## 1. Introduction
An polygon in $`𝕊^3`$ or $`𝔼^3`$ is specified by its set of vertices $`v=(v_1,\mathrm{},v_n).`$ This vertices are joined in cyclic order by edges $`e_1,\mathrm{},e_n`$, where $`e_i`$ is the directed geodesic segment from $`v_i`$ to $`v_{i+1}.`$ Two polygons $`P=(v_1,\mathrm{},v_n)`$ and $`Q=(w_1,\mathrm{},w_n)`$ are identified if there exists an orientation preserving isometry sending each $`v_i`$ to $`w_i.`$
The side-length $`r_i`$ of a polygon is defined to be the length of the geodesic segment $`e_i.`$ We say that a polygon is regular if all of its side-lengths are equal.
Let $`r=(r_1,\mathrm{},r_n)`$ be a tuple of real numbers such that $`0<r_i<\pi i,`$ following , we will denote by $`𝒫_r^{𝕊^3}`$ the configuration space of all polygons in $`𝕊^3`$ with side-lengths $`r.`$ The moduli space of polygons in $`𝕊^3`$ with side-lengths $`r`$ is defined to be $`𝒫_r^{𝕊^3}/\mathrm{Iso}=𝒫_r^{𝕊^3}/\mathrm{SO}(4)`$.
In $`𝔼^3,`$ similarly, if $`r=(r_1,\mathrm{},r_n)`$ is a tuple of positive real numbers, we denote by $`𝒫_r^{𝔼^3}`$ the configuration space of all polygons in $`𝔼^3`$ with side-lengths $`r.`$ The moduli space of polygons in $`𝔼^3`$ with side-lengths $`r`$ is defined to be $`𝒫_r^{𝔼^3}/\mathrm{Iso}=𝒫_r^{𝔼^3}/𝔼^+(3).`$
We will drop the superscript and simply write $`𝒫_r`$ when we want to talk about polygons in both $`𝕊^3`$ and $`𝔼^3.`$
We are interested in the question : for which tuple $`r`$ the polygon with side-lengths $`r`$ is the most flexible when it is moved in the space provided that its side-lengths are fixed ? First of all we need to specified what is the exact meaning of “flexible”. When a polygon is moved, it always lies in its configuration space. Therefore, it is natural to think that a polygon with side-lengths $`r`$ is more flexible than another one with side-lengths $`r^{}`$ if the corresponding configuration space $`𝒫_r`$ is “bigger” than $`𝒫_r^{}`$ in some sense. Fortunately, there is a standard symplectic structure on the moduli spaces $`𝒫_r/\mathrm{Iso}`$, see , and we will measure the flexibility of a polygon with side-lengths $`r`$ by the symplectic volume of its moduli space.
The main results are the followings:
###### Theorem 1.1.
Among all the spherical or Euclidean polygons with fixed perimeter, the regular polygon is one of the most flexible. Moreover in the case $`n`$ is even, the regular $`n`$-gon is the unique one with this property.
###### Theorem 1.2.
Among all the spherical polygons the regular one with side-length $`\pi /2`$ is the unique one which is the most flexible.
The rest of this paper is organized as follows. In section 2, we will recall the symplectic structure on the moduli space $`𝒫_r^{𝕊^3}/\mathrm{SO}(4)`$ and using Witten’s formula to find its symplectic volume in terms of Bernoulli polynomials. As a corollary we also derive the positivity of certain trigonometric sums which has been proved by analytic method. Section 3 is devoted to the study of the moduli space $`𝒫_r^{𝔼^3}/𝔼^+(3),`$ in particular we derive an explicit formula for its symplectic volume. The main results are proved in section 4 by a detail study of the symplectic volume of as a function of $`r.`$
## 2. Volume of the moduli space of spherical polygons
We first briefly recall the identification between $`𝒫_r^{𝕊^3}/\mathrm{SO}(4)`$ and the moduli space of flat $`\mathrm{SU}(2)`$ connections on a punctured sphere with fixed holonomies around the punctures. Let $`S_n:=S^2\{p_1,\mathrm{},p_n\}`$ be an $`n`$-punctured sphere. For a tuple of numbers $`r=(r_1,\mathrm{},r_n)`$ such that $`0<r_i<\pi i,`$ we will denote by $`(S_n,r)`$ the moduli space of flat $`\mathrm{SU}(2)`$ connections on $`S_n`$ modulo gauge equivalence such that the holonomy around $`p_i`$ is conjugate to $`\left(\begin{array}{cc}e^{ir_i}& 0\\ 0& e^{ir_i}\end{array}\right).`$
It is well-known that $`(S_n,r)(S_n,r)/\mathrm{SU}(2),`$ where
$$(S_n,r):=\{(g_1,\mathrm{},g_n)\mathrm{SU}(2)^n|g_1\mathrm{}g_n=\mathrm{Id},\mathrm{Tr}(g_j)=2\mathrm{cos}(r_j)j\}$$
is the representation space consists of all representations from $`\pi _1(S_n)`$ to $`\mathrm{SU}(2)`$ such that the image of the loop around the puncture $`p_i`$ is conjugate to $`\left(\begin{array}{cc}e^{ir_i}& 0\\ 0& e^{ir_i}\end{array}\right)`$ and the group $`\mathrm{SU}(2)`$ act diagonally by conjugation.
On the other hand, following , we may identify $`(S_n,r)`$ with the configuration space $`𝒫_{r,0}^{𝕊^3}`$ of based polygons, i.e. polygons having the first vertex $`v_1=\mathrm{Id}\mathrm{SU}(2)𝕊^3.`$ For $`(g_1,\mathrm{},g_n)(S_n,r)`$ we construct a based polygon by defining $`g_0:=\mathrm{Id}`$ and let $`v_j:=g_0\mathrm{}g_{j1}`$ for $`1jn.`$ Under this identification, we have
$$(S_n,r)/\mathrm{SU}(2)𝒫_{r,0}^{𝕊^3}/\mathrm{SO}(3)𝒫_r^{𝕊^3}/\mathrm{SO}(4).$$
Therefore, we have a natural identification $`𝒫_r^{𝕊^3}/\mathrm{SO}(4)(S_n,r).`$ According to Atiyah-Bott , the moduli space of flat connections on a surface of genus $`g`$ with specified holonomies around the boundary is a finite dimensional symplectic manifold, possibly with singularities. Witten compute the symplectic volume of the moduli space of flat connections. Since then, other people have given simplified and rigorous proofs (see ). As a special case of Witten’s theorem, we have the followings (see ):
###### Theorem 2.1.
(Witten’s formula) Denote by $`V^{𝕊^3}(r)`$ the symplectic volume of $`𝒫_r^{𝕊^3}/\mathrm{SO}(4)`$\- the moduli space of polygons of side-lengths $`r=(r_1,\mathrm{},r_n).`$ We have:
$$V^{𝕊^3}(r)=\frac{2^{n1}}{\pi }\underset{k=1}{\overset{\mathrm{}}{}}\frac{\mathrm{sin}kr_1\mathrm{}\mathrm{sin}kr_n}{k^{n2}}.$$
For $`I\{1,2,\mathrm{}n\},`$ let $`\overline{I}`$ denote the complement of $`I,`$ $`|I|`$ be the cardinality of $`I`$ and $`r_I:=_{iI}r_i.`$ It has been shown that the moduli space of polygons of side-lengths $`r=(r_1,\mathrm{},r_n)`$ is non-empty if and only if :
$$r_Ir_{\overline{I}}+(|I|1)\pi I\{1,2,\mathrm{}n\},|I|\text{odd}.$$
See for a proof using quantum cohomology. This result can also be proved in a similar line as in the case of spherical polygons in $`𝕊^2`$ in .
Combining this result with Witten’s formula we obtained the positivity of a trigonometric series.
###### Corollary 2.2.
If $`r=(r_1,\mathrm{},r_n)`$ be a tuple of real numbers such that $`0<r_i<\pi `$ for all $`i,`$ Then
$$\underset{k=1}{\overset{\mathrm{}}{}}\frac{\mathrm{sin}kr_1\mathrm{}\mathrm{sin}kr_n}{k^{n2}}0.$$
Moreover, the inequality is strict if and only if $`r_Ir_{\overline{I}}+(|I|1)\pi `$ for all $`I\{1,2,\mathrm{}n\}`$, $`|I|`$ odd.
We can find a closed form expression of the volume function in terms of the Bernoulli polynomials.
###### Proposition 2.3.
The symplectic volumes of the moduli space of $`2n`$-gons and $`(2n+1)`$-gons are given respectively by:
$$V^{𝕊^3}(r)=\frac{(2\pi )^{2n3}}{2(2n2)!}\underset{I\{1,2,\mathrm{}2n\}}{}(1)^{|I|}B_{2n2}(\left\{(\underset{iI}{}r_i\underset{i\overline{I}}{}r_i)/2\pi \right\})$$
$$V^{𝕊^3}(r)=\frac{(2\pi )^{2n2}}{(2n1)!}\underset{I\{1,2,\mathrm{}2n+1\},|I|\text{odd}}{}B_{2n1}(\left\{(\underset{iI}{}r_i\underset{i\overline{I}}{}r_i)/2\pi \right\}).$$
Where, $`B_n`$ are the Bernoulli polynomials defined by the generating functions
$$\frac{te^{xt}}{e^t1}=\underset{n=0}{\overset{\mathrm{}}{}}B_n(x)\frac{t^n}{n!}$$
and $`\{\}`$ denote the fractional part.
###### Proof.
In the proof of this proposition we will use the following formulae ( see , 1.443) :
$$\underset{k=1}{\overset{\mathrm{}}{}}\frac{\mathrm{cos}k\pi x}{k^{2n}}=(1)^{n1}\frac{(2\pi )^{2n}}{2(2n)!}B_{2n}(\frac{x}{2})(1)$$
$$\underset{k=1}{\overset{\mathrm{}}{}}\frac{\mathrm{sin}k\pi x}{k^{2n+1}}=(1)^{n1}\frac{(2\pi )^{2n+1}}{2(2n+1)!}B_{2n+1}(\frac{x}{2})(2)$$
Where, $`B_n`$ is the Bernoulli polynomials defined above and the formulae hold for $`0x2.`$
We first prove the $`2n`$-gon case. Notice that we can write,
$`\mathrm{sin}kr_1\mathrm{}\mathrm{sin}kr_{2n}`$ $`=`$ $`{\displaystyle \frac{(e^{ikr_1}e^{ikr_1})\mathrm{}(e^{ikr_{2n}}e^{ikr_{2n}})}{(2i)^{2n}}}`$
$`=`$ $`{\displaystyle \frac{_{I\{1,2,\mathrm{}2n\}}(1)^{|I^{}|}e^{ik(_{iI}r_i_{iI^{}}r_i)}}{(2i)^{2n}}}`$
$`=`$ $`{\displaystyle \frac{_{I\{1,2,\mathrm{}2n\}}(1)^{|I|}\mathrm{cos}k(_{iI}r_i_{iI^{}}r_i)}{(2i)^{2n}}}.`$
In the last line, we have replaced $`|I^{}|`$ by $`|I|`$ since $`|I|+|I^{}|=2n.`$ Substitute this in to Witten’s formula we get:
$$V^{𝕊^3}(r)=\frac{(1)^n}{2\pi }\underset{I\{1,2,\mathrm{}2n\}}{}(1)^{|I|}\underset{k=1}{\overset{\mathrm{}}{}}\frac{\mathrm{cos}k(_{iI}r_i_{iI^{}}r_i)}{k^{2n2}}$$
Now using (1) above, we can deduce the result.
Now consider the $`2n+1`$-gon case. Similarly, we can write
$`\mathrm{sin}kr_1\mathrm{}\mathrm{sin}kr_{2n+1}`$ $`=`$ $`{\displaystyle \frac{(e^{ikr_1}e^{ikr_1})\mathrm{}(e^{ikr_{2n+1}}e^{ikr_{2n+1}})}{(2i)^{2n+1}}}`$
$`=`$ $`{\displaystyle \frac{_{I\{1,2,\mathrm{}2n+1\}}(1)^{|I^{}|}e^{ik(_{iI}r_i_{iI^{}}r_i)}}{(2i)^{2n+1}}}`$
$`=`$ $`{\displaystyle \frac{_{I\{1,2,\mathrm{}2n+1\},|I|\text{odd}}\mathrm{sin}k(_{iI}r_i_{iI^{}}r_i)}{(2i)^{2n}}}.`$
We can obtain the formula by plugging this into Witten’s formula and then use formula (2).
## 3. Volume of the moduli space of Euclidean polygons
In Euclidean space, the moduli space of polygons $`𝒫_r^{𝔼^3}/𝔼^+(3)`$ also has a symplectic structure which can be briefly described as follows.
Let $`r=(r_1,\mathrm{},r_n)`$ be a tuple of positive numbers, denote by $`𝒫_{r,0}^{𝔼^3}`$ the configuration space of based polygons, i.e. polygons having the first vertex $`v_1=0𝔼^3.`$ Consider the map:
$`\mathrm{\Phi }_r:𝒫_{r,0}^{𝔼^3}`$ $``$ $`(S^2)^n`$
$`(v_1,\mathrm{},v_n)`$ $``$ $`({\displaystyle \frac{v_2}{r_1}},{\displaystyle \frac{v_3v_2}{r_2}},\mathrm{},{\displaystyle \frac{v_nv_{n1}}{r_{n1}}},{\displaystyle \frac{v_n}{r_n}}).`$
If we defined $`\mu _r:(S^2)^n𝐑:(u_1,\mathrm{},u_n)r_1u_1+\mathrm{}+r_nu_n`$, then $`\mathrm{\Phi }_r`$ define an $`\mathrm{SO}(3)`$-equivariant map between $`𝒫_{r,0}^{𝔼^3}`$ and $`\mu _r^1(0)`$ which results in a diffeomorphism between $`𝒫_{r,0}^{𝔼^3}/\mathrm{SO}(3)`$ and $`\mu _r^1(0)/\mathrm{SO}(3).`$
It is well-known that the unit sphere $`S^2`$ is a symplectic manifold where the symplectic $`2`$-form is the volume form. Consider the symplectic manifold $`(S^2)^n`$ with the weighted symplectic structure $`r_1,\mathrm{},r_n`$, i.e., the symplectic $`2`$-form is $`r_1\omega _1+\mathrm{}+r_n\omega _n,`$ where $`\omega _i`$ is the volume form on the $`i^{th}`$ component. It is well-known that $`\mu _r`$ is the moment map for $`(S^2)^n`$ with the diagonal action of $`\mathrm{SO}(3).`$ Therefore, we can identify $`𝒫_r^{𝔼^3}/𝔼^+(3)𝒫_{r,0}^{𝔼^3}/\mathrm{SO}(3)`$ with the weighted symplectic quotient of $`(S^2)^n`$ by the diagonal action of $`\mathrm{SO}(3).`$
Moreover, we can deduce from a result of L. Jeffrey , Theorem 6.6 (see also the discussion at the end of ) that the moduli space of polygons in $`𝕊^3`$ and $`𝔼^3`$ with the same side-lengths $`r`$ are symplectomorphic provided that $`r_i`$ are sufficiently small. So the volume of the moduli space of polygons in $`𝕊^3`$ and $`𝔼^3`$ are the same if the side-lengths are sufficiently small. Using this fact, we can express the volume of $`𝒫_r^{𝔼^3}/𝔼^+(3)`$ in a quite explicit form.
###### Proposition 3.1.
The symplectic volumes of the moduli space of normalized Euclidean $`2n`$-gon and $`(2n+1)`$-gon are given respectively by:
$$V^{𝔼^3}(r)=\frac{1}{4(2n3)!}\underset{I\{1,2,\mathrm{}2n\}}{}(1)^{|I|}|\underset{iI}{}r_i\underset{i\overline{I}}{}r_i|^{2n3}$$
$$V^{𝔼^3}(r)=\frac{1}{2(2n2)!}\underset{I\{1,2,\mathrm{}2n+1\},|I|\text{odd}}{}\mathrm{sign}(\underset{iI}{}r_i\underset{i\overline{I}}{}r_i)(\underset{iI}{}r_i\underset{i\overline{I}}{}r_i)^{2n2}.$$
###### Proof.
Notice that if we scale $`r`$ by a scalar $`\lambda >0`$ then $`𝒫_{\lambda r}^{𝔼^3}/𝔼^+(3)`$ is the weighted symplectic quotient of $`(S^2)^n`$ for the weight $`\lambda r.`$ Moreover, we know that $`𝒫_r^{𝔼^3}/𝔼^+(3)`$ is of dimension $`2(n3)`$, it follows that $`V^{𝔼^3}(\lambda r)=\lambda ^{n3}V^{𝔼^3}(r).`$
We first, assume that all $`r_i`$ are sufficiently small, from the discussion above, we know that the volume of $`𝒫_r^{𝔼^3}/𝔼^+(3)`$ can also be expressed in terms of Bernoulli polynomials as in Proposition 2.3. When all $`r_i`$ are are sufficiently small, we have:
$$\{(\underset{iI}{}r_i\underset{i\overline{I}}{}r_i)/2\pi \}=\{\begin{array}{cc}(_{iI}r_i_{i\overline{I}}r_i)/2\pi ,& \text{i}f_{iI}r_i_{i\overline{I}}r_i0\\ 1+(_{iI}r_i_{i\overline{I}}r_i)/2\pi ,& \text{i}f_{iI}r_i_{i\overline{I}}r_i<0.\end{array}$$
We recall the two following property of Bernoulli polynomials :
a) $`B_n(1x)=(1)^nB_n(x)`$
b) $`B_n(x)=(B+x)^n`$, where the symbol $`B^n`$ is the Bernoulli number $`B_n(0).`$
By the homogeneity of $`V^{𝔼^3}(r)`$, we know that to find $`V^{𝔼^3}(r)`$ it is enough to collect the terms of order $`(n3)`$ from the formulae of Proposition 2.3. Now using this two properties of Bernoulli polynomials and notice that $`B_1=1/2`$, one can easily obtain the required results. So our formulae hold in the case where all $`r_i`$ are sufficiently small. But as notice before, $`V^{𝔼^3}(r)`$ are homogenous of degree $`n3`$ with respect to multiplication by a positive scalar, therefore the formulae hold for all $`r.`$
As a special case of our result, when all $`r_i=1,`$ we obtain the following formula which has been shown by direct computation in .
###### Corollary 3.2.
The volume of the moduli space of regular polygon with side-length $`1`$ in $`𝔼^3`$ is given by:
$$V^{𝔼^3}(1)=\frac{1}{2(n3)!}\underset{k=0}{\overset{[n/2]}{}}(1)^k\left(\genfrac{}{}{0pt}{}{n}{k}\right)(n2k)^{n3}.$$
###### Proof.
The $`2n`$-gon case is straightforwards, we will give the proof for the case of $`2n+1`$-gon. From Proposition 3.1, we get
$$V^{𝔼^3}(1)=\frac{1}{2(2n2)!}(\underset{k\text{o}dd,k>n}{}\left(\genfrac{}{}{0pt}{}{2n+1}{k}\right)(2n+12k)^{2n2}\underset{k\text{o}dd,kn}{}\left(\genfrac{}{}{0pt}{}{2n+1}{k}\right)(2n+12k)^{2n2}).$$
By changing the index in the first sum from $`k`$ to $`(2n+1k)`$ we get,
$$V^{𝔼^3}(1)=\frac{1}{2(2n2)!}(\underset{k\text{e}ven,kn}{}\left(\genfrac{}{}{0pt}{}{2n+1}{k}\right)(2n+12k)^{2n2}\underset{k\text{o}dd,kn}{}\left(\genfrac{}{}{0pt}{}{2n+1}{k}\right)(2n+12k)^{2n2})$$
$`=\frac{1}{2(2n2)!}_{k=0}^n(1)^k\left(\genfrac{}{}{0pt}{}{2n+1}{k}\right)(2n+12k)^{2n2}.`$
This finishes the proof. ∎
## 4. Proof of the main results
In this section we will give the proofs of Theorem 1.1 and 1.2 by study the maximum of the volume function which is the trigonometric series given by Theorem 2.1. From Proposition 2.3, we can easily find out that for $`n=3`$, the moduli space of polygons has volume $`1,`$ no matter what the side-lengths are. So we only need to deal with the non-trivial cases and may assume that $`n>3.`$
Proof of Theorem 1.1. It is enough to prove the spherical case.
For a fixed $`n,`$ consider the set of all polygon whose perimeter is $`P.`$ To prove the first part of our theorem, we must show the following :
Claim: Let
$$f(x_1,\mathrm{},x_n)=\underset{k=1}{\overset{\mathrm{}}{}}\frac{\mathrm{sin}kx_1\mathrm{}\mathrm{sin}kx_n}{k^{n2}}$$
be a function on the domain $`D=\{(x_1,\mathrm{},x_n)|0<x_i<\pi ,_1^nx_i=P\}.`$ Then $`f(x)f(x^{})`$ for all $`xD`$ where $`x^{}=(\frac{P}{n},\mathrm{},\frac{P}{n}).`$
To prove this claim we will show that starting from any point $`xD,xx^{}`$ we can build a sequence $`\{x^k\}_{k=1..\mathrm{}},x^kD,`$ such that:
i) $`x^1=x.`$
ii) $`f(x)`$ is non-decreasing on the segment connecting $`x^i`$ and $`x^{i+1}.`$
iii) $`lim_k\mathrm{}x^k=x^{}.`$
Once we have done this, it follows immediately that $`f(x)lim_k\mathrm{}f(x^k)=f(x^{}).`$
We construct the sequence $`\{x^k\}`$ as follows. Put $`x^1=x.`$ Suppose that we already have $`x^k,`$ we will find $`x^{k+1}`$ as follows:
\- If $`x^k=x^{},`$ put $`x^{k+1}=x^{}.`$
\- If $`x^k=(x_1^k,\mathrm{},x_n^k)x^{},`$ suppose that $`x_m^k=\mathrm{min}_{i=1..n}x_i^k\text{and}x_M^k=\mathrm{max}_{i=1..n}x_i^k,`$ we then specified $`x^{k+1}`$ by :
$$x_i^{k+1}:=\{\begin{array}{ccc}x_i^k& \text{for}& im,M\\ \frac{x_m^k+x_M^k}{2}& \text{for}& i=m\text{or}i=M.\end{array}$$
It is not hard to see that the sequence constructed this way converges to $`x^{}.`$ It remains to prove that it satisfies (ii) above. Without loss of generality, it is enough to show that if $`x=(x_1,x_2,\mathrm{},x_n)`$ and $`x_1=\mathrm{max}_{i=1..n}x_i,`$ $`x_2=\mathrm{min}_{i=1..n}x_i`$ then $`f(x)`$ is non-decreasing on the segment $`l(t):=(x_1\frac{t(x_1x_2)}{2},x_2+\frac{t(x_1x_2)}{2},x_3,\mathrm{},x_n),0t1.`$
When restricted to the segment $`l(t),`$ our function has the form :
$`f(t)`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{sin}k(x_1\frac{t(x_1x_2)}{2})\mathrm{sin}k(x_2+\frac{t(x_1x_2)}{2})\mathrm{}\mathrm{sin}kx_n}{k^{n2}}}`$
$`=`$ $`{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\mathrm{cos}k(x_1x_2t(x_1x_2))\mathrm{cos}k(x_1+x_2))\mathrm{sin}kx_3\mathrm{}\mathrm{sin}kx_n}{2k^{n2}}}.`$
We compute :
$$f^{}(t)=(x_1x_2)\underset{k=1}{\overset{\mathrm{}}{}}\frac{\mathrm{sin}k(x_1x_2t(x_1x_2))\mathrm{sin}kx_3\mathrm{}\mathrm{sin}kx_n}{2k^{n3}}.$$
Since $`x_1>x_2,`$ by Corollary 2.2, we deduce that $`f^{}(t)0,0t1.`$ So we obtain the first part of the theorem.
For $`n`$ even, by Corollary 2.2, it is easy to see that $`f^{}(t)`$ above is strictly positive for $`x`$ sufficiently close to $`x^{}.`$ We conclude that $`x^{}`$ is the unique point where $`f`$ attains its maximum value.
For $`n`$ odd, in order for the function $`f(x)`$ not to be identically zero on the domain $`D,`$ by Corollary 2.2, we must have $`P(n1)\pi .`$ Even with this assumption, in certain cases, $`x^{}`$ may not be the unique point where $`f`$ attains its maximum value. For example, when $`n=5`$ and $`P=4\pi ,`$ then it is not hard to check that
$$f(\frac{4\pi }{5},\frac{4\pi }{5},\frac{4\pi }{5},\frac{4\pi }{5},\frac{4\pi }{5})=f(\frac{4\pi }{5}+ϵ,\frac{4\pi }{5}ϵ,\frac{4\pi }{5},\frac{4\pi }{5},\frac{4\pi }{5})$$
for all $`ϵ>0`$ sufficiently small.
Proof of Theorem 1.2. It follows from the proof of Theorem 1.1 that a regular polygon is the unique one which is the most flexible in the set of all polygons with the same perimeter. So, by Theorem 1.1, it is enough to restrict ourself to the regular polygons.
Let $`0<x<\pi ,`$ and denote by $`V_n(x)`$ the symplectic volume of the moduli space of regular spherical $`n`$-gons with side $`x`$. From Theorem 2.1, we know that
$$V_n(x)=\frac{2^{n1}}{\pi }\underset{k=1}{\overset{\mathrm{}}{}}\frac{(\mathrm{sin}kx)^n}{k^{n2}}.$$
As
$$V_n^{}(x)=\frac{2^{n1}}{\pi }\underset{k=1}{\overset{\mathrm{}}{}}\frac{n(\mathrm{sin}kx)^{n1}\mathrm{cos}kx}{k^{n3}}=\frac{2^{n2}n}{\pi }\underset{k=1}{\overset{\mathrm{}}{}}\frac{(\mathrm{sin}kx)^{n2}\mathrm{sin}2kx}{k^{n3}}.$$
If $`0<x<\pi /2,`$ Corollary 2.2 tells us that $`V_n^{}(x)0.`$
In the case $`\pi /2<x<\pi ,`$ we write :
$$V_n^{}(x)=\frac{2^{n2}n}{\pi }\underset{k=1}{\overset{\mathrm{}}{}}\frac{(\mathrm{sin}kx)^{n2}\mathrm{sin}2k(\pi x)}{k^{n3}}$$
and again by Corollary 2.2, we get $`V_n^{}(x)0.`$
Notice that $`V_n^{}(\pi /2)=0`$ and moreover, by using the second part of Corollary 2.2, it is not hard to check that $`V_n^{}(x)`$ is strictly positive (resp. negative) when $`x`$ is sufficiently close to $`\pi /2`$ on the left (resp. right). From all this we conclude that $`\pi /2`$ is the unique global maximum of $`V_n(x)`$ on the interval $`(0,\pi ).`$ So, we get the conclusion of Theorem 1.2.
###### Remark.
Theorem 1.2 seems to agree with the intuitive reasoning that as the side-lengths are $`\pi /2,`$ each vertex moves on the great sphere, i.e., on the biggest possible space, therefore the polygon is the most flexible. |
warning/0507/math-ph0507036.html | ar5iv | text | # Random discrete Schrödinger operators from Random Matrix Theory
## 1 Introduction
Dyson’s Coulomb gas model for the spectral fluctuations of random matrix ensembles was recently formulated in terms of ensembles of symmetric, real tridiagonal matrices ; see also . These ensembles share the property that the diagonal matrix elements are independent, identically distributed Gaussian random variables, while the off diagonal elements are independent random variables whose probability distribution function (PDF) depends both on the position within the matrix and on the inverse temperature $`\beta `$. We consider these matrix ensembles as ensembles of discrete Schrödinger operators, with random on-site potentials (diagonal matrix elements) and random hopping amplitudes (off-diagonal elements) with prescribed PDF. This viewpoint has been the starting point of recent studies into characterizing the largest eigenvalues of the limiting matrix ensemble in terms of a certain stochastic Schrödinger operator .
The interest in these operators stems also from the fact that they are analogous to a class of operators for which the random potential diminishes as a power law $`|x|^\alpha `$, where $`x`$ marks the position along the chain. Similar systems were thoroughly discussed in the mathematical literature (see e.g, ), where it was proved that a decaying diagonal disorder with $`\alpha <1/2`$ induces localization and the spectrum is pure point. However, for $`\alpha >1/2`$ the states are extended and the spectrum is absolutely continuous. The behavior at the critical power $`\alpha =1/2`$ depends on the details of the potential, and the eigenstates can be either power-law localized or extended. The model we study here is related to this critical class, but not exactly, since in the present case the transition amplitudes are also random variables. We show that in this model, the parameter $`\beta `$ determines the spectral properties: in the regime $`0\beta <2`$ the spectrum is pure -point, and the eigenstates are power law localized, while for $`\beta 2`$, the eigenstates are extended and the spectrum is singular continuous with a $`\beta `$ dependent spectral measure with a Hausdorff dimension $`1\frac{2}{\beta }`$.
We start with a short survey of the relevant information from Random Matrix Theory (RMT). The random matrix ensembles $`GOE,GUE`$ and $`GSE`$ are ensembles of $`N\times N`$ real symmetric, complex hermitian or hermitian real quaternion matrices, respectively, whose matrix elements are independently distributed random Gaussian variables with joint distribution proportional to
$$\mathrm{exp}(c\mathrm{Tr}H^2).$$
(1)
The probability distribution functions of their eigenvalues $`\lambda _1,\mathrm{},\lambda _N`$ can be written in a concise form
$$P_\beta (\lambda _1,\mathrm{},\lambda _N)=\frac{1}{G_{\beta N}}\mathrm{exp}\left(\frac{1}{2}\underset{j=1}{\overset{N}{}}\lambda _j^2\right)\underset{1j<kN}{}|\lambda _j\lambda _k|^\beta .$$
(2)
Here, $`\beta =1,2,4`$ is used for the $`GOE,GUE,GSE`$ ensembles, respectively, $`G_{\beta N}`$ are known normalization constants, and $`c`$ in (1) has been chosen to equal $`1/2,1/2,1/4`$ for the $`GOE,GUE,GSE`$, respectively. It can be shown from (2), or alternatively directly from the definitions of the ensembles by studying their resolvent, that to leading order the normalized spectral density is supported in the interval $`[\sqrt{2\beta N},\sqrt{2\beta N}]`$ and it assumes the “semi-circle” law
$$\rho (\lambda )=\frac{2}{\pi }\frac{1}{\sqrt{2\beta N}}\sqrt{1\frac{\lambda ^2}{2\beta N}}.$$
(3)
Recently, a systematic way to construct the ensembles corresponding to arbitrary (positive) $`\beta `$ was introduced , and it is based on the following observation . Any real symmetric matrix $`𝒜GOE`$ can be orthogonally transformed to a tridiagonal form
$`_N=\left(\begin{array}{ccccccc}a_1& b_1& & & & & \\ b_1& a_2& b_2& & & & \\ & b_2& a_3& b_3& & & \\ & & & & & & \\ & & & & & & \\ & & & & b_{N2}& a_{N1}& b_{N1}\\ & & & & & b_{N1}& a_N\end{array}\right)`$ (11)
The probability distribution function of the matrix elements of the corresponding tridiagonal matrix $`_N`$ has the following properties:
* The diagonal elements $`\{a_n\}`$ are real, independent, identically distributed, Gaussian random variables.
* The off-diagonal elements $`\{b_n\}`$ are non-negative, independently distributed random variables, with PDF
$$P_{GOE}(b_n)=\chi _n(b_n)\frac{2}{\mathrm{\Gamma }(\frac{n}{2})}(b_n)^{n1}\mathrm{e}^{b_n^2}.$$
(12)
The surprising new result is that by distributing the off diagonal matrix elements using the PDF
$$P_\beta (b_n)=\chi _{\beta n}(b_n)\frac{2}{\mathrm{\Gamma }(\frac{\beta n}{2})}(b_n)^{\beta n1}\mathrm{e}^{b_n^2},$$
(13)
the eigenvalue PDF of $`_N`$ is given by (2) for any positive $`\beta `$. Thus, the study of the tridiagonal ensembles (denoted by $`G\beta E`$) provides a convenient way to interpolate between the classical random matrix ensembles with the discrete $`\beta =1,2,4`$. (A similar method was recently applied in to Dyson’s ensembles of unitary matrices, to get Circular $`\beta `$ Ensembles; this was used in to calculate eigenvalue statistics for CMV matrices).
Denoting by $`_\beta `$ the expectation value with respect to the $`G\beta E`$ measures, we can easily find,
$$b_n_\beta =\frac{\mathrm{\Gamma }\left(\frac{\beta n+1}{2}\right)}{\mathrm{\Gamma }\left(\frac{\beta n}{2}\right)}=\sqrt{\frac{\beta n}{2}}\left(1\frac{1}{4\beta n}\right)+𝒪\left(\frac{1}{n^{\frac{3}{2}}}\right)$$
(14)
and,
$$\left(b_nb_n\right)^2_\beta =\frac{1}{4}+𝒪\left(\frac{1}{n}\right).$$
(15)
Thus, for large $`n`$, the PDF (13) limits to the Dirac distribution $`\delta (u_n1)`$ in the normalized variable defined by $`b_n=\sqrt{\frac{\beta n}{2}}u_n`$. This also shows that by scaling the matrix elements $`_N\sqrt{2/\beta N}_N`$, the new off diagonal elements decay as $`n^{1/2}`$ where $`n`$ is counted from the bottom row of the matrix.
Once the matrix $`𝒜`$ under consideration is in tridiagonal form (11), a simple recursion relation can be written for the characteristic polynomial $`D_N(\lambda ):=det\left(\lambda I𝒜\right)=det\left(\lambda I_N\right)`$. Denoting the determinant of the top $`n\times n`$ sub-block of $`\lambda I_N`$ by $`D_n(\lambda )`$, expansion by the last row shows
$$D_n=(\lambda a_n)D_{n1}b_{n1}^2D_{n2};1nN,$$
(16)
subject to the initial conditions
$$D_1=0,D_0=1.$$
(17)
We remark that by computing the zeros of the characteristic polynomial for the tridiagonal matrices (11) one is sampling from the *correlated* PDF (2).
The matrix $`_N`$, in the limit $`N\mathrm{}`$ can be considered as a representation of a discrete quantum hamiltonian which governs the dynamics of a quantum particle hopping randomly between sites on the half line. The distribution of the “on-site potentials” and “hopping amplitudes” are provided by the PDF of the $`a_n`$ and the $`b_n`$ respectively. In the mathematics literature, this is referred to as a discrete random Schrödinger operator, or a random Jacobi matrix. We address the following questions: *i.* Whether, for almost all realizations, the eigenfunctions of the random hamiltonian are localized, or in other words, if the spectrum is continuous or discrete. *ii.* In what way the localization depends on the parameter $`\beta `$.
Consider the matrix (11) for a finite $`N`$. The eigenvectors $`𝐯=(v_1,\mathrm{},v_N)`$ satisfy
$$𝐯=\lambda 𝐯b_{n1}v_{n1}+(a_n\lambda )v_n+b_nv_{n+1}=0,1nN,$$
(18)
with the boundary conditions
$$v_0=v_{N+1}=0.$$
(19)
The homogenous boundary conditions (19) can be satisfied only for $`N`$ discrete values of $`\lambda `$, and this set coincides with the zeros of the characteristic polynomial $`p_N(\lambda )`$.
Ignoring the boundary conditions for a while, the recursion relations (18) can be solved for any $`\lambda `$ and $`N`$, by two independent vectors $`𝐱`$ and $`𝐲`$. The Wronskian
$$W(x,y)=b_n\left(x_{n+1}y_nx_ny_{n+1}\right)$$
(20)
is independent of $`n`$ and therefore the growth rate of the vectors compensate each other so that the Wronskian remains constant. It is convenient to chose one of the vectors, say x as the solution which satisfies the *initial* conditions,
$$x_0=0,x_1=1.$$
(21)
It can be computed (for any $`\lambda `$) by forward iterations of (18). Comparing the two initial value problems (16) , (19) and (18), (21), we find that
$$D_n=x_{n+1}\underset{m=1}{\overset{n}{}}b_m,$$
(22)
which can be proved by direct substitution. Note that the forward iterations usually pick up the solution with the fastest growing rate. An independent solution of the recursion relation can be obtained by imposing the condition $`y_{N+1}=0,y_N=1`$ at an arbitrary value of $`N`$ and perform a backward iteration of (18). This can be done for every $`\lambda `$, and in most cases the solution to be picked up is the one for which $`|y_n|`$ is the fastest increasing solution when $`n`$ is decreasing (for $`0nN`$). The Wronskian relation (20) implies
$$y_0=\frac{b_Nx_{N+1}}{b_0}.$$
(23)
Before addressing the effect of randomness it is useful and instructive to study first the one parametric family of mean hamiltonians which are obtained by replacing $`a_n`$ and $`b_n`$ in (11) by their $`G\beta E`$ expectation values. This way we can better appreciate the effect of randomness on the quantum dynamics. We shall show that the eigenfunctions of the mean hamiltonians are extended, and the spectra are absolutely continuous for all $`\beta >0`$.
The mean hamiltonians $`_\beta `$ are tridiagonal matrices with vanishing diagonal matrix elements. The off diagonal terms are given by (14), and, to leading order, are proportional to $`\sqrt{n}`$. Thus, for large $`n`$, the recursion relations for the components of an eigenvector are:
$$\sqrt{n1}x_{n1}+\sqrt{n}x_{n+1}=\sqrt{2}\stackrel{~}{\lambda }x_n,$$
(24)
where $`\stackrel{~}{\lambda }=\frac{\lambda }{\sqrt{\beta }}`$. The solution of this recursion relation subject to the initial condition $`x_0=0,x_1=1`$ can be written in terms of the normalized eigenfunctions of the one dimensional harmonic oscillator
$$x_{n+1}=u_n(\stackrel{~}{\lambda })=\left(\frac{1}{\sqrt{\pi }n!2^n}\right)^{\frac{1}{2}}\mathrm{e}^{\frac{\stackrel{~}{\lambda }^2}{2}}H_n(\stackrel{~}{\lambda }),$$
(25)
with $`u_1(\stackrel{~}{\lambda })=0`$. The completeness and orthonormality of the Hermite polynomials implies that for any real $`\lambda ,\mu `$,
$$\underset{m=0}{\overset{\mathrm{}}{}}u_m(\lambda )u_m(\mu )=\delta (\lambda \mu ).$$
(26)
This proves that the spectrum of the operator $`_N_\beta `$ for $`N\mathrm{}`$ is absolutely continuous and supported on the entire real line, for all $`\beta >0`$. For finite matrices, the boundary condition $`v_{N+1}=0`$ is satisfied if $`\stackrel{~}{\lambda }`$ is chosen as one of the zeros of the Hermite polynomial $`H_N(\stackrel{~}{\lambda })`$. For finite but large $`N`$ the spectrum is located in an interval of size $`2\sqrt{2N}`$ centered at $`\lambda =0`$. The normalized spectral density $`\rho (\mu =\stackrel{~}{\lambda }/\sqrt{2N})`$ is supported on the interval $`[1,1]`$, and approaches the semi-circle law
$$\rho (\mu )=\frac{2}{\pi }\sqrt{1\mu ^2}$$
(27)
in the limit $`N\mathrm{}`$.
In the subsequent paragraphs, we shall show that, in contrast with the eigenfunctions of the mean Schrödinger operators which are delocalized, the eigenfunctions of the disordered operators are power law localized for the $`G\beta E`$ ensembles with $`\beta <2`$. Beyond the critical value $`\beta =2`$ the eigenfunctions of $`_N`$ cannot be normalized and the spectrum is continuous. However, the disorder has the effect that now the spectrum is singular continuous with a spectral measure which has a $`\beta `$ dependent Hausdorff dimension $`1\frac{2}{\beta }`$.
A prominent quantity of interest in the study of random Schrödinger operators is the mean growth rate of the eigenvectors $`𝐱`$ . It is related to the properties of the characteristic polynomial by
$$_\beta \frac{1}{n}\mathrm{log}\left|\frac{x_1}{x_{n+1}}\right|_\beta =\frac{1}{n}\mathrm{log}|x_{n+1}|_\beta =\frac{1}{n}\mathrm{log}|D_n|_\beta +\frac{1}{n}\underset{m=1}{\overset{n}{}}\mathrm{log}|b_m|_\beta .$$
(28)
Thus, the mean Lyapunov exponent $`_\beta `$ which characterizes the Anderson model, is expressed in terms of the expectation value of the logarithm of the characteristic polynomial of the $`G\beta E`$ ensemble. Since the latter is known from RMT, and the mean value of the rightmost term in (28) can be evaluated directly, the mean Lyapunov exponent for this model can be written down for any value of $`\lambda `$. Using the exact PDF for the $`b_n`$, we get
$$\frac{1}{n}\underset{m=1}{\overset{n}{}}\mathrm{log}|b_m|=\frac{1}{n}\underset{m=1}{\overset{n}{}}\frac{1}{2}\frac{\mathrm{\Gamma }^{}\left(\frac{m\beta }{2}\right)}{\mathrm{\Gamma }\left(\frac{m\beta }{2}\right)}=\frac{1}{2}\left(\mathrm{log}\frac{n\beta }{2}1+\left(\frac{1}{2}\frac{1}{\beta }\right)\frac{\mathrm{log}n}{n}\right)+𝒪\left(\frac{1}{n}\right).$$
(29)
This is derived starting with the identity
$$\underset{j=0}{\overset{N1}{}}\mathrm{\Gamma }(\alpha +1+jc)=c^{cN(N1)/2+N(\alpha +1/2)}(2\pi )^{N(c1)/2}\underset{p=1}{\overset{c}{}}\underset{j=0}{\overset{N1}{}}\mathrm{\Gamma }\left(\frac{\alpha +p}{c}+j\right)$$
valid for $`c𝐙_+`$. Let
$$f_n(\alpha ,c)=c^{\alpha n}\underset{p=0}{\overset{c1}{}}\frac{G(n+(\alpha p)/c+1)G(p/c+1)}{G(np/c+1)G((\alpha p)/c+1)},$$
where $`G`$ denotes the Barnes $`G`$-function. Using the asymptotic formula
$$\mathrm{log}\left(\frac{G(N+a+1)}{G(N+b+1)}\right)\underset{N\mathrm{}}{}(ba)N+\frac{ab}{2}\mathrm{log}2\pi +\left((ab)N+\frac{a^2b^2}{2}\right)\mathrm{log}N+\mathrm{o}(1)$$
it follows that
$$f_n(\alpha ,c)\mathrm{exp}(\alpha n\mathrm{log}n)c^{\alpha n}e^{\alpha n}n^{(c1)\alpha /2+\alpha ^2/2c}\underset{p=0}{\overset{c1}{}}\frac{G(p/c+1)}{G((\alpha p)/c+1)}.$$
Taking logarithms of both sides, differentiating with respect to $`\alpha `$, then setting $`\alpha =1+c`$ gives (29). Furthermore, if $`c`$ is rational, $`c=s/r`$ for $`s,r𝐙_+`$, the identity
$$f_{rn}(\alpha ,s/r)=\underset{\nu =0}{\overset{r1}{}}\frac{f_n(\alpha +s\nu /r,s)}{f_n(s\nu /r,s)}$$
proves that (29) remains valid for $`n`$ a multiple of $`r`$. Thus if the limit is to exist for non-integer $`c`$ it must be given by (29) .
The $`G\beta E`$ expectation value of $`\mathrm{log}|D_n(\lambda )|`$ is given by
$`{\displaystyle \frac{1}{n}}\mathrm{log}|D_n(\lambda )|_\beta `$ $`=`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{l=1}{\overset{n}{}}}\mathrm{log}|\lambda \lambda _l|_\beta `$ (30)
$`=`$ $`{\displaystyle dy\rho _\beta (y)\mathrm{log}|y\lambda |}={\displaystyle \frac{1}{2}}\left(\mathrm{log}{\displaystyle \frac{n\beta }{2}}1\right)+𝒪\left({\displaystyle \frac{1}{n}}\right),`$
where the exact spectral density was replaced by its semi-circle limit (3), and $`\lambda \sqrt{2n\beta }`$. Substituting in (28), we find that
$$_\beta =\frac{\mathrm{log}n^{\frac{1}{2}(\frac{1}{2}\frac{1}{\beta })}}{n}.$$
(31)
Thus, on average, the components of the eigenvectors $`𝐱`$ behave as
$$|x_n|^2n^{\frac{1}{\beta }\frac{1}{2}}.$$
Using (23) we expect the other solution of the recursion relation to exhibit a mean decay rate of
$$|y_n|^2n^{\frac{1}{\beta }\frac{1}{2}}.$$
This result suggests the following scenario: The eigenvectors are square normalizable when $`\beta 2`$ which would imply that the spectrum is pure point in this $`\beta `$ domain. A transition to a continuous spectrum would be expected in the complementary domain. Indeed, with a little more effort, we could show that this is true.
Consider the matrices:
$`S_n^\lambda =\left(\begin{array}{cc}\frac{\lambda a_n}{b_n}& \frac{b_{n1}}{b_n}\\ 1& 0\end{array}\right)`$ (34)
and their product
$$T_n^\lambda =S_n^\lambda S_{n1}^\lambda \mathrm{}S_1^\lambda .$$
$`T_n^\lambda `$ have the property that for the eigenvectors $`𝐱`$
$$\left(\begin{array}{c}x_{n+1}\\ x_n\end{array}\right)=T_n^\lambda \left(\begin{array}{c}x_1\\ x_0\end{array}\right).$$
(35)
In fact, the above holds for any vector satisfying equation (18) (with any boundary conditions). Thus, in order to control the asymptotics of $`𝐱`$, it is reasonable to try and control the asymptotics of $`T_n^\lambda `$. By an adaptation of methods from (for details see ), it is possible to prove
Proposition 1. For any $`\lambda `$
$$\underset{n\mathrm{}}{lim}\frac{\mathrm{log}T_n^\lambda ^2}{\mathrm{log}n}=\frac{1}{\beta }\frac{1}{2}$$
with probability one.
With this at hand, a direct application of the methods of (see also ) imply
Theorem 1. For any $`\lambda `$, with probability one, equation (11) has a solution $`𝐲`$ satisfying
$$|y_n|^2n^{(\frac{1}{2}+\frac{1}{\beta })}.$$
(36)
Any solution to equation (18) that is linearly independent from $`𝐲`$, satisfies
$$|x_n|^2n^{\frac{1}{\beta }\frac{1}{2}}.$$
(37)
The asymptotic behavior of eigenfunctions is intimately connected with the spectral measure of $`H`$, associated with the vector $`\delta _1=(1,0,0,\mathrm{})`$. This is the object defined by
$$d\mu (E)=\mathrm{w}\mathrm{lim}_N\mathrm{}\underset{n=1}{\overset{N}{}}\left|\delta _1\psi _n\right|^2\delta (\lambda \lambda _n)$$
(38)
where $`\lambda _n`$ are the eigenvalues of $`_N`$, and $`\psi _n`$ are the corresponding normalized eigenvectors. using the technique of spectral averaging (see e.g. ), the following can be shown to ensue from the theory of subordinacy ().
Theorem 2. For any $`\beta `$ the essential spectrum of $`H`$ is $``$.
If $`\beta <2`$, then, with probability one, $`\mu `$ is pure point with eigenfunctions decaying as
$$|x_n|^2n^{(\frac{1}{2}+\frac{1}{\beta })}.$$
If $`\beta 2`$, then with probability one, for any $`\epsilon >0`$, $`\mu `$ is absolutely continuous with respect to $`(1\frac{2}{\beta }\epsilon )`$-dimensional Hausdorff measure and singular with respect to $`(1\frac{2}{\beta }+\epsilon )`$-dimensional Hausdorff measure.
Thus, we see that, as long as $`\beta <2`$, $`H`$ has square-summable eigenfunctions whereas for $`\beta 2`$ the spectrum is purely continuous. This spectral transition at $`\beta =2`$ from pure-point to continuous spectrum is, in a certain sense, continuous in $`\beta `$, since the decay rate of the eigenvalues ($`\frac{1}{2}+\frac{1}{\beta }`$) changes continuously in $`\beta `$. This is in striking contrast with the transition expected in the Anderson model. Thus, for the case studied here, the Inverse Participation Ratio, for example, should change continuously from $`1`$ (when $`\beta =0`$) to $`0`$ (for $`\beta =\mathrm{}`$).
An interesting feature of this continuous transition is the connection between the ‘extendedness’ of states and the level repulsion observed on finite scales. Equation (2) above implies that the probability of finding pairs of close eigenvalues, of the finite dimensional matrix $`_N`$, diminishes as $`\beta `$ increases. Theorem 1 says that as $`\beta `$ increases, the states decay at a slower rate. These two facts are complementary: We expect slower decay rate to be connected with stronger level repulsion. Our analysis confirms this expectation and gives physical meaning to this aspect of the eigenvalue statistics of $`G\beta E`$.
*Acknowledgements:* This research of US was supported by the EPSRC grant GR/T06872/01 and by the Institute of Advance Studies, Bristol University. The work of PJF was supported by the Australian Research Council. The work of JB was supported in part by THE ISRAEL SCIENCE FOUNDATION (grants no. 188/02 and 1169/06) and by Grant no. 2002068 from the United States-Israel Binational Science Foundation (BSF), Jerusalem, Israel. We would like to thank J.P. Keating, J. Marklof, M. Aizenman, S. Warzel, R. Sims, S. Molchanov and Y. Last for inspiring discussions and comments. US would like to thank the School of Mathematics at the University of Bristol for the hospitality extended during his stay there, and PJF would like to thank F. Mezzardi for facilitating this collaboration by hosting his visit to Bristol.
Bibliography |
warning/0507/cond-mat0507332.html | ar5iv | text | # Three types of spectra in one-dimensional systems with random correlated binary potential
## Abstract
The stationary one-dimensional tight-binding Schrödinger equation with a weak diagonal long-range correlated disorder in the potential is studied. An algorithm for constructing the discrete binary on-site potential exhibiting a hybrid spectrum with three different spectral components (absolutely continuous, singular continuous and point) ordered in any predefined manner in the region of energy and/or wave number is presented. A new approach to generating a binary sequence with the long-range memory based on a concept of *additive* Markov chains (Phys. Rev. E 68, 061107 (2003)) is used.
In recent years, the problem of transport of electro-magnetic waves (and various excitations in solids) in one-dimensional (1D) systems with the long-range correlated disorder has attracted much attention IK ; T ; IM . The important significance of this problem is due to exciting results that revise a commonly accepted belief that any randomness (no matter how weak the randomness is) in 1D structures results in the Anderson localization An . It was also believed that 1D systems could not display a complex dynamic feature such as the metal-insulator transition, which gives rise to an appearance of the mobility edges. However, in Refs. Fl ; Bo , a highly nontrivial role of correlations was shown, the divergency of the localization length for some specific values of energy was observed. Using a perturbative approach, it was also shown in Ref. IK that the position of the mobility edges and the windows of transparency can be controlled by the form of the binary correlator of a scattering potential, which was supposed to be continuous, of a long-range and Gaussian. If the correlations are of a long-range, i.e. they decay by the power law, a continuum of extended states may appear in the energy spectrum. In other words, the long-range correlations are necessary for the observation of the metal-insulator transition and co-existence of different type of spectra. A quite simple method was used therein to construct a random potential with a long-range correlation. From the experimental point of view, the importance of recently obtained results may be explained by a strong impact upon the creation of a new class of electron devices and electromagnetic waveguides, which can be used as window filters having new transmission properties.
In this paper we point out a new method for constructing the long range correlated sequences of a two-valued site potential $`\epsilon (n)`$ with a given correlator and prescribed probability distribution function (PDF), not only Gaussian. For this potential the method given in IK does not work. An attempt to construct a correlated dichotomous sequence where the metal-insulator transition can be observed was made in Ref. CBIS . Later, Ref. CBIS2 , this result was retracted.
We study a stationary one-dimensional tight-binding Schrödinger equation (the Anderson model An ),
$$\psi _{n+1}+\psi _{n1}+(E\epsilon (n))\psi _n=0,$$
(1)
with a site potential $`\epsilon (n)`$ taking on two different values $`\epsilon _0`$ and $`\epsilon _1`$.
Equation (1) is a prototype model describing a propagation of excitations (electromagnetic waves, electrons, photons, phonons) in deterministic ordered or random disordered systems (solids or layered super-lattices). The type of ordering, i.e., the correlations in the site potential $`\epsilon (n)`$ determine the spectrum of excitation. The wave functions of excitations in such systems are usually characterized by the Lyapunov exponent $`\mathrm{\Lambda }`$. Some typical examples of sequences having different types of spectra are (spectrum is a set of allowed energies of excitations in the system):
1. A system with a regular periodic variation of the site potential $`\epsilon (n)`$ possess an *absolutely continuous* (AC) spectrum and describes extended, delocalized states having bounded wave functions with the Lyapunov exponent $`\mathrm{\Lambda }=0`$.
2. A system with a random non-correlated potential $`\epsilon (n)`$ has the allowed states with the exponentially localized wave functions and displays a *point* (discrete) spectrum with the positive Lyapunov exponent (the Anderson localization).
3. A potential constructed using the deterministic quasi-periodic Fibonacci chain (see the definition given below, Eq. (10)) exhibits a *singular continuous* (SC) Cantor-like spectrum characterized by the power-law localized wave functions with $`\mathrm{\Lambda }=0`$. The singular continuous spectrum differs essentially from the absolutely continuous one. This spectrum corresponds to a singular continuous integrated density of states (the number of states having the energies smaller then a given one). The integrated density of states is the continuous function of energy, but its derivation equals to zero almost everywhere. In terms of the forbidden and allowed energy zones, this means that the spectrum consists almost completely of forbidden gaps with the infinite number of allowed energies between them. In other words, the spectrum of quasi-periodic potential sequence possesses a fractal structure.
In each of these structures only one pure type of spectrum is presented. A nontrivial result of co-existence of continuous and point spectra in one dimensional chain with correlated disorder was presented in Ref. IK .
Correlated properties in the site potential $`\epsilon (n)`$ are determined by the pair correlation function,
$$C(r)=<\epsilon (n)\epsilon (n+r)>_s<\epsilon (n)>_s^2,<f(\epsilon (n))>_s=\underset{N\mathrm{}}{lim}\frac{1}{N}\underset{n=1}{\overset{N}{}}f(\epsilon (n)),$$
(2)
where the Cezaro average $`<>_s=lim_N\mathrm{}\frac{1}{N}_{n=1}^N`$ can be treated as spatial one. In the Born approximation the Lyapunov exponent is expressed, see Ref. Ta ; GF ; L ; IK , in terms of the Fourier transform $`\stackrel{~}{K}`$ of this two-point correlation function,
$$\mathrm{\Lambda }(E)=\frac{(\epsilon _1\epsilon _0)^2\stackrel{~}{K}(2k)}{32\mathrm{sin}^2k},k[\pi ,\pi ],K(r)=\frac{C(r)}{C(0)}.$$
(3)
The correlation function $`K(r)`$ and its Fourier transform are connected by the following relations,
$$\stackrel{~}{K}(k)=1+2\underset{r=1}{\overset{\mathrm{}}{}}K(r)\mathrm{cos}(kr),K(r)=\frac{1}{\pi }_0^\pi \stackrel{~}{K}(k)\mathrm{cos}(kr)𝑑k$$
(4)
Here the evenness of the functions $`K(r)`$ and $`\stackrel{~}{K}(k)`$ is taking into account. The energy $`E`$ of the eigenstate and the wave number $`k`$ in Eq. (3) in zero approximation on the strength of disorder are related by the simple formula,
$$E=2\mathrm{cos}k.$$
(5)
Thus, there exists the relation between the correlations in the on-site potential $`K(r)`$ and localization properties of eigenstates expressed in terms of the Lyapunov exponent $`\mathrm{\Lambda }(E)`$ (or by means of integrated density of states). This observation enables one to construct random correlated sequences with prescribed spectral properties and provides a recipe for designing filters of arbitrary complexity. Starting from a desirable spectral dependence $`\mathrm{\Lambda }(E)`$ (or integrated density of states) we have to solve the inverse problem of construction a sequence of ”symbols” $`\epsilon (n)`$. This program was partially implemented in Ref. IK , where the one-dimensional tight-binding Schrödinger equation with *two* kind of spectra (absolutely continuous and point) and Gauss distribution of the site potential was examined. However, it is known RS , that in the general case the space of states can be decomposed into the direct sum of subspaces with the wave functions belonging to the point, singular continuous and absolutely continuous parts of spectra. This statement is closely connected with a possibility for presenting an integrated density of states, as for any monotone function, in the sum of three functions: the stepwise, singular and absolutely continuous ones. Here we point to a general method, which is the extension of that used in Ref. IK , and construct a particular example of 1D system having all three above-mentioned types of spectra Note .
To this end we use a new instrument of constructing the correlated sequence of $`\epsilon (n)`$ with a given correlation function $`K(r)`$ using an *additive* Markov chain uya ; uyakm . Here we give a short description of this method.
Consider a homogeneous binary sequence of symbols, $`\epsilon (i)=\{\epsilon _0,\epsilon _1\}`$, $`i\text{Z}=\mathrm{},2,1,0,1,2,\mathrm{}`$. To determine the $`N`$-step Markov chain we have to introduce the *conditional probability* function $`P(\epsilon (i)\epsilon (iN),\epsilon (iN+1),\mathrm{},\epsilon (i1))`$. It is a probability of occurring the definite symbol $`\epsilon _i`$ (for example, $`\epsilon (i)=\epsilon _1`$) after $`N`$-word $`T_{N,i}`$, where $`T_{N,i}`$ stands for the sequence of symbols $`\epsilon (iN),\epsilon (iN+10),\mathrm{},\epsilon (i1)`$. The *additive* Markov chain is characterized by the conditional probability function of the form
$$P(\epsilon (i)=\epsilon _1T_{N,i})=p_1+\underset{r=1}{\overset{N}{}}F(r)(\epsilon (ir)<\epsilon >).$$
(6)
We refer to the amount $`F(r)`$ (introduced first in Ref. mel ) as the *memory function*. It describes the strength of influence of previous symbol $`\epsilon (ir)`$ $`(r=1,\mathrm{},N)`$ upon a generated one, $`\epsilon (i)`$. Here $`p_1`$ is the relative part of symbols $`\epsilon _1`$ among the total number of symbols in the whole sequence ($`p_1`$ is not necessarily equal to $`1/2`$). There is a single-valued relation between the memory function $`F(r)`$ and the correlation function $`K(r)`$ of the Markov chain, see Ref. mel ; melg ,
$$K(r)=(\epsilon _1\epsilon _0)\underset{r^{}=1}{\overset{N}{}}F(r^{})K(rr^{}),r1.$$
(7)
This equation allows finding the memory function $`F(r)`$ and effective constructing the Markov chain with the obtained conditional probability $`P(..)`$ introduced in Eq. (6).
Let us present a ”*Recipe*” and enumerate the consecutive steps of the general method for constructing the correlated sequence of $`\epsilon _n`$ with desired spectra:
1. Dividing the interval of energy $`E`$ (or, that is the same, of the wave number $`k`$) to the regions where we want to have different types of spectra.
2. Prescribing $`\mathrm{\Lambda }`$ or/and $`\stackrel{~}{K}(2k)`$ in these intervals, we construct ”by hand” the total Fourier transform of correlation function $`\stackrel{~}{K}(2k)`$ for all values of $`k[\pi ,\pi ]`$.
3. Calculating $`K(r)`$ as inverse transform of $`\stackrel{~}{K}(2k)`$.
4. Solving Eq. (7), we determine the memory function $`F(r)`$ and the conditional probability function $`P(..)`$.
5. Constructing sequence of $`\epsilon _0`$ and $`\epsilon _1`$ with obtained function $`P(..)`$.
To construct the total correlation function $`K(r)`$ we have to calculate the ”patrial” ones. Let us present here some simple examples of the correlation functions corresponding to the above-mentioned sequences exhibiting the pure spectra of periodic, random and quasi-periodic chains.
It follows from definition (2) that the correlation function of a chain of alternating potentials $`\epsilon (n)=\mathrm{}\epsilon _0\epsilon _1\epsilon _0\epsilon _1\mathrm{}`$ has the form:
$$C_{a,2}(r)=\frac{(1)^r}{4}(\epsilon _1\epsilon _0)^2.$$
(8)
It follows from the same definition (2) that the correlation function of the sequence with the same potentials $`\epsilon (n)=\mathrm{}\epsilon _0\epsilon _0\epsilon _0\mathrm{}`$ is equal to zero for all distances $`r`$. Nevertheless, the ratio $`C(r)/C(0)=K(r)`$ should be defined as unity. Only in this case Eq. (3) gives the correct value of the Lyapunov exponent.
The correlation function of the non-biased non-correlated random sequence of $`\epsilon _n`$ is
$$C_p(r)=(\epsilon _1\epsilon _0)^2\{\begin{array}{cc}1/4,r=0\hfill & \\ 0,r0.\hfill & \end{array}$$
(9)
Let us derive the correlation function of the quasi-periodic chain of potentials $`\epsilon _n`$ constructed using the Fibonacci chain,
$$\mathrm{\Phi }_\xi (r)=[(r+1)\phi +\xi ][r\phi +\xi ].$$
(10)
as,
$$\epsilon _r(\xi )=\epsilon _0+(\epsilon _1\epsilon _0)\mathrm{\Phi }_\xi (r),\xi [0,1),$$
(11)
Here $`\phi `$ is the so-called golden number, $`\phi =(\sqrt{5}1)/2`$, and $`\xi `$ selects some concrete sequence from among all possible Fibonacci sequences. The geometrical meaning of this construction is explained in Fig. 1. To determine the correlation function of this sequence it is convenient to use an ergodic theorem AA about the homogeneity of distribution of $`y_n=n\phi [n\phi ]`$ on the interval $`0y1`$ for any irrational number $`\phi `$ and calculate correlation function $`C(r)`$ using ensemble average $`<>_a`$ instead of space average (2),
$$C(r)=<\epsilon _n\epsilon _{n+r}>_a<\epsilon _n>_a^2<f(\epsilon _n)>_a=_0^1f(\epsilon _n(\xi ))𝑑\xi .$$
(12)
From a simple geometric consideration, see explanation in Fig. 1, we obtain the following result,
$$C_s(r)=(\epsilon _1\epsilon _0)^2\{\begin{array}{cc}2\phi 1\{\phi r\},\hfill & \{\phi r\}1\phi ,\hfill \\ 3\phi 2,\hfill & 1\phi \{\phi r\}<\phi ,\hfill \\ 2\phi 2+\{\phi r\},\hfill & \phi \{\phi r\}.\hfill \end{array}$$
(13)
Here $`\{x\}`$ is the fractional part of $`x`$. The function $`K_s(r)=C_s(r)/C_s(0)`$ of continuous argument $`r`$ is shown in Fig. 2. Actually, the correlation function is defined for the integer arguments, which values are incommensurable with the period $`1/\phi `$ of function $`C_s(r)`$. The function $`C_s(r)`$ of the discrete argument is shown in Fig. 3. Numerical calculations of $`C_s(r)`$ carried out using spatial averaging (2) exhibit a good agreement with the analytical one (13).
Below we construct an example of the sequence with a prescribed spectral property. Let us suppose that the Schrödinger Eq. (1) has all three kind of spectra. Let $`k_{i1}`$ and $`k_{i2}`$ be the wave vectors corresponding to the top and bottom edges of the bands where the singularly continuous, $`i=s`$, absolutely continuous, $`i=a`$, and point (discrete), $`i=p`$, spectra take place. In fact, we do not need to calculate Fourier transforms $`\stackrel{~}{K}(k)`$ of the correlation function as it is indicated above in the paragraph 3 of the ”*Recipe*”. We can calculate the Fourier transform of the composed correlation function $`K(r)`$ presented in $`k`$-space via characteristic functions $`\chi _{[k_{i1},k_{i2}]}(k)`$ of the intervals $`[k_{i1},k_{i2}]`$,
$$\stackrel{~}{K}(k)=\underset{i=\{a,s,p\}}{}\chi _{[k_{i1},k_{i2}]}(k)\stackrel{~}{K}_i(k),$$
(14)
where characteristic function $`\chi _{[k_1,k_2]}(k)`$ is equal to $`1`$ if $`k[k_1,k_2]`$ and equals $`0`$ otherwise. Straightforward calculations yield:
$$K(r)=\underset{i}{}\left[\frac{k_{i2}k_{i1}}{\pi }K_i(r)+\frac{1}{\pi }\underset{r^{}=1}{\overset{\mathrm{}}{}}[K_i(r+r^{})+K_i(rr^{})]\frac{\mathrm{sin}r^{}k_{i2}\mathrm{sin}r^{}k_{i1}}{r^{}}\right].$$
(15)
Here the functions $`(\mathrm{sin}r^{}k_{i2}\mathrm{sin}r^{}k_{i1})/r^{}`$ are the Fourier transforms of the characteristic functions. In the case $`i=p`$, the contribution of the point spectrum to Eq. (15) is reduced to one term only,
$$[\mathrm{}]=\{\begin{array}{cc}\frac{\mathrm{sin}rk_{p2}\mathrm{sin}rk_{p1}}{\pi r},\hfill & r0,\hfill \\ \frac{k_{i2}k_{i1}}{\pi },\hfill & r=0.\hfill \end{array}$$
(16)
Note that in deriving Eq. (15) we did not use the concrete form of correlation functions Eqs. (8), (9) and (13). In Fig. 4 the result of numerical simulation of the normalized Lyapunov exponent $`\mathrm{\Lambda }^{}(E)=32\mathrm{sin}^2(k)\mathrm{\Lambda }(E)/(\epsilon _1\epsilon _0)^2`$ is presented for the system exhibiting two types of spectrum, namely, absolutely continuous with $`\mathrm{\Lambda }^{}=0`$ at $`0<E<1`$ and point spectrum $`\mathrm{\Lambda }^{}=1.5`$ at $`1<E<2`$ Note2 .
Thus, we have proposed a method for constructing 1D sequences of sites $`\epsilon (n)`$ with a correlated disorder exhibiting a hybrid spectrum with three different spectral components ordered in any predefined manner in the region of energy $`E`$ and/or wave number $`k`$.
Authors are grateful to L. A. Pastur for the discussion which was helpful in the formulation of the problem under study. We thank S. A. Gredeskul, A. A. Krokhin, and F. M. Izrailev for useful and enlightening discussions. |
warning/0507/hep-th0507163.html | ar5iv | text | # Brane Decay from the Origin of Time
## Abstract
We present a novel scenario where matter in a spacetime originates from a decaying brane at the origin of time. The decay could be considered as a “Big Bang”-like event at $`X^0=0`$. The closed string interpretation is a time-dependent spacetime with a semi-infinite time direction, with the initial energy of the brane converted into energy flux from the origin. The open string interpretation can be viewed as a string theoretic non-singular initial condition.
Introduction – Perhaps the most ambitious problem in cosmology is the question of the initial conditions of the Universe. In this paper we present a completely new scenario: brane nucleation and decay from the origin of the spacetime. Various proposals for initial conditions have been formulated in different frameworks (see Carroll:2004pn for a recent discussion). The most famous is the no-boundary proposal, where an expanding universe emerges from an instanton in the Euclidean geometry, thus removing the boundary from the spacetime manifold nobdry . Alternatively, it has been speculated that a preceding long pre-Big Bang phase Gasperini:2002bn ended in a Big Crunch, with the Universe re-emerging into a Big Bang. It has been hoped that this so-called Big Bounce could be reliably described, once the true quantum gravitational effects would be properly understood, perhaps in the context of string theory. This proposal then motivated a variant Khoury:2001wf where the Big Bounce was attributed to a collision of branes moving in a higher dimensional space, either as a singular event or in repeated cycles. The associated spacetimes were modeled in string theory as Lorentzian boost orbifolds, but as of yet it is still an open question whether a detailed understanding of the associated stringy effects Durin:2005ix could lend support to the bounce idea.
A related development is the AdS/CFT duality, or more generally the open-closed duality of string theory. Based on models of stacks of stable D-branes, one can show a duality between the open string theory on the brane and closed string theory in a higher dimensional spacetime. In the bulk spacetime, the additional spacelike direction can be related to renormalization group flow in the lower dimensional effective field theory of open strings on the brane. The branes in question have a timelike worldvolume. It was hoped that one could develop analogous models by replacing the timelike branes with spacelike branes Gutperle:2002ai ; the open string physics on the brane would give rise to a time direction, so that the dual closed string interpretation would be string theory in a higher dimensional time-dependent spacetime.
Natural candidates for the spacelike branes are unstable D-branes or D-brane -anti D-brane configurations. The decay of these brane configurations can be modeled in open string theory in different ways, e.g. as conformal field theory of the open string worldsheet. In this case, the brane decay can be modeled by a rolling tachyonic background Sen:2002nu introduced into the action as an exactly marginal deformation. On the other hand, for the dual spacetime interpretation one needs to find time-dependent (S-brane) solutions of supergravity which are sourced by the associated unstable D-brane configuration. This programme is hampered by several technical difficulties on both sides. It is difficult to construct generic exactly marginal tachyonic deformations. Few examples of S-brane spacetimes are known. Furthermore, string dynamics in such backgrounds is technically very involved, obscuring the problem of identifying interesting scaling limits for simplified dualities.
The prototype scenarios for brane decay come in two basic variants. In the first Sen:2002nu , the brane decay has a finite time origin, but the decay phase is preceded by a build-up phase, where the unstable brane is first created by a carefully fine-tuned closed string initial state. In fact, the whole creation-and-decay process is completely time reflection symmetric. In the second variant Larsen:2002wc , the decay starts at past infinity and continues through all time until future infinity. The first variant has the virtue that the timescale of the decay is a tunable parameter and the decay has a finite origin. But this is achieved at the expense of a highly fine-tuned creation phase. Alternatively, a prescription for brane nucleation in an empty spacetime has been proposed in Lambert:2003zr . The second variant has the virtue of only having decay, but the downside is not having a finite initial time and no tunable timescale.
We now propose a new model where the brane decay starts from a finite time origin, has a tunable timescale for the decay, and no creation phase. This opens up exciting possibilities for scenarios where the unstable brane on one hand provides initial conditions which can be phrased in terms of open strings on the brane and on the other hand has a dual interpretation of a time-dependent spacetime with a finite origin of time. The details have some flavor of the no-boundary proposal in terms of involving physics in the complex time plane, some flavor of the pre-Big Bang idea, and some flavor of Lorentzian string orbifold models, the AdS/CFT duality and perhaps even holography. This paper presents the main features of our model, a more detailed discussion will follow in KKLN .
As this paper was prepared, other new ideas about cosmological singularities were reported in othercoolideas .
The model – We begin by recalling the basic pre-Big Bang idea: to resolve the initial singularity to extend the past history of the Universe, as sketched in Figure 1(a).
The concept of “time” breaks down at the Big Bang singularity. In Big Crunch / Big Bang scenarios one hopes to continue time’s arrow across, by a resolution mechanism Seiberg:2002hr . However, one can ask if time’s arrow could be taken to point in multiple directions from the Big Bang. We could imagine several universes being created, such as is suggested by Figure 1(b). It might be that this allows for a different kind of resolution of the initial singularity. An additional ingredient which can be added to such considerations is CPT. One could postulate that the two Universes are simply CPT reflections of one another. Then it would be possible to identify them under the CPT reflection. An analogous proposal has been made before in the elliptic interpretation of the de Sitter space EdS .
Consider then the two main models for brane decay; for simplicity, we focus on D-branes in bosonic string theory. The first variant, mentioned above, is the full S-brane, associated in the boundary CFT of the open string worldsheet with the spatially homogenous tachyonic deformation
$$T(X^0)=\lambda \mathrm{cosh}(X^0/\sqrt{\alpha ^{}})$$
(1)
where the parameter $`0\lambda \frac{1}{2}`$ controls the lifetime of the brane: the brane exists for a finite time of the order
$$\mathrm{\Delta }X^02\mathrm{ln}(\mathrm{sin}(\pi \lambda )).$$
(2)
The decay starts at $`X^0=0`$ and is preceded by a formation phase, where the brane forms out of a carefully fine-tuned initial closed string configuration, which can be mapped to the final state by time reversal. The situation is analogous to building a bomb by reversing its explosion. From the target space point of view, this appears unnatural. It is possible to isolate the decay phase by considering instead a deformation $`T(X^0)=\lambda e^{X^0}`$, corresponding to the second variant, the half-S-brane. Now the parameter $`\lambda `$ has no physical meaning as it can be absorbed into a shift of the origin of the time coordinate.
The full S-brane in Lorentzian spacetime is related by a Wick rotation $`X^0iX`$ to an infinite periodic array of smeared D-branes in Euclidean space. The tachyonic deformation becomes periodic $`\lambda _\mathrm{\Sigma }\mathrm{cos}(X/\sqrt{\alpha ^{}})`$ and alters the boundary conditions at the endpoint of the open string. Increasing $`\lambda `$ from initial value $`0`$ to a final value $`1/2`$ smoothly deforms the open string boundary conditions from Neumann to Dirichlet, so that the open strings become attached to a periodic array of D-branes localized at $`X=2\pi \sqrt{\alpha ^{}}(n+\frac{1}{2})`$, with $`nZ`$. Conversely, as $`\lambda `$ is decreased from $`1/2`$, the branes are smeared over the spacelike direction transverse to the branes. Wick rotating back to the original Lorentzian direction, the smearing corresponds to slowing down the decay, with the lifetime $`\mathrm{\Delta }X^0`$ related to the smearing of the branes. Ref. Gaiotto:2003rm proposed that generic, suitably symmetric, D-brane configurations in the complexified time plane correspond to time-dependent closed string backgrounds. The D-brane configurations act as sources for the spacetime fields and backreact to the geometry. It is of great interest to find spacetime solutions which are sourced by the brane arrays. Additional discussion and some example solutions can be found eg. in Jones:2004rg .
Let us now return to the spacetimes of Figure 1 with the origin of time $`X^0`$ at the Big Bang. Could one consider brane decay in such backgrounds? In the case of (a), if the Crunch and the Bang are time reverses of each other, one could attempt to construct a full S-brane centered at the Big Bang. However, the spatial slices pass through zero size, and one would need to take this into account. For example, if the spacetime is modeled as a Misner space, one could attempt to construct a full brane solution where the tachyon rolls with respect to the Misner time rather than with respect to the Minkowski time. In the case of (b), if the two branches of the spacetime emanating from $`X^0`$ are mapped to each other under T-reversal, it should be again be possible to introduce the full S-brane. But now there is an additional complication because the spacetime is not globally time-orientable. Let us simplify the problem by not requiring that the spatial sections undergo expansion. As a simple toy model, one can consider the Lorentzian orbifold $`R^{1,d}/Z_2\times M_D`$, where $`Z_2`$ acts as a (C)PT-reflection
$$(X^0,X^1,\mathrm{},X^d)(X^0,X^1,\mathrm{}X^d)$$
(3)
in $`R^{1,d}`$, and $`M_D`$ is a Euclidean space of dimension $`D=25d`$ (this would be replaced by $`9d`$ in the superstring). A priori $`0d25`$. This orbifold has been investigated in Balasubramanian:2002ry and Biswas:2003ku in bosonic string and type II superstring theories. The quantization of untwisted and twisted sectors is straightforward. It was found that ghosts are absent at tree level. In supersymmetric theory, the partition function vanishes just as in the corresponding Euclidean orbifold. There are no closed causal curves after the proper definition of the time function on the fundamental domain. The covering space is then a Minkowski space, except that the time’s arrow points in opposite directions on the two half-spaces, as in Figure 1(b). (If $`d>0`$, the spatial slice at $`X^0=0`$ is reduced in half due to the P-reflection). Moreover, there is no dangerous backreaction as the stress tensor vanishes everywhere after the Big Bang initial slice.
Since the tachyonic deformation on the open string worldsheet is symmetric with respect to the time reversal, the full brane survives the orbifold $`Z_2`$ identification. Consider the closed string boundary state description of the full brane. The starting point of the boundary state construction is the Wick rotation to the Euclidean signature. The boundary state $`|B`$ can be written in terms of the underlying $`SU(2)_L\times SU(2)_R`$ symmetry with the generators $`J^\pm =e^{\pm iX/\sqrt{\alpha ^{}}};J^3=iX,`$ as a linear combination of Ishibashi states $`|j,m,n`$ labeled by the $`SU(2)`$ quantum numbers. The $`Z_2`$ reflection $`XX`$ acts on the $`SU(2)`$ generators by
$$J^\pm J^{};J^3J^3.$$
(4)
The boundary state is symmetric with respect to interchanging $`m,m`$, so it is invariant under (4). The same will be true after the inverse Wick rotation back to Lorentzian signature. But there is an important subtlety involved, first discussed in Biswas:2003ku . When the boundary state $`|B`$ is expressed in the standard oscillator and momentum basis, it is a linear superposition of on-shell closed string states $`|\psi _i,k_i`$ with $`k_i^2=M_i^2`$. But, after the inverse Wick rotation, the $`Z_2`$ reflection acts in time, so it also acts on center-of-mass energy as $`k^0k^0`$. Therefore, the on-shell states $`|\psi _i,k_i`$ in the sum must be doubled to reflect the symmetry of the choice of branch $`k^0=\pm \omega _\stackrel{}{k}`$:
$$|\psi _i,k_{iE}\left(\begin{array}{c}|\psi _i,k_i^0,\stackrel{}{k}_i\\ |\psi _i,k_i^0,\stackrel{}{k}_i\end{array}\right).$$
(5)
For more discussion, see Biswas:2003ku and KKLN .
The physical interpretation of the resulting boundary state is different from that of the standard full S-brane. Rather than corresponding to formation and decay of an unstable brane, it corresponds to a brane decaying into closed strings propagating in opposite directions in the covering space, as depicted in Figure 2 ($`X^0,\stackrel{~}{X}^0>0`$ coordinatize the two sheets).
On the fundamental domain (corresponding to the upper half space), we then obtain a brane at the origin of time, decaying into closed strings. This necessitates a prescription for preparing the brane at the origin of time. Ref. Lambert:2003zr proposed a Hartle-Hawking contour integration prescription for spontaneously nucleating a brane at $`X^0=0`$ which subsequently decays. We can use the same prescription on the fundamental domain. On the covering space, it corresponds to a double contour, with branches approaching the origin along the imaginary time axes of $`X^0,\stackrel{~}{X}^0`$ and then proceeding along real time axis into the opposite $`X^0,\stackrel{~}{X}^0>0`$ directions. The physical interpretation is otherwise the same as in Lambert:2003zr , except that now there is no spacetime before $`X^0=\stackrel{~}{X}^0=0`$ \- the nucleation of the brane can truly be viewed as a ”Big Bang”-like event. The average total number density and total energy density of the emitted closed strings can be calculated in a similar manner as in Lambert:2003zr for the Hartle-Hawking contour, with similar results KKLN .
If all the closed strings in the spacetime originate from the brane decay, the latter would then serve as an initial condition. However, the spacetime is the covering space of an orbifold, so it contains twisted sector strings as well. How would they be connected with the decaying brane? In other words, what is their role in the initial condition proposal – can the twisted strings also be associated with a decaying brane?
It turns out that on the spacetime orbifold there is another class of S-branes or decaying branes to consider. We call these fractional S-branes. They are constructed as follows. Start again by the Wick rotation $`X^0iX`$ to Euclidean signature. Formally in the bosonic theory one can also obtain an array of D-branes at $`X=2\pi \sqrt{\alpha ^{}}n`$ by setting $`\lambda =\frac{1}{2}`$. Tachyonic deformation to $`\lambda <0`$ is not well defined in the bosonic theory because it leads to the region where the tachyon effective potential is unbounded from below. However, in superstring theory the effective potential is bounded from below and symmetric about $`T=0`$, and deformation to $`\lambda =\frac{1}{2}`$ is well defined. We will however continue with the bosonic theory, in order to keep our proposal clear – we expect that our proposal can be extended to superstrings. The boundary state of the decaying brane is usually constructed by first compactifying $`X`$ at selfdual radius, $`XX+2\pi \sqrt{\alpha ^{}}`$, then the deformation acts as an $`SU(2)`$ rotation, and finally a projection back to infinite radius is performed to obtain the final boundary state. Now consider the $`Z_2`$ orbifold identification where $`Z_2`$ acts by $`XX`$. Then $`X=0`$ becomes an orbifold fixed point. In the array of D-branes at $`X=2\pi \sqrt{\alpha ^{}}n`$, the brane at $`X=0`$ must be replaced by a fractional D-brane. If we compactify on a selfdual radius as before, after the $`Z_2`$ identification to $`S^1/Z_2`$ there are two fixed points, at $`X=0`$ and $`X=\pi \sqrt{\alpha ^{}}`$. Ref. Oshikawa:1996dj studied the conformal field theory of a free boson on this orbifold and found that there are 8 fractional boundary states at the fixed point $`X_0=0,\pi \sqrt{\alpha ^{}}`$, where the Dirichlet or Neumann untwisted sector boundary states are combined with twisted sector boundary states. Extending and elaborating a previous analysis by Recknagel:1998ih , we have shown how the tachyonic deformation interpolates through the 8 fractional boundary states as $`\lambda `$ is varied. It is notable that the process has a dual description in terms of a free boson on a circle – without orbifold singularities – where it corresponds to moving an ordinary D-brane around the circle. We have also constructed a projection back to the infinite radius.
We have also considered a more straightforward approach, starting with the open string action with tachyonic deformations at the opposite endpoints,
$$\delta S=\lambda _{_1\mathrm{\Sigma }}\mathrm{cos}(X/\sqrt{\alpha ^{}})\stackrel{~}{\lambda }_{_2\mathrm{\Sigma }}\mathrm{cos}(X/\sqrt{\alpha ^{}})$$
and computed the annulus partition function with the $`Z_2`$ reflection, using the fermionization technique of Polchinski:1994my . After reinterpreting this in closed string variables, the result factorizes into a form corresponding to closed string propagation between two deformed boundary states. The result is in agreement with the construction starting from the selfdual radius. Details will appear in KKLN .
Discussion – We have considered S-branes on a spacetime orbifold and shown that this leads to a new class of S-branes that we have called fractional S-branes. We can see three interesting directions for further study and potential applications. i) Our construction is a toy model for a spacetime with a (fractional) S-brane as the stringy initial state at the origin of time. In particular, a space-filling brane decay provides homogeneous initial conditions. Customarily homogeneity in brane cosmological models is obtained by a collision of two branes which must be initially aligned to be parallel, amounting to a careful fine-tuning. In our proposal there is only a single brane. As the mathematical tools and general understanding of the properties of S-branes and associated spacetime solutions develops, we expect our construction to be one step towards big-bang cosmological toy models where a decaying brane at the finite origin of time leads to a holographic interpretation of the spacetime that is created. ii) It would be interesting to investigate our proposal in the context of string/brane gas cosmology Brandenberger:1988aj . Instead of a space-filling infinite brane, one could start with all the spacelike dimensions being compact. Initially the brane decay produces massive (non-winding and winding) closed strings which cascade to lighter modes, interact and presumably thermalize in the end. The total energy of the system is the initial mass of the unstable brane. The brane decay might be used to ”explain” the origin of the hot string gas which drives the expansion of some of the compact directions. iii) We can also speculate on the relation of our construction to the discussion of the emergence of the arrow of time in Carroll:2004pn , starting from a generic low-entropy state on an ”initial” spatial slice. One of the issues with this scenario was the nature of the initial state of the slice – since the scenario leads to arrows of time pointing away from the slice, we could argue that the initial condition could be chosen to be the fractional S-brane that we constructed here. |
warning/0507/math0507429.html | ar5iv | text | # An introduction to joinings in ergodic theory
## 1. What are we talking about?
### 1.1. Dynamical systems and stationary processes
We call here a *dynamical system* any quadruple of the form $`(X,𝒜,\mu ,T)`$, where $`(X,𝒜,\mu )`$ is a Lebesgue space (or equivalently a standard Borel space equipped with a probability measure $`\mu `$) and $`T`$ an automorphism of $`(X,𝒜,\mu )`$: $`T`$ is a one-to-one measurable transformation of $`X`$ satisfying, for any measurable subset $`A`$ of $`X`$,
$$\mu (T^1A)=\mu (TA)=\mu (A).$$
Throughout this text, we will often use simply the symbol $`T`$ to designate the dynamical system $`(X,𝒜,\mu ,T)`$. We will also often need a second dynamical system $`S`$, which will stand for the quadruple $`(Y,,\nu ,S)`$
Such dynamical systems, which are the objects of interest in measurable ergodic theory, can be considered from a rather probabilistic viewpoint by studying *stationary processes*. Indeed, any measurable map $`\xi _0:XE`$ gives rise to a stationary process $`\xi =(\xi _k)_k`$ defined by $`\xi _k:=\xi _0T^k`$. (The set $`E`$ in which $`\xi `$ takes its values could be any standard Borel space; most of the time, it will be a finite or countable alphabet.) Conversely, to any $`E`$-valued stationary process $`\xi =(\xi _k)_k`$, we can associate in a canonical way a dynamical system: take $`X_\xi :=E^{}`$, the sample space of the whole process, equipped with its Borel $`\sigma `$-algebra $`𝒜_\xi `$ and the law $`\mu _\xi `$ of the process. Since $`\xi `$ is stationary, $`\mu _\xi `$ is invariant by the shift $`T_\xi `$: $`(x_k)_k(x_k^{})_k`$ where $`x_k^{}:=x_{k+1}`$.
###### Definition 1.1.
In the dynamical system $`(X,𝒜,\mu ,T)`$, the stationary process $`\xi =(\xi _0T^k)_k`$ is called a *generating process* if the $`\sigma `$-algebra it generates is, modulo $`\mu `$, the whole $`\sigma `$-algebra $`𝒜`$.
Observe that in the dynamical system $`(X_\xi ,𝒜_\xi ,\mu _\xi ,T_\xi )`$, the process $`\xi `$ (obtained by looking at the coordinates) is always a generating process.
The two dynamical systems $`T`$ and $`S`$ are said to be *isomorphic* if we can find invariant subsets<sup>1</sup><sup>1</sup>1In the sequel, we will not explicitely mention this restriction to subsets of full measure; it should be understood that every map is really defined on a subset of probability 1. $`X_0X`$, $`Y_0Y`$, with $`\mu (X_0)=\nu (Y_0)=1`$, and a measurable one-to-one map $`\psi :X_0Y_0`$ with, for all $`B`$, $`\mu (\psi ^1B)=\nu (B)`$ and such that $`S\psi =\psi T`$. If $`\xi `$ is a generating process for the dynamical system $`T`$, then $`T`$ is isomorphic to the dynamical system $`T_\xi `$ defined above. Recall the following well-known result of ergodic theory, which shows that studying dynamical systems is essentially the same thing as studying stationary processes taking values in a countable alphabet. (A nice proof can be found in .)
###### Theorem 1.2.
Every dynamical system admits a generating process taking its values in a countable alphabet.
### 1.2. Products and factors: arithmetic of dynamical systems
In his famous article initiating the theory of joinings , Furstenberg observes that a kind of arithmetic can be done with dynamical systems: Indeed, there are two natural operations in ergodic theory which present some analogy with the integers.
First, if we are given two dynamical systems $`(X,𝒜,\mu ,T)`$ and $`(Y,,\nu ,S)`$, we can construct their *direct product* $`(X\times Y,𝒜,\mu \nu ,T\times S)`$ where $`T\times S:(x,y)(Tx,Sy)`$. Note that the Cartesian product $`X\times Y`$ is here equipped with the direct product of the probability measures $`\mu `$ and $`\nu `$ , which is always $`T\times S`$-invariant. Acting on equivalence classes of isomorphic dynamical systems, the direct product operation is commutative, associative, and possesses a neutral element which is the trivial system reduced to one singleton. Translated in the context of stationary processes, if $`\xi `$ and $`\zeta `$ are respectively $`E`$-valued and $`F`$-valued stationary processes (not necessarily defined on the same probability space), their direct product $`\xi \zeta `$ is the $`E\times F`$-valued stationary process whose law is the direct product of the laws of $`\xi `$ and $`\zeta `$. That is to say, we are making the two processes $`\xi `$ and $`\zeta `$ live together in an independant (hence stationary) way.
Second, we say that $`(Y,,\nu ,S)`$ is a *factor* of $`(X,𝒜,\mu ,T)`$ if we can find a measurable map $`\pi :XY`$ satisfying
* $`\pi (\mu )=\nu `$ ;
* $`\pi T=S\pi `$.
Such a map is called a *homomorphism* of dynamical systems. Heuristically, the existence of such a homomorphism means that we can see the system $`S`$ inside the system $`T`$ by looking at $`\pi (x)`$. To this factor is associated the *factor $`\sigma `$-algebra* $`\pi ^1()`$. This $`\sigma `$-algebra has the property that if $`A`$ is measurable with respect to $`\pi ^1()`$, then so are $`T(A)`$ and $`T^1(A)`$. Conversely, each sub-$`\sigma `$-algebra $``$ of $`𝒜`$ which is invariant by $`T`$ and $`T^1`$ is a factor $`\sigma `$-algebra. Indeed, such a $`\sigma `$-algebra is always generated by some stationary process $`\xi `$ living in the system $`T`$ (this is an adaptation of Theorem 1.2 to the case of a factor $`\sigma `$-algebra), and the dynamical system $`T_\xi `$ associated to $`\xi `$ is a factor of $`T`$: Take $`\pi (x):=(\mathrm{},\xi _1(x),\xi _0(x),\xi _1(x),\mathrm{})`$.
Clearly, both systems $`S`$ and $`T`$ are factors of the direct product $`T\times S`$: These systems are seen by looking at the two coordinates in the Cartesian product. The analogy with the integers introduced by Furstenberg now gives two ways to define “$`T`$ and $`S`$ are relatively prime”.
###### Property 1.
$`T`$ and $`S`$ have no common factor except the trivial system reduced to one point.
###### Property 2.
Each time $`T`$ and $`S`$ appear as factors in some dynamical system, then their direct product $`T\times S`$ also appears as a factor<sup>2</sup><sup>2</sup>2To be exact, Property 2 should be stated in the following, stronger form: Each time $`T`$ and $`S`$ appear as factors in some dynamical system through the respective homomorphisms $`\pi _T`$ and $`\pi _S`$, $`T\times S`$ also appears as a factor through a homomorphism $`\pi _{T\times S}`$ such that $`\pi _X\pi _{T\times S}=\pi _T`$ and $`\pi _Y\pi _{T\times S}=\pi _S`$, where $`\pi _X`$ and $`\pi _Y`$ are the projections on the coordinates in the Cartesian product $`X\times Y`$..
The relationships between these two properties are exposed in Section 3. Property 2 has been called by Furstenberg the *disjointness* property of $`T`$ and $`S`$. The study of disjointness turns out to be quite enlightening, as it refers to *all the possible ways that two systems can be seen living together inside a third one*. As explained in the next subsection, this is precisely the theory of joinings.
The large scale of applications of disjointness and joinings is already visible in Furstenberg’s article where the concepts of disjointness and joinings where introduced. Indeed, Furstenberg’s initial motivations ranged from the classification of dynamical systems (how disjointness can be used to characterized some classes of dynamical systems) to a question in Diophantine approximation (which multiplicative semigroups of integers $`S`$ have the property that any real number is the limit of a sequence of rational numbers with denominators belonging to $`S`$?), passing through a filtering problem in probability theory (if $`(\xi _n)`$ and $`(\zeta _n)`$ are two real-valued stationary processes, can we recover $`(\xi _n)`$ from the knowledge of $`(\xi _n+\zeta _n)`$?). Since Furstenberg’s article, disjointness and joinings have been widely studied, and many other related notions have been introduced. In the present work, we shall mainly concentrate on some links between joinings and other ergodic properties of dynamical systems. For a more complete treatment of ergodic theory via joinings, we refer the readers to Eli Glasner’s book .
### 1.3. Joinings
###### Definition 1.3.
Let $`(X,𝒜,\mu ,T)`$ and $`(Y,,\nu ,S)`$ be two dynamical systems. A *joining* of $`T`$ and $`S`$ is a probability measure $`\lambda `$ on the Cartesian product $`X\times Y`$, whose marginals on $`X`$ and $`Y`$ are $`\mu `$ and $`\nu `$ respectively, and which is invariant by the product transformation $`T\times S`$.
The set of joinings of $`T`$ and $`S`$, denoted by $`J(T,S)`$, is never empty since it always contains the product measure $`\mu \nu `$. Each $`\lambda J(T,S)`$ provides a new dynamical system $`(X\times Y,𝒜,\lambda ,T\times S)`$, which will be denoted by $`(T\times S)_\lambda `$ to avoid confusion with the direct product $`T\times S`$. Both $`T`$ and $`S`$ are factors of $`(T\times S)_\lambda `$. Conversely, if $`T`$ and $`S`$ appear as factors in a third system $`(Z,𝒞,\rho ,R)`$ through respectively the maps $`\pi _X:ZX`$ and $`\pi _Y:ZY`$, we can consider the map $`\pi :ZX\times Y`$ defined by $`\pi (z):=(\pi _X(z),\pi _Y(z))`$. Then the probability measure $`\lambda :=\pi (\rho )`$ on $`X\times Y`$ satisfies the conditions to be a joining of $`T`$ and $`S`$, and the system $`(T\times S)_\lambda `$ is a factor of $`R`$.
### Topology on the set of joinings
The set of joinings of $`T`$ and $`S`$ is equipped with the topology defined by
(1)
$$\lambda _n\underset{n\mathrm{}}{}\lambda A𝒜,B,\lambda _n(A\times B)\underset{n\mathrm{}}{}\lambda (A\times B).$$
This topology is metrizable: An example of distance defining it is given by
$$d(\lambda ,\lambda ^{}):=\underset{m,n0}{}\frac{1}{2^{m+n}}|\lambda (A_m\times B_n)\lambda ^{}(A_m\times B_n)|,$$
where $`(A_m)_{m0}`$ and $`(B_n)_{n0}`$ are countable algebras generating respectively the $`\sigma `$-algebras $`𝒜`$ and $``$. With this topology, the set of joinings $`J(T,S)`$ is turned into a *compact metrizable topological space*. It is interesting to observe that, when $`X`$ and $`Y`$ are compact metric spaces, this topology coincides with the weak topology on the set of joinings.
All these definitions can be extended in an obvious way to the notion of *joining of a finite or countable family of dynamical systems* $`(T_i)_{iI}`$. If all the $`T_i`$’s are copies of a single dynamical system $`T`$, we speak of *self-joinings* of $`T`$. The set of self-joinings of order $`k`$ of $`T`$ (joinings of $`k`$ copies of $`T`$) is denoted by $`J_k(T)`$, and we simply write $`J(T)`$ for $`J_2(T)`$.
### Ergodic joinings
We denote by $`J_e(T,S)`$ the set of *ergodic* joinings of $`T`$ and $`S`$. Since every factor of an ergodic system is still ergodic, a necessary condition for $`J_e(T,S)`$ not to be empty is that both $`T`$ and $`S`$ be ergodic. The following proposition shows that the converse is true.
###### Proposition 1.4.
Given two ergodic dynamical systems $`T`$ and $`S`$, there exists at least one ergodic joining of $`T`$ and $`S`$.
###### Proof.
We start from the only joining whose existence is known, $`\lambda :=\mu \nu `$. This joining may not be ergodic, but in this case we consider its *decomposition in ergodic components*: There exists a probability measure $`P`$ on the set of all $`T\times S`$-invariant ergodic probability measures such that
$$\lambda =\lambda _\omega 𝑑P(\omega ).$$
Denoting by $`\lambda _\omega ^1`$ (respectively $`\lambda _\omega ^2`$) the marginals of the ergodic component $`\lambda _\omega `$ on $`X`$ (respectively $`Y`$), we get
$$\mu =\lambda _\omega ^1𝑑P(\omega ),$$
and
$$\nu =\lambda _\omega ^2𝑑P(\omega ).$$
Since $`\mu `$ and $`\nu `$ are ergodic, this implies that for $`P`$-almost every $`\omega `$, $`\lambda _\omega ^1=\mu `$ and $`\lambda _\omega ^2=\nu `$. In other words, for $`P`$-almost every $`\omega `$, $`\lambda _\omega `$ is an ergodic joining of $`T`$ and $`S`$. ∎
Note that this proof implies a stronger result: When $`T`$ and $`S`$ are ergodic, every joining of $`T`$ and $`S`$ is a convex combination of ergodic joinings. Since an ergodic measure cannot be written as a non trivial convex combination of invariant measures, we get the following.
###### Proposition 1.5.
If $`T`$ and $`S`$ are ergodic, the set of ergodic joinings of $`T`$ and $`S`$ is the set of extremal points in the compact convex set $`J(T,S)`$.
## 2. From disjointness to isomorphism
In this section, we present the two extreme cases concerning the joinings of two dynamical systems $`T`$ and $`S`$. The first one is when the set of joinings $`J(T,S)`$ is reduced to the singleton $`\{\mu \nu \}`$. Heuristically, this means that $`T`$ and $`S`$ have nothing in common which could lead to a non-trivial joining (what they have in common when they are not disjoint is developped in Sections 3.3 and 3.4). To the opposite, if $`T`$ and $`S`$ are isomorphic, we can construct very special joinings of $`T`$ and $`S`$, namely the joinings supported on graphs of isomorphisms.
### 2.1. Product measure and disjointness
###### Definition 2.1.
The dynamical systems $`(X,𝒜,\mu ,T)`$ and $`(Y,,\nu ,S)`$ are *disjoint* if the product measure $`\mu \nu `$ is the only joining of $`T`$ and $`S`$. We write in this case $`TS`$.
We leave as an exercise to the reader the verification of the equivalence of this definition with the strong form of Property 2. By the way, why was it necessary to state this strong form? (Hint: There exists a non-trivial dynamical system $`T`$ which is isomorphic to the direct product $`T\times T`$.)
We now give the most elementary example of disjointness: The identity is disjoint from any ergodic dynamical system.
###### Proposition 2.2.
If $`T`$ is the identity map on $`X`$ and $`S`$ is any ergodic dynamical system, then $`TS`$.
###### Proof.
Let $`\lambda `$ be any joining of $`T=\text{Id}`$ and $`S`$. For any $`A𝒜`$, $`B`$, invariance of $`\lambda `$ by $`T\times S`$ gives
$`n1,\lambda (A\times B)`$ $`=`$ $`\lambda (A\times S^nB)`$
$`=`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{k=0}{\overset{n1}{}}}\lambda (A\times S^kB)`$
$`=`$ $`{\displaystyle _{A\times Y}}{\displaystyle \frac{1}{n}}{\displaystyle \underset{k=0}{\overset{n1}{}}}\mathrm{𝟙}_B(S^ky)d\lambda (x,y).`$
Since $`S`$ is ergodic, $`\frac{1}{n}_{k=0}^{n1}\mathrm{𝟙}_B(S^ky)`$ converges to $`\nu (B)`$ $`\nu `$-almost everywhere, hence $`\lambda `$-almost everywhere. This implies
$$\lambda (A\times B)=\mu (A)\nu (B),$$
hence $`\lambda =\mu \nu `$. ∎
Another useful way to state the preceding result is the following:
###### Proposition 2.3.
Let $`T`$ and $`S`$ be two dynamical systems, with $`S`$ ergodic, and $`\lambda J(T,S)`$. If $`\lambda `$ is invariant by $`\text{Id}\times S`$, then $`\lambda =\mu \nu `$.
#### 2.1.1. Application: Furstenberg’s multiple recurrence theorem for weakly mixing transformations
As a nice application of Proposition 2.3, we give below an elegant proof of Furstenberg’s multiple recurrence theorem in the (easy) case of weakly mixing transformations . This proof was given by V. Ryzhikov in . Recall that $`T`$ is weakly mixing if $`T\times T`$ is ergodic, and that this implies
* $`T\times S`$ is ergodic for any ergodic system $`S`$ (see also subsection 4.1);
* $`T^k`$ is weakly mixing for any $`k0`$.
###### Theorem 2.4.
Let $`T`$ be a weakly mixing dynamical system. For any integer $`k1`$, and any measurable subsets $`A_0,\mathrm{},A_k`$, we have
(2)
$$\frac{1}{n}\underset{j=1}{\overset{n}{}}\mu (A_0T^jA_1T^{2j}A_2\mathrm{}T^{kj}A_k)\underset{n\mathrm{}}{}\mu (A_0)\mu (A_1)\mathrm{}\mu (A_k).$$
###### Proof.
Consider for any $`n1`$ the self-joining $`\lambda _n`$ of $`T`$ of order $`k+1`$ defined by
$$\lambda _n(A_0\times A_1\times \mathrm{}\times A_k):=\frac{1}{n}\underset{j=1}{\overset{n}{}}\mu (A_0T^jA_1T^{2j}A_2\mathrm{}T^{kj}A_k).$$
Note that the result we want to prove is equivalent to the following translation in the language of joinings:
$$\lambda _n\underset{n\mathrm{}}{}\mu \mu \mathrm{}\mu .$$
By compacity, it is enough to verify that if $`\lambda `$ is any cluster point of the sequence $`(\lambda _n)`$, then $`\lambda `$ is the product measure, which we shall prove by induction on $`k`$. Observe that such a cluster point is always $`\text{Id}\times T\times T^2\times \mathrm{}\times T^k`$-invariant. Thus Proposition 2.3 gives the result when $`k=1`$ by ergodicity of $`T`$. Now take $`k>1`$ and suppose the result is true for $`k1`$. This induction hypothesis ensures that the marginal of $`\lambda `$ on the last $`k`$ coordinates is the product measure. Hence $`\lambda `$ can be seen as a joining of $`\text{Id}`$ with $`T\times T^2\times \mathrm{}\times T^k`$. Since $`T\times T^2\times \mathrm{}\times T^k`$ is ergodic (because $`T`$ is weakly mixing), Proposition 2.3 shows that the result is still true for $`k`$. ∎
#### 2.1.2. Disjointness and classification of systems
As we have already mentioned, one of the reasons for the introduction of disjointness by Furstenberg was the classification of classes of dynamical systems. And indeed there are many results in this direction: For example, Furstenberg proved that a dynamical system has zero entropy if and only if it is disjoint from any Bernoulli dynamical system. It also follows from his work that $`T`$ is weakly mixing if and only if it is disjoint from any rotation on the circle. Some further explanations on how disjointness can characterize such classes of dynamical systems will be given in Section 3.3.2 (see in particular Theorem 3.6).
#### 2.1.3. Disjointness and pointwise convergence of ergodic averages
In was introduced the so-called *weak disjointness* property of dynamical systems:
###### Definition 2.5.
$`(X,𝒜,\mu ,T)`$ and $`(Y,,\nu ,S)`$ are said to be *weakly disjoint* if, given any function $`f`$ in $`L^2(\mu )`$ and any function $`g`$ in $`L^2(\nu )`$, there exist a set $`A`$ in $`𝒜`$ and a set $`B`$ in $``$ such that
* $`\mu (A)=\nu (B)=1`$
* for all $`xA`$ and for all $`yB`$, the sequence
(3)
$$\left(\frac{1}{N}\underset{n=0}{\overset{N1}{}}f\left(T^nx\right)g\left(S^ny\right)\right)_{N>0}$$
converges.
(In fact it is sufficient to check this statement for only dense families of $`f`$ and $`g`$ in their respective $`L^2`$ spaces.) The main motivation for the introduction of this property was the study of non-conventional ergodic averages: Suppose that $`S`$ and $`T`$ act on the same probability space $`X`$, and consider for $`f`$ and $`g`$ in $`L^2(\mu )`$ the averages
$$\frac{1}{n}\underset{k=0}{\overset{n1}{}}f(T^kx)g(S^kx).$$
When do these averages converge $`\mu `$-a.e.? If we assume that $`T`$ and $`S`$ are weakly disjoint, then almost-everywhere pointwise convergence holds.
As we could guess from the name of the property, if $`T`$ and $`S`$ are disjoint, they are weakly disjoint. The link between joinings and convergence of ergodic sums like in (3) comes from the following remark: We can always assume that $`T`$ and $`S`$ are continuous transformations of compact metric spaces $`X`$ and $`Y`$ respectively (indeed, any measurable system is isomorphic to such a transformation on a compact metric space: see e.g. ). Then the set of probability measures on $`X\times Y`$ equipped with the topology of weak convergence is metric compact. If moreover $`T`$ and $`S`$ are ergodic, we can easily find subsets $`X_0X`$ and $`Y_0Y`$ with $`\mu (X_0)=\nu (Y_0)=1`$, such that for all $`(x,y)X_0\times Y_0`$, any cluster point of the sequence
$$\delta _n(x,y)=\frac{1}{n}\underset{k=0}{\overset{n1}{}}\delta _{(T^kx,S^ky)}$$
is automatically a joining of $`T`$ and $`S`$. When $`T`$ and $`S`$ are disjoint, there is therefore only one possible cluster point to the sequence $`\delta _n(x,y)`$ which is $`\mu \nu `$. This ensures that for continuous functions $`f`$ and $`g`$, convergence of (3) to the product of the integrals of $`f`$ and $`g`$ holds, hence $`T`$ and $`S`$ are weakly disjoint.
It is not difficult to show that weak disjointness is really weaker than disjointness. Indeed, there are many examples of dynamical systems which are self-weakly disjoint (weakly disjoint from themselves): For example, irrational rotations and Chacon’s transformation are self-weakly disjoint (see for details and other examples). However, a non trivial dynamical system is never disjoint from itself, as we will see in the next paragraph.
### 2.2. Joinings supported on graphs
Suppose that $`S`$ is a factor of $`T`$, with $`\pi :XY`$ a homomorphism. Then $`\pi `$ gives rise to the existence of a very special joining of $`T`$ and $`S`$, which we denote by $`\mathrm{\Delta }_\pi `$, and which is defined by
$$\mathrm{\Delta }_\pi (A\times B):=\mu (A\pi ^1B).$$
###### Lemma 2.6.
The joining $`\mathrm{\Delta }_\pi `$ has the following two properties:
* $`\mathrm{\Delta }_\pi (G)=1`$, where $`G:=\{(x,\pi (x)),xX\}`$ is the graph of $`\pi `$;
* $`\{X,\mathrm{}\}𝒜\{Y,\mathrm{}\}mod\mathrm{\Delta }_\pi `$, where the notation ‘$`𝒞𝒟mod\lambda `$’ means that for each $`C`$ in the $`\sigma `$-algebra $`𝒞`$, there exists a $`D`$ in the $`\sigma `$-algebra $`𝒟`$ such that $`\lambda (C\mathrm{\Delta }D)=0`$.
###### Proof.
Take a refining and generating sequence $`(𝒫_n)`$ of finite measurable partitions in $`(Y,,\nu )`$. (That is to say, $`𝒫_{n+1}`$ is finer than $`𝒫_n`$, and if $`y`$ and $`y^{}`$ are distinct points in $`Y`$, there exists some $`n`$ for which $`y`$ and $`y^{}`$ do not belong to the same atom of $`𝒫_n`$; this ensures in particular that the $`\sigma `$-algebra generated by the partitions $`𝒫_n`$ is $``$.) Then $`G`$ can be written as
$$G=\underset{n}{}\underset{B𝒫_n}{}\pi ^1(B)\times B.$$
But for any $`n`$, it is clear from the definition of $`\mathrm{\Delta }_\pi `$ that
$$\mathrm{\Delta }_\pi \left(\underset{B𝒫_n}{}\pi ^1(B)\times B\right)=1,$$
which proves that $`\mathrm{\Delta }_\pi `$ is supported on the graph of $`\pi `$. For the second property, it is easy to verify that $`\mathrm{\Delta }_\pi `$ identifies any $`X\times B`$, where $`B`$, with $`\pi ^1B\times Y`$. ∎
It is quite remarkable to see that the converse is true, i.e. we can characterize the fact that $`S`$ is a factor of $`T`$ by the existence of some special joining.
###### Proposition 2.7.
Suppose that there exists a joining $`\lambda J(T,S)`$ such that
$$\{X,\mathrm{}\}𝒜\{Y,\mathrm{}\}mod\lambda .$$
Then $`S`$ is a factor of $`T`$ and we can find a homomorphism $`\pi :XY`$ such that $`\lambda `$ is supported on the graph of $`\pi `$.
###### Proof.
Take again a refining and generating sequence $`(𝒫_n)`$ of finite measurable partitions in $`(Y,,\nu )`$. For each $`n`$, $`𝒫_n`$ can be identified via $`\lambda `$ with some $`\overline{𝒫_n}`$ in $`X`$, where $`(\overline{𝒫_n})`$ is a refining, but not necessarily generating, sequence of partitions in $`(X,\mu )`$, Now, for $`\mu `$-almost every $`xX`$, if we denote by $`\overline{B_n}`$ the atom of $`\overline{𝒫_n}`$ to which $`x`$ belongs and $`B_n`$ the atom of $`𝒫_n`$ identified via $`\lambda `$ with $`\overline{B_n}`$, then $`_nB_n`$ is a singleton. Denote by $`\pi (x)`$ the only element of this set; we let the reader check that $`x\pi (x)`$ satisfies the required properties. ∎
If we assume further that $`T`$ and $`S`$ are *isomorphic* (in other words, the homomorphism $`\pi `$ from $`X`$ to $`Y`$ is invertible and $`\pi ^1`$ is also a homomorphism of dynamical systems; in this case we will rather use the symbol $`\varphi `$ instead of $`\pi `$), the joining $`\mathrm{\Delta }_\varphi `$ now satisfies $`\{X,\mathrm{}\}=𝒜\{Y,\mathrm{}\}mod\mathrm{\Delta }_\varphi `$, where the notation ‘$`𝒞=𝒟mod\lambda `$’ means that we have both $`𝒞𝒟mod\lambda `$ and $`𝒟𝒞mod\lambda `$. The analog of Proposition 2.7 is also true, and we summarize all these results in the following theorem.
###### Theorem 2.8.
Let $`T`$ and $`S`$ be two dynamical systems. $`S`$ is a factor of $`T`$ if and only if there exists some joining $`\lambda J(T,S)`$ such that $`\{X,\mathrm{}\}𝒜\{Y,\mathrm{}\}mod\lambda `$; in this case $`\lambda `$ is supported on the graph of some homomorphism. $`T`$ and $`S`$ are isomorphic if and only if there exists some joining $`\lambda J(T,S)`$ such that $`\{X,\mathrm{}\}=𝒜\{Y,\mathrm{}\}mod\lambda `$; in this case $`\lambda `$ is supported on the graph of some isomorphism.
#### 2.2.1. Application: isomorphism of discrete-spectrum transformations
There are very deep applications of the characterization of isomorphic systems by joinings. Krieger’s theorem, stating that any dynamical system with entropy less than $`\mathrm{log}_2(k)`$ is isomorphic to the shift map on $`\{1,\mathrm{},k\}`$ with an appropriate invariant measure, and Ornstein’s theorem characterizing the dynamical systems which are isomorphic to some Bernoulli shift, can both be proved using this result (see , Chapter 7). Here we illustrate the power of Theorem 2.8 by giving Lemańczyk’s elegant proof of a well-known theorem due to Halmos and von Neumann. This proof is taken from . Recall that a system $`T`$ is said to have *discrete spectrum* if the subspace of $`L^2`$ spanned by the eigenvectors of $`ffT`$ is dense in $`L^2`$.
###### Theorem 2.9.
If $`T`$ is ergodic and has discrete spectrum, and if $`S`$ is spectrally isomorphic to $`T`$, then $`T`$ and $`S`$ are isomorphic.
###### Proof.
First note that $`S`$ is ergodic and also has discrete spectrum with the same eigenvalues as $`T`$, since it is spectrally isomorphic to $`T`$. Let $`\lambda `$ be an ergodic joining of $`T`$ and $`S`$. For an eigenvalue $`\alpha `$ of $`T`$, let $`f`$ be an eigenvector associated to $`\alpha `$ in the system $`T`$, and $`g`$ associated to $`\alpha `$ in the system $`S`$. Then both functions $`(x,y)f(x)`$ and $`(x,y)g(y)`$ are eigenvectors in $`(T\times S)_\lambda `$. Since $`(T\times S)_\lambda `$ is an ergodic system, $`\alpha `$ is a simple eigenvalue in $`(T\times S)_\lambda `$, hence $`(x,y)g(y)`$ and $`(x,y)f(x)`$ are proportional. This implies that $`(x,y)f(x)`$ belongs to $`L^2(\{X,\mathrm{}\},\lambda )`$; but since eigenvectors are dense in $`L^2(𝒜,\mu )`$, this in turn implies $`L^2(𝒜\{Y,\mathrm{}\},\lambda )L^2(\{X,\mathrm{}\},\lambda )`$. For the same reason, the inverse inclusion is also true, which proves that $`\{X,\mathrm{}\}=𝒜\{Y,\mathrm{}\}mod\lambda `$, and the two systems are isomorphic. ∎
#### 2.2.2. Self-joinings, commutant and minimal self-joinings
When studying self-joinings of a single dynamical system $`T`$, we are obviously in a situation where the two systems that we want to join together are isomorphic. Now, thinking of what an isomorphism between $`T`$ and itself is, we see that this has to be an automorphism of the Lebesgue space $`(X,𝒜,\mu )`$, which commutes with $`T`$. We call *commutant of $`T`$*, and denote by $`C(T)`$, the set
$$C(T):=\{S\text{Aut}(X,𝒜,\mu ):ST=TS\}.$$
For each $`SC(T)`$, let us denote by $`\mathrm{\Delta }_S`$ the self-joining of $`T`$ defined by
$$\mathrm{\Delta }_S(A\times B):=\mu (AS^1B).$$
This leads to an identification of $`C(T)`$ with some subset of $`J_2(T)`$, namely the subset of all $`\lambda J_2(T)`$ such that
$$𝒜\{X,\mathrm{}\}=\{X,\mathrm{}\}𝒜mod\lambda .$$
Observe that for each $`T`$ the commutant of $`T`$ contains at least the powers $`T^n`$, $`n`$. Besides, note that each joined system $`(T\times T)_{\mathrm{\Delta }_S}`$ for $`SC(T)`$ is isomorphic to $`T`$ (an isomorphism is given by $`x(x,Sx)`$). In particular, $`(T\times T)_{\mathrm{\Delta }_S}`$ is ergodic as soon as $`T`$ is ergodic. This gives for each ergodic $`T`$ a family of “obvious” ergodic self-joinings: the $`\mathrm{\Delta }_{T^n}`$, $`n`$. In the case where $`T`$ is weakly mixing, the product measure $`\mu \mu `$ is also considered as an obvious ergodic self-joining. In 1979, Rudolph introduced in the important notion of “minimal self joinings”, which says that $`T`$ has no other ergodic self-joinings of order 2 than the obvious ones. The notion of “$`k`$-fold minimal self-joinings” refers to joinings of order $`k`$:
###### Definition 2.10.
For $`k2`$, we say that an ergodic $`T`$ has *$`k`$-fold minimal self-joinings* if for each ergodic joining $`\lambda `$ of $`k`$ copies of $`T`$, we can partition the set $`\{1,\mathrm{},k\}`$ of coordinates into two subsets $`J_1,\mathrm{},J_{\mathrm{}}`$ such that
1. for $`j_1`$ and $`j_2`$ belonging to the same $`J_i`$, the marginal of $`\lambda `$ on the coordinates $`j_1`$ and $`j_2`$ is some $`\mathrm{\Delta }_{T^n}`$;
2. for $`j_1J_1,\mathrm{},j_{\mathrm{}}J_{\mathrm{}}`$, the coordinates $`j_1\mathrm{},j_{\mathrm{}}`$ are independent.
We say that $`T`$ has minimal self-joinings (MSJ) if $`T`$ has $`k`$-fold minimal self-joinings for every $`k2`$.
In fact, it has been proved by Glasner, Host and Rudolph that for a weakly mixing $`T`$, 3-fold minimal self-joinings implies $`k`$-fold minimal self-joinings for all $`k`$. (See also a proof in .) A very important question which remains open in the area is the following:
###### Question 2.11.
Does there exist a transformation $`T`$ which has 2-fold but not 3-fold minimal self-joinings?
(This question is closely related to the problem of 2-fold and 3-fold mixing: See Section 4.3.)
There are many examples of dynamical systems with the MSJ property, the simplest of which is probably the well known Chacon’s transformation (see ). One immediate consequence of this property is that $`C(T)`$ is reduced to the powers of $`T`$. This in turn implies that the transformation $`T`$ has no square root (i.e. there is no automorphism $`S`$ such that $`SS=T`$). It can be noticed that the first example of a transformation with no square root, given by Ornstein in 1970, belongs to the class of *mixing rank-one<sup>3</sup><sup>3</sup>3We shall not define in this text what a rank-one system is; the reader unfamiliar with this notion can find a good introduction in systems*, which all turned out to have the MSJ property, as King proved in 1988 . Another joining proof of this result has been given by Ryzhikov .
If, in condition (1) of the definition of MSJ, we replace $`\mathrm{\Delta }_{T^n}`$ by $`\mathrm{\Delta }_S`$, for some $`SC(T)`$, we get the weaker, but important, notion of *simplicity*. In the case of joinings of order 2, this property was introduced by Veech , who proved a useful result on the structure of factor $`\sigma `$-algebras of 2-simple systems. The general definition of simplicity was later given by Del Junco and Rudolph . We also refer to Thouvenot for a presentation of properties of simple systems.
As a nice application of joinings to the study of $`C(T)`$, we should also mention here Ryzhikov’s proof of King’s so-called *Weak Closure Theorem*, stating that if $`T`$ is rank one, any $`S`$ in $`C(T)`$ is the limit of some sequence $`T^{n_k}`$ of the powers of $`T`$ .
## 3. Joinings and factors
In this section we discuss the links and the differences between the two properties proposed in the introduction to define “$`T`$ and $`S`$ are relatively prime”. The starting point of the theory is the existence of a special joining arising when $`T`$ and $`S`$ have a common factor.
### 3.1. Relatively independent joining above a common factor
We first define this special joining in the particular case of two systems which obviously have a common factor, since they both arise as a joining of a fixed system $`R`$ with other systems.
Let $`(X,𝒜,\mu ,T)`$, $`(Y,,\nu ,S)`$ and $`(Z,𝒞,\rho ,R)`$ be three dynamical systems. Suppose that $`\lambda _T`$ is a joining of $`T`$ and $`R`$, and $`\lambda _S`$ is a joining of $`S`$ and $`R`$. Then we can put all these systems together so that the marginal on $`(X\times Z)`$ is $`\lambda _T`$ and the marginal on $`(Y\times Z)`$ is $`\lambda _S`$: First pick $`z`$ according to the probability law $`\rho `$, then pick $`x`$ and $`y`$ according to their conditionnal law knowing $`z`$ in the respective joinings $`\lambda _T`$ and $`\lambda _S`$, but independently of each other. More precisely, consider the 3-fold joining of $`T`$, $`S`$ and $`R`$ denoted by $`\lambda _T_R\lambda _S`$ defined by setting, for all $`A𝒜`$, $`B`$ and $`C𝒞`$
(4)
$$\lambda _T_R\lambda _S(A\times B\times C):=_C𝔼_{\lambda _T}[\mathrm{𝟙}_{xA}|z]𝔼_{\lambda _S}[\mathrm{𝟙}_{yB}|z]𝑑\rho (z).$$
We get in this way a joining of the two systems $`(T\times R)_{\lambda _T}`$ and $`(S\times R)_{\lambda _S}`$ identifying the coordinate $`z`$ on both systems. This joining is called the *relatively independent joining of $`(T\times R)_{\lambda _T}`$ and $`(S\times R)_{\lambda _S}`$ above their common factor $`R`$*.
Now turn to the more general situation of two systems $`T`$ and $`S`$ sharing a common factor $`R`$. This means that we have homomorphisms of dynamical systems $`\pi _T:XZ`$ and $`\pi _S:YZ`$, which give us special joinings $`\lambda _{\pi _T}J(T,R)`$ and $`\lambda _{\pi _S}J(S,R)`$ (see section 2.2). We can then consider the relatively independent joining $`\lambda _{\pi _T}_R\lambda _{\pi _S}`$, which is a 3-fold joining of $`T`$, $`S`$ and $`R`$. The marginal of $`\lambda _{\pi _T}_R\lambda _{\pi _S}`$ on $`X\times Y`$ is now a joining of $`T`$ and $`S`$, called the *relatively independent joining of $`T`$ and $`S`$ above their common factor $`R`$*, and denoted by $`\mu _R\nu `$. This joining has the strong property of identifying the projections of both systems on the common factor.
In summary, the $`\sigma `$-algebra $`𝒞`$ can be viewed as a sub-$`\sigma `$-algebra of both $`𝒜`$ and $``$, any $`𝒞`$-measurable function can be viewed as a function defined on $`Z`$, and we have the following integration formula:
$$_{X\times Y}f(x)g(y)d(\mu _R\nu )(x,y)=_Z𝔼[f(x)|𝒞](z)𝔼[g(y)|𝒞](z)𝑑\rho (z).$$
###### Proposition 3.1.
Let $`\mu _R\nu `$ be the relatively independent joining of $`T`$ and $`S`$ above their common factor $`R`$. Then
$$\pi _T(x)=\pi _S(y)\mu _R\nu \text{-a.e.}$$
###### Proof.
By construction, $`\mu _R\nu `$ is the marginal on $`X\times Y`$ of the 3-fold joining $`\lambda _{\pi _T}_R\lambda _{\pi _S}`$ of $`T`$, $`S`$ and $`R`$ whose marginal on $`X\times Z`$ (respectively $`Y\times Z`$) is $`\lambda _{\pi _T}`$ (respectively $`\lambda _{\pi _S}`$). Recall that, from Lemma 2.6, we have $`z=\pi _T(x)`$ $`\lambda _{\pi _T}`$-a.e. and $`z=\pi _S(y)`$ $`\lambda _{\pi _S}`$-a.e. Therefore, $`z=\pi _T(x)=\pi _S(y)`$ $`\lambda _{\pi _T}_R\lambda _{\pi _S}`$-a.e., hence $`\pi _T(x)=\pi _S(y)`$ $`\mu _R\nu `$-a.e. ∎
Note that what we did here with two systems could have been done with any finite or countable family of systems sharing a common factor.
### 3.2. Self-joinings and factors
Any factor $`S`$ of $`T`$ can be considered as a common factor of two copies of $`T`$, thus gives rise to a special self-joining of $`T`$: The relatively independent self-joining of $`T`$ above its factor $`S`$. As in ordinary arithmetic of integer numbers, every dynamical system which is not reduced to one point has at least two factors: $`T`$ itself is of course a factor of $`T`$, and the system reduced to one point is also a factor of $`T`$. Note that the relatively independent joining of $`T`$ above $`T`$ itself is the self-joining $`\mathrm{\Delta }_{\text{Id}}`$ of $`T`$ supported on the graph of $`\text{Id}`$, and that the relatively independent joining of $`T`$ above the factor reduced to one point is the product measure $`\mu \mu `$. Now, we want to show that any other factor of $`T`$ is incompatible with the 2-fold MSJ property.
###### Proposition 3.2.
Let $`(X,𝒜,\mu ,T)`$ be an ergodic dynamical system where $`X`$ is not finite. If $`T`$ has another non-trivial factor than $`T`$ itself, then $`T`$ does not have the 2-fold MSJ property.
###### Proof.
Let $`S`$ be a non-trivial factor of $`T`$, and consider $`\mu _S\mu `$ the relatively independent self-joining of $`T`$ above $`S`$ constructed from the homomorphism $`\pi :XY`$. The system defined by this joining may not be ergodic, but if we assume further that $`T`$ has the 2-fold MSJ property, the ergodic decomposition of $`\mu _S\mu `$ has the form
$$\mu _S\mu =\theta \mu \mu +(1\theta )\underset{n}{}p_n\mathrm{\Delta }_{T^n},$$
where $`0\theta 1`$, $`p_n0`$ and $`_np_n=1`$. Denote by $`x`$ and $`x^{}`$ the two coordinates in the Cartesian square of $`X`$, and remember that, by Proposition 3.1, we have $`\pi (x)=\pi (x^{})`$ $`\mu _S\mu `$-a.e. Since $`S`$ is not the factor reduced to one point, this would not be possible if $`\mu \mu `$ appeared in the ergodic decomposition, hence $`\theta =0`$. Moreover, if the ergodic decomposition of $`\mu _S\mu `$ was reduced to $`\mathrm{\Delta }_{\text{Id}}`$ , $`S`$ would be $`T`$ itself. Hence there is some $`n0`$ such that $`p_n>0`$. In other words, there exists some $`n0`$ such that, under $`\mu _S\mu `$,
(5)
$$\mu _S\mu (x^{}=T^nx)>0.$$
This implies that, in the factor $`S`$, $`y=S^ny`$ with positive probability. But $`S`$ is ergodic (because $`T`$ is), hence the space $`Y`$ on which $`S`$ acts is finite. In turn, this implies that there exists at least one $`yY`$ such that
(6)
$$\mu _S\mu \left((x^{}=T^nx)(\pi (x)=y)\right)>0.$$
Now, observe that $`\mu _S\mu `$ conditioned on $`(\pi (x)=y)`$ is the direct product of $`\mu `$ conditioned on $`(\pi (x)=y)`$ with itself. Thus, (6) implies that $`\mu `$ conditioned on $`(\pi (x)=y)`$ has at least one atom. But $`\mu (\pi (x)=y)>0`$, hence $`\mu `$ itself has an atom, hence $`T`$ is periodic. ∎
###### Corollary 3.3.
If $`T`$ is ergodic aperiodic with the 2-fold MSJ property, then any non constant stationary process living in the system $`T`$ generates the whole $`\sigma `$-algebra.
### 3.3. Disjointness and lack of common factor
We now turn back to the case of two dynamical systems $`T`$ and $`S`$ sharing a common factor $`R`$. As we have already noticed in the case of self-joinings, Proposition 3.1 ensures that if $`R`$ is not reduced to one point, the relatively independent joining of $`T`$ and $`S`$ above $`R`$ is not the product measure $`\mu \nu `$, hence $`T`$ and $`S`$ are not disjoint. This means that Property 2 (the disjointness property) implies Property 1 (the lack of common factor). It had been asked by Furstenberg whether these two properties were, in fact equivalent. The discovery by Rudolph of systems with the MSJ property allowed him to answer negatively to this question .
###### Proposition 3.4.
There exist two non disjoint systems $`S`$ and $`T`$ without any other common factor than the trivial one.
###### Proof.
Let $`\xi =(\xi _k)`$ be a stationary process generating a system $`T`$ with the MSJ property, and let $`(\xi _k^{})`$ be an independent copy of $`(\xi _k)`$. We now define another stationary process $`\zeta `$ by setting
$$\zeta _k:=\{\xi _k,\xi _k^{}\}.$$
(This means that $`\zeta _k`$ tells us which are the two values taken by $`\xi _k`$ and $`\xi _k^{}`$, but neither which one is taken by $`\xi _k`$ nor which one is taken by $`\xi _k^{}`$.) Denote by $`S`$ the system generated by $`\zeta `$: $`S`$ is a factor of the Cartesian product $`T\times T`$. Then $`T`$ and $`S`$ are certainly not disjoint, since the two processes $`\xi `$ and $`\zeta `$ are clearly not independent of each other. However, suppose that $`S`$ and $`T`$ share a common factor which is not the system reduced to one point. Then, since $`T`$ has the MSJ property, this factor has to be $`T`$ itself. Therefore, we can find a third copy $`\xi ^{\prime \prime }`$ of $`\xi `$, living in the Cartesian square $`T\times T`$, and which is measurable with respect to the factor $`\sigma `$-algebra generated by $`\zeta `$. By the 3-fold MSJ property of $`T`$, either $`\xi ^{\prime \prime }`$ is equal to one of the processes $`\xi `$ or $`\xi ^{}`$, possibly shifted by some constant integer $`k`$, or $`\xi ^{\prime \prime }`$ is independent of $`(\xi ,\xi ^{})`$. The former case is impossible since $`\zeta `$, hence $`\xi ^{\prime \prime }`$, is invariant by the “flip” $`(x,x^{})(x^{},x)`$, which is clearly not true for any shifted copy of $`\xi `$ or $`\xi ^{}`$. The latter case is also impossible since $`(\xi )`$ and $`(\xi ^{})`$ generate the whole $`\sigma `$-algebra in $`T\times T`$. ∎
#### 3.3.1. The fundamental lemma of non-disjointness
There exists however an important result which tells us how to derive some information on factors from the non-disjointness of two systems. This theorem appeared for the first time in year 2000, in two publications .
###### Theorem 3.5.
If $`T`$ and $`S`$ are not disjoint, then $`S`$ has a non trivial common factor with some joining of a countable family of copies of $`T`$.
###### Proof.
Start with a joining $`\lambda `$ of $`(X,𝒜,\mu ,T)`$ and $`(Y,,\nu ,S)`$: Then $`S`$ is a factor of $`(T\times S)_\lambda `$. Consider a countable family of copies of $`(T\times S)_\lambda `$, and denote by $`\lambda _{\mathrm{}}`$ their relatively independent joining over their common factor $`S`$. Since this joining identifies all the projections on $`Y`$, it can be viewed as a probability measure on $`Y\times X^{}`$, and is characterized by
$$\lambda _{\mathrm{}}(B\times A_0\times \mathrm{}\times A_k\times X\times X\mathrm{})=_B𝔼_\lambda [\mathrm{𝟙}_{A_0}|y]\mathrm{}𝔼_\lambda [\mathrm{𝟙}_{A_k}|y]𝑑\nu (y).$$
Observe that this probability measure is invariant under the shift on each $`y`$-fiber:
$$(y,x_0,x_1,\mathrm{})(y,x_1,x_2\mathrm{}).$$
Moreover, $`\lambda _{\mathrm{}}`$ conditioned on such a fiber is a product measure: The infinite direct product of $`\lambda `$ conditioned on $`y`$. A relative version of Kolmogorov 0-1 law (precisely stated in , Lemma 9) then tells us that the $`\sigma `$-algebra of shift-invariant events coincides modulo $`\lambda _{\mathrm{}}`$ with the $`\sigma `$-algebra of $`y`$-measurable events.
Now, let $`f`$ be a bounded measurable function defined on $`X`$, which we also consider as a function on $`X\times Y`$ and on $`Y\times X^{}`$ by setting $`f(x,y):=f(x)`$, and $`f(y,x_0,x_1,\mathrm{}):=f(x_0)`$. Applying Birkhoff ergodic theorem to $`f`$ in the system defined by the shift and $`\lambda _{\mathrm{}}`$ gives
(7)
$$\underset{n\mathrm{}}{lim}\frac{1}{n}\underset{k=0}{\overset{n1}{}}f(x_k)=𝔼_\lambda _{\mathrm{}}[f|y]=𝔼_\lambda [f|y]\lambda _{\mathrm{}}\text{-a.e.}$$
Since $`T`$ and $`S`$ are not disjoint, we could have chosen $`\lambda \mu \nu `$, and then $`f`$ such that $`𝔼_\lambda [f|y]`$ is not constant modulo $`\lambda `$. Then this conditional expectation generates in $`S`$ a non-trivial factor, which by (7) can be identified to some factor of the joining of a countable family of copies of $`T`$ defined by the marginal of $`\lambda _{\mathrm{}}`$ on $`X^{}`$. ∎
#### 3.3.2. Application: disjointness of classes of dynamical systems
We get from Theorem 3.5 some important disjointness results. Note that some properties of dynamical systems are stable under the operations of taking joinings and factors. We will call these properties *stable properties*. This is e.g. the case of the zero-entropy property: We know that a factor of a zero-entropy system still has zero entropy, and that any joining of zero-entropy systems also has zero entropy. Take now any $`K`$-system $`S`$: We know that any non-trivial factor of $`S`$ has positive entropy, hence $`S`$ cannot have a non-trivial common factor with a joining of copies of a zero-entropy $`T`$.
The same argument also applies to the disjointness of discrete-spectrum systems with weakly mixing systems, since discrete spectrum is a stable property, and weakly mixing systems are characterized by the fact that they do not have any discrete-spectrum factor. We thus have the following theorem:
###### Theorem 3.6.
Any zero-entropy system is disjoint from any $`K`$-system. Any discrete-spectrum system is disjoint from any weakly mixing system.
### 3.4. Joinings and $`T`$-factors
The result of Theorem 3.5 leads to the introduction of a special class of factors when some dynamical system $`T`$ is given: For any other dynamical system $`S`$, call *$`T`$-factor* of $`S`$ any common factor of $`S`$ with a joining of countably many copies of $`T`$. We also call *$`T`$-factor $`\sigma `$-algebra* any factor $`\sigma `$-algebra of $`S`$ associated to a $`T`$-factor of $`S`$. Another way to state Theorem 3.5 is then the following: If $`S`$ and $`T`$ are not disjoint, then $`S`$ has a non-trivial $`T`$-factor. In fact, the proof of this theorem gives a more precise result: For any joining $`\lambda `$ of $`S`$ and $`T`$, for any bounded measurable function $`f`$ on $`X`$, the factor $`\sigma `$-algebra of $`S`$ generated by the function $`𝔼_\lambda [f(x)|y]`$ is a $`T`$-factor $`\sigma `$-algebra of $`S`$.
With the notion of $`T`$-factor, we can extend Theorem 3.5 in the following way, showing the existence of a special $`T`$-factor $`\sigma `$-algebra of $`S`$ concentrating anything in $`S`$ which could lead to a non trivial joining between $`T`$ and $`S`$.
###### Theorem 3.7.
Given two dynamical systems $`(X,𝒜,\mu ,T)`$ and $`(Y,,\nu ,S)`$, there always exists a maximum $`T`$-factor $`\sigma `$-algebra of $`S`$, denoted by $`_T`$.
Under any joining $`\lambda `$ of $`T`$ and $`S`$, the $`\sigma `$-algebras $`𝒜\{\mathrm{},Y\}`$ and $`\{\mathrm{},X\}`$ are independent conditionally to the $`\sigma `$-algebra $`\{\mathrm{},X\}_T`$.
The proof of the theorem is based on the two following lemmas.
###### Lemma 3.8.
Let $`(B_i)_{iI}`$ be a countable family of events in $``$, such that for all $`i`$, $`B_i`$ belongs to some $`T`$-factor $`\sigma `$-algebra $`_i`$ of $`S`$. Then there exists a $`T`$-factor $`\sigma `$-algebra of $`S`$ containing all the $`B_i`$’s.
###### Proof.
For each $`iI`$, we have a joining $`\lambda _i`$ of $`S`$ with a countable family $`(T_{i,n})_n`$ of copies of $`T`$, such that $`_i_n𝒜_{i,n}mod\lambda _i`$. Let us denote by $`R_i`$ the dynamical system defined by $`\lambda _i`$, and by $`\lambda `$ the relatively independent joining of all the $`R_i`$, $`iI`$, over their common factor $`S`$. We can see $`\lambda `$ as a joining of $`S`$ with the countable family $`(T_{i,n})_{(i,n)I\times }`$ and, for each $`i`$ we have
$$_i\underset{(i,n)I\times }{}𝒜_{i,n}mod\lambda .$$
We conclude that the factor $`\sigma `$-algebra of $`S`$ generated by all the $`_i`$’s is a $`T`$-factor $`\sigma `$-algebra, which certainly contains all the $`B_i`$’s. ∎
###### Lemma 3.9.
Let $``$ be a factor $`\sigma `$-algebra of $`S`$. If there exists a joining $`\lambda `$ of $`T`$ and $`S`$ under which the $`\sigma `$-algebras $`𝒜\{\mathrm{},Y\}`$ and $`\{\mathrm{},X\}`$ are not independent conditionally to $`\{\mathrm{},X\}`$, then there exists a $`T`$-factor $`\sigma `$-algebra $`^{}`$ of $`S`$, not contained in $``$.
###### Proof.
The hypothesis of the lemma implies the existence of a bounded measurable function $`f`$ on $`X`$ such that, on a set of positive $`\nu `$-measure,
$$𝔼_\lambda [f(x)|]𝔼_\lambda [f(x)|].$$
The factor $`\sigma `$-algebra $`^{}`$ of $`S`$ generated by the function $`𝔼_\lambda [f(x)|]`$ is not contained in $``$; but we saw in the proof of Theorem 3.5 that $`^{}`$ is a $`T`$-factor $`\sigma `$-algebra. ∎
###### Proof of Theorem 3.7.
In order to prove the existence of a maximum $`T`$-factor $`\sigma `$-algebra, we define
$$_T:=\{B:B\text{ belongs to a }T\text{-factor }\sigma \text{-algebra of }S\},$$
and we claim that it is a $`T`$-factor $`\sigma `$-algebra. Since $`(Y,,\nu )`$ is a Lebesgue space, the $`\sigma `$-algebra $``$ equipped with the metric $`d(B,C):=\nu (B\mathrm{\Delta }C)`$ is separable (as usual, we identify subsets $`B`$ and $`C`$ of $`Y`$ when $`\nu (B\mathrm{\Delta }C)=0`$). There exists a countable family $`(B_i)_{iI}`$ dense in $`_T`$, and, thanks to Lemma 3.8, there exists a $`T`$-factor $`\sigma `$-algebra $``$ containing all the $`B_i`$’s. By density, we have $`_T`$ but, since $`_T`$ contains all the $`T`$-factor $`\sigma `$-algebras, we have $`_T=`$. This proves the first assertion of Theorem 3.7. The second one is just the application of Lemma 3.9 to $`_T`$. ∎
The notion of $`T`$-factor was introduced in in order to study some aspects of the weak disjointness of dynamical systems. It was used there to prove that if $`T`$ satisfies the property of being *self-weakly disjoint to all order $`k2`$* (which is the natural generalization of self weak-disjointness to the case of $`k`$ copies of $`T`$), then $`T`$ is weakly disjoint from *any other dynamical system.* The main open question in this area is whether the self weak disjointness of $`T`$ alone is sufficient to ensure this universal weak disjointness property.
#### A characterization of disjointness?
We know from Theorem 3.5 that if $`T`$ and $`S`$ are not disjoint, then the maximum $`T`$-factor of $`S`$ is not trivial. One may ask whether the converse is true, but in fact it is not difficult to find a counterexample: Indeed, take $`(\sigma _n)_n`$ any stationary process taking its values in $`\{0,1\}`$, and let $`(\tau _n)_n`$ be an i.i.d. process with $`\tau _n`$ taking each value 0 or 1 with probability $`1/2`$, independent of $`(\sigma _n)`$. Then, setting
$$\stackrel{~}{\tau _n}:=\tau _n+\sigma _nmod2,$$
we get another i.i.d. process $`(\stackrel{~}{\tau }_n)`$ with the same distribution as $`(\tau _n)`$. The $`(\sigma _n)`$ process is clearly measurable with respect to the $`\sigma `$-algebra generated by $`(\tau _n)`$ and $`(\stackrel{~}{\tau }_n)`$, which means that the dynamical system $`S`$ generated by $`(\sigma _n)`$ is a factor of a self-joining of the Bernoulli shift generated by $`(\tau _n)`$. Therefore, denoting by $`T`$ this Bernoulli shift, we get that any dynamical system $`S`$ generated by a two-valued stationary process is a $`T`$-factor. But we could have taken for $`(\sigma _n)`$ any zero-entropy process, which is disjoint from any Bernoulli shift (Theorem 3.6).
However, we can notice that Theorem 3.5 gives a stronger conclusion when $`T`$ and $`S`$ are not disjoint. Indeed, this hypothesis being symmetric in $`S`$ and $`T`$, we get in that case that both the maximum $`T`$-factor of $`S`$ and the maximum $`S`$-factor of $`T`$ are not trivial. In the counterexample given above, $`S`$ having zero entropy, the maximum $`S`$-factor of a Bernoulli shift is trivial. So we ask the following question:
###### Question 3.10.
Can we find two disjoint systems $`S`$ and $`T`$ such that both the maximum $`T`$-factor of $`S`$ and the maximum $`S`$-factor of $`T`$ are not trivial?
#### A remark on stable properties
Analyzing the proof of Theorem 3.7, we can observe that the existence of the maximum $`T`$-factor $`\sigma `$-algebra is in fact a corollary of the following: Being a $`T`$ factor is a *stable property* (in the sense given in Section 3.3.2). And indeed, instead of the notion of $`T`$-factor we could take any stable property and we would get by the same argument the existence of a maximum factor $`\sigma `$-algebra satisfying this stable property. Some classical factor $`\sigma `$-algebras are particular cases of this type: If we start from the stable property of having zero entropy, we get the *Pinsker $`\sigma `$-algebra* of the system, and if we start from the discrete-spectrum property, we get the *Kronecker factor $`\sigma `$-algebra*.
## 4. Self-Joinings and mixing
We now turn to the relationships between the study of joinings and some (weak and strong) mixing properties. We begin by a beautiful proof by Ryzhikov of a well-known result concerning weak mixing.
### 4.1. Weak mixing and ergodicity of products
There are many different (but equivalent!) ways to define the property of weak mixing. A very concise one is the following: We say that $`T`$ is weakly mixing if the Cartesian square $`T\times T`$ is ergodic. It is a standard result in ergodic theory that this property implies the seemingly stronger one:
###### Theorem 4.1.
Suppose that $`T\times T`$ is ergodic. Then for any ergodic $`S`$, $`T\times S`$ is ergodic.
The classical proof of this implication makes use of the spectral theory of the action of the unitary operators $`U_T:ffT`$ and $`U_S:ggS`$ on the subspaces of $`L^2(\mu )`$ and $`L^2(\nu )`$ of functions with zero mean. It can be shown that if $`T`$ and $`S`$ are ergodic, the Cartesian square $`T\times S`$ is ergodic if and only if $`U_T`$ and $`U_S`$ do not have any common eigenvalue.
We propose here an alternative approach by joinings to prove Theorem 4.1, due to Ryzhikov . Let us start with a simple lemma on relatively independent joinings.
###### Lemma 4.2.
Let $`\lambda `$ be a joining of $`T`$ and $`S`$, and let $`\lambda _S\lambda `$ be the relatively independent self-joining of $`(T\times S)_\lambda `$ above its factor $`S`$. If the marginal of $`\lambda _S\lambda `$ on $`X\times X`$ is the product measure $`\mu \mu `$, then $`\lambda `$ is the product measure $`\mu \nu `$.
###### Proof.
For all $`A𝒜`$, we have by definition of $`\lambda _S\lambda `$
$$\lambda _S\lambda (A\times Y\times A)=_Y\left(𝔼_\lambda [\mathrm{𝟙}_A|y]\right)^2𝑑\nu (y).$$
Thus, if the marginal on $`X\times X`$ is the product measure, we get that for all $`A𝒜`$,
$$\left(\mu (A)\right)^2=_Y\left(𝔼_\lambda [\mathrm{𝟙}_A|y]\right)^2𝑑\nu (y)=\left(_Y𝔼_\lambda [\mathrm{𝟙}_A|y]𝑑\nu (y)\right)^2,$$
which is possible only if $`𝔼_\lambda [\mathrm{𝟙}_A|y]`$ is a constant $`\nu `$-a.e., hence $`\lambda =\mu \nu `$. ∎
###### Proof of Theorem 4.1.
Assuming that $`\mu \mu `$ is an ergodic self-joining of $`T`$, we are going to show that for all ergodic $`S`$, $`\mu \nu J_e(T,S)`$, by checking that it is an extremal point in $`J(T,S)`$, and applying Proposition 1.5.
Suppose that for some $`0<p<1`$ and some joinings $`\lambda _1`$ and $`\lambda _2`$ of $`T`$ and $`S`$, we have
$$\mu \nu =p\lambda _1+(1p)\lambda _2.$$
Note that the relatively independent self-joining of $`\mu \nu `$ above the factor $`S`$ is nothing but $`\mu \nu \mu `$, which can be decomposed as follows:
$$\mu \nu \mu =p^2\lambda _1_S\lambda _1+(1p)^2\lambda _2_S\lambda _2+p(1p)\lambda _1_S\lambda _2+p(1p)\lambda _2_S\lambda _1.$$
Since $`\mu \mu `$ is ergodic, each term of the RHS must have $`\mu \mu `$ as marginal on $`X\times X`$. But this implies that $`\lambda _1`$ and $`\lambda _2`$ satisfy the hypothesis of Lemma 4.2, hence both $`\lambda _1`$ ad $`\lambda _2`$ are the product measure $`\mu \nu `$. ∎
### 4.2. Characterization of mixing by joinings
Strong mixing is also very easily characterized in term of joinings: By definition, the dynamical system $`T`$ is *strongly mixing* (or simply *mixing*) if $`A,B𝒜`$,
$$\mu (AT^nB)\underset{n\mathrm{}}{}\mu (A)\mu (B),$$
but this means precisely that the sequence $`(\mathrm{\Delta }_{T^n})`$ converges in the topology of joinings to $`\mu \mu `$. An elegant application of this characterization was found and used by Ornstein in , where he constructed the first example of a transformation with no square root.
###### Theorem 4.3 (Ornstein’s criterion for mixing).
The dynamical system $`T`$ is mixing if and only if $`T`$ is weakly mixing and there exists some real number $`\theta >0`$ such that
(8)
$$A,B𝒜,\underset{n+\mathrm{}}{lim\; sup}\mu (AT^nB)\theta \mu (A)\mu (B).$$
###### Proof.
The fact that mixing implies weak mixing can be proved by the following argument using joinings: If $`T`$ is mixing, convergence of $`(\mathrm{\Delta }_{T^n})`$ to $`\mu \mu `$ also implies convergence of $`(\mathrm{\Delta }_{(T\times T)^n})`$ (sequence of self-joinings of $`T\times T)`$ to $`\mu \mu \mu \mu `$. And this gives that any $`T\times T`$-invariant set in $`X\times X`$ has measure zero or one, i.e. $`T`$ is weakly mixing. Condition (8) is obviously necessary for $`T`$ to be mixing.
Conversely, let us suppose that $`T`$ is weakly mixing and that (8) holds. Then any cluster point $`\lambda `$ of the sequence $`(\mathrm{\Delta }_{T^n})`$ in $`J(T)`$ satisfies $`\lambda \theta \mu \mu `$, hence is absolutely continuous with respect to $`\mu \mu `$. Since $`T`$ is weakly mixing, $`\mu \mu `$ is ergodic, therefore $`\lambda =\mu \mu `$. This proves that
$$\underset{n\mathrm{}}{lim}\mathrm{\Delta }_{T^n}=\mu \mu ,$$
i.e. $`T`$ is mixing. ∎
### 4.3. 2-fold and 3-fold mixing
We finish this presentation by a long-standing open question in ergodic theory, and some of its relationships with joining theory. Recall that the property of (strong) mixing is also called *2-fold mixing*, because it involves only two subsets of the space. A strengthening of this definition involving 3 (or $`k3`$) sets gives rise to the property of *3-fold mixing* (respectively *$`k`$-fold mixing*):
###### Definition 4.4.
The dynamical system $`T`$ is said to be *3-fold mixing* if $`A,B,C𝒜`$,
$$\underset{n,m\mathrm{}}{lim}\mu (AT^nBT^{(n+m)}C)=\mu (A)\mu (B)\mu (C).$$
Again, this definition can easily be translated into the language of joinings: $`T`$ is 3-fold mixing when the sequence $`(\mathrm{\Delta }_{T^n,T^{n+m}})`$ converges in $`J^3(T)`$ to $`\mu \mu \mu `$ as $`n`$ and $`m`$ go to infinity, where $`\mathrm{\Delta }_{T^n,T^{n+m}}`$ is the obvious generalization of $`\mathrm{\Delta }_{T^n}`$ to the case of self-joinings of order 3.
Whether 2-fold mixing implies 3-fold mixing was asked long ago by Rohlin , and the question is still open today. However, some important results have been established around this problem. Some of them could make us think that this implication is false, by exhibiting counterexamples to the same question but in a generalized context. In this category, we must in particular cite Ledrappier’s example in of a $`^2`$-action which is 2-fold but not 3-fold mixing. More recently, it was noticed in that if we consider the property of being mixing relatively to a factor $`\sigma `$-algebra (roughly speaking, this means that we replace the measure by the conditional expectation with respect to this $`\sigma `$-algebra), then we can find examples of 2-fold but not 3-fold mixing. On the other hand, there are also many results which go in the opposite direction. We know that 2-fold mixing implies 3-fold mixing in many cases: For example when $`T`$ has singular spectrum (Host, ), when $`T`$ is rank one (Kalikow, ) or even for finite rank systems (Ryzhikov, ), when $`T`$ is generated by a Gaussian system (Leonov, ; see also Totoki ). In fact, in many of these works it is shown that a stronger result holds for the respective class of systems under consideration, and this result concerns joining theory: For finite rank systems, as well as for singular spectrum systems, it is proved that any ergodic self-joining $`\lambda `$ of order 3 which has pairwise independent marginals is necessarily the product measure. As far as Gaussian systems are concerned, this property has been established for an important subclass of the zero-entropy Gaussian processes, called GAG . It is not known whether it can be extended to all zero-entropy Gaussian processes. More generally, we may ask the following important question:
###### Question 4.5.
Does there exist a zero-entropy dynamical system $`T`$ and an ergodic joining $`\lambda J_3(T)`$ for which the coordinates are pairwise independent but which is different from $`\mu \mu \mu `$?
Why a negative answer to this question would solve Rohlin’s multiple mixing problem is explained by the following remarks: First, if we had a 2-fold mixing but not 3-fold mixing dynamical system $`T`$, the sequence $`(\mathrm{\Delta }_{T^n,T^{n+m}})`$ would have a cluster point $`\lambda \mu \mu \mu `$. Since $`T`$ is 2-fold mixing, the coordinates would be pairwise independent for $`\lambda `$. And we could find an ergodic joining satisfying those two properties by taking an ergodic component of $`\lambda `$. Second, Thouvenot proved that if we could find a counterexample to Rohlin’s question, then we could find one in the zero-entropy category. This restriction to zero entropy is crucial since we can easily construct examples answering Question 4.5 if we drop it: Start from the $`(1/2,1/2)`$ Bernoulli shift $`T`$ generated by the i.i.d. process $`(\xi _p)_p`$, taking its values in $`\{0,1\}`$. Take an independent copy $`(\xi _p^{})_p`$ of this process and set
$$\xi _p^{\prime \prime }:=\xi _p+\xi _p^{}mod2.$$
Then the three processes $`\xi `$, $`\xi ^{}`$ and $`\xi ^{\prime \prime }`$ have the same distribution, are pairwise independent, but the self-joining of $`T`$ of order 3 that we get in this way is not the product measure.
Observe also that a negative answer to Question 4.5 would imply another important result in joining theory: It would show that 2-fold MSJ implies $`k`$-fold MSJ for all $`k2`$. Indeed, any 2-fold MSJ but not 3-fold MSJ would have an ergodic joining $`\lambda `$ of order 3 for which the coordinates would be pairwise independent, but not independent.
The purpose of the following is to show that, if ever we could find some system satisfying the requirements of Question 4.5, it must be of a different nature than the construction of Ledrappier for an action of $`^2`$, or than the relative example given in . Indeed, each of these two constructions leads to some “joining” $`\lambda `$, for which we get 3 pairwise independent coordinates $`x_1`$, $`x_2`$ and $`x_3`$, and for which the global independence is denied by the following property: We have some random variable $`\xi `$ defined on $`X`$, taking its value in $`\{0,1\}`$ (each value being taken with probability $`1/2`$), and satisfying
(9)
$$\xi (x_3)=\xi (x_1)+\xi (x_2)mod2(\lambda \text{-a.e.}).$$
There are two ways of interpreting equation (9): First, the random variable $`\xi `$ takes its values in a finite alphabet $`A`$, and there exists a function $`f:A\times AA`$ such that
$$\xi (x_3)=f(\xi (x_1),\xi (x_2))(\lambda \text{-a.e.}).$$
Second, the random variable $`\xi `$ takes its values in a compact abelian group $`G`$, and we have
$$\xi (x_3)=\xi (x_1)+\xi (x_2)(\lambda \text{-a.e.}).$$
We want to show now that there cannot exist an example answering Question 4.5 of either of those types, by proving the following two propositions.
###### Proposition 4.6.
Suppose that there exist some joining $`\lambda J_3(T)`$, with pairwise independent coordinates, and some random variable $`\xi :XA`$ (finite alphabet) such that
$$\xi (x_3)=f(\xi (x_1),\xi (x_2))(\lambda \text{-a.e.}).$$
Then
* either the factor of $`T`$ generated by $`\xi `$ is periodic,
* or this factor has an entropy at least $`\mathrm{log}2`$.
###### Proposition 4.7.
The dynamical system $`T`$ is assumed to be ergodic. Suppose that for some joining $`\lambda J_3(T)`$ with pairwise independent marginals, there exists some measurable map $`\xi :XG`$ (compact Abelian group) which is not $`\mu `$-a.e. constant, and such that
$$\xi (x_3)=\xi (x_1)+\xi (x_2)(\lambda \text{-a.e.})$$
then $`T`$ has positive entropy.
###### Lemma 4.8.
Let $`\xi _1`$, $`\xi _2`$ and $`\xi _3`$ be 3 random variables taking their values in the finite alphabet $`A`$, identically distributed, pairwise independent, and such that there exists a function $`f:A\times AA`$ satisfying
$$\xi _3=f(\xi _1,\xi _2)(\text{a.e.})$$
Then their common probability law is the uniform law on the set of values which are taken with positive probability.
###### Proof.
Taking a subset of $`A`$ instead of $`A`$ if necessary, we can assume that for all $`iA`$,
$$p_i:=P(\xi _1=i)>0.$$
For all $`i`$ and $`k`$ in $`A`$, there exists a unique $`j=j_k(i)`$ such that $`f(i,j_k(i))=k`$. Moreover, we have
$$p_k=\underset{i}{}p_iP(\xi _3=k|\xi _1=i)=\underset{i}{}p_ip_{j_k(i)}.$$
In other words, $`\pi =M\pi `$ where $`\pi :=(p_i)_{iA}`$ and $`M`$ is the bistochastic matrix whose entries are $`m_{k,i}:=p_{j_k(i)}`$. Taking the limit as $`s\mathrm{}`$ in the equality $`\pi =M^s\pi `$, we get that $`\pi `$ has to be the uniform law on $`A`$. ∎
###### Proof of Proposition 4.6.
From the preceding lemma, we know that $`\xi `$ is uniformly distributed on the subset of values which are taken with positive probability. But for any $`m1`$, Lemma 4.8 also applies to
$$\xi _0^{m1}:x(\xi (x),\xi (Tx),\mathrm{},\xi (T^{m1}x)).$$
Let $`i_0,\mathrm{},i_{m2}`$ in $`A`$ be such that
$$\mu \left(\xi (x)=i_0,\xi (Tx)=i_1,\mathrm{},\xi (T^{m2}x)=i_{m2}\right)>0.$$
Then, the conditional probability law of $`\xi (T^{m1}x)`$ knowing the above event is uniform on the set of values to which it assigns a positive mass. Moreover, we easily check that the cardinal $`a_m`$ of the support of this conditional probability law does not depend on the choice of $`i_0,\mathrm{},i_{m2}`$, and satisfies $`a_ma_{m1}`$. Thus,
* either the sequence $`(a_m)`$ reaches 1, and in this case the factor generated by $`\xi `$ is periodic;
* or $`a_m`$ is always greater than or equal to 2, and then the entropy of the factor generated by $`\xi `$ is at least $`\mathrm{log}2`$.
###### Lemma 4.9.
Let $`\xi _1`$, $`\xi _2`$ and $`\xi _3`$ be 3 random variables taking their values in the compact Abelian group $`G`$, identically distributed, pairwise independent, and satisfying
$$\xi _3=\xi _1+\xi _2(\text{a.e.})$$
Then their common probability law is the Haar measure on some compact subgroup of $`G`$.
###### Proof.
Denote by $`\nu `$ the probability law of $`\xi _1`$ on $`G`$. For $`\nu `$-almost every $`gG`$, $`\nu `$ is invariant under the addition of $`g`$, since knowing $`\xi _2=g`$, both $`\xi _1`$ and $`\xi _3=\xi _1+g`$ are distributed according to $`\nu `$. Let us define
$$H:=\{gG:\nu \text{ is invariant under the addition of }g\}.$$
This is a closed subgroup of $`G`$ containing the support of $`\nu `$. Thus $`\nu `$ is a probability law on $`H`$, invariant under every translation on $`H`$: $`\nu `$ is the Haar measure on $`H`$. ∎
###### Proof of Proposition 4.7.
Let us consider $`\psi :XG^{}`$ defined by
$$\psi (x):=(\xi (T^nx))_n.$$
From Lemma 4.9, the probability law of $`\psi (x)`$ is the Haar measure on some compact subgroup $`H`$ of $`G^{}`$. Moreover, $`H`$ is invariant by the shift, denoted by $`\sigma `$, on $`G^{}`$. The restriction of $`\sigma `$ to $`H`$ is therefore an ergodic endomorphism of $`H`$, which is a factor of $`T`$ (This is the factor generated by $`\xi `$). A theorem by Juzvinskii tells us that any such ergodic group endomorphism has to be a $`K`$-system; hence $`T`$ has positive entropy. ∎
Observe that the two situations exposed in Propositions 4.6 and 4.7 can be seen as particular cases of a third one, from which I do not know whether some similar conclusion can in general be derived.
###### Question 4.10.
What can be said about a dynamical system $`T`$ satisfying the following: There exists an ergodic self-joining $`\lambda J_3(T)`$, with pairwise independent marginals, and a map $`f:X\times XX`$ such that
$$x_3=f(x_1,x_2)(\lambda \text{-a.e.})\mathrm{?}$$
## Acknowledgements
I want to express my special thanks to Jean-Paul Thouvenot and Valery Ryzhikov, from whom I learnt most of what I know on joinings of dynamical systems; to Anthony Quas for his help concerning Propositions 4.6 and 4.7; to El Houcein El Abdalaoui, Élise Janvresse and Yvan Velenik for their support and advice in the preparation of this review.
I am also grateful to the referees for having corrected several mistakes in the first version of this work, and suggested substantial improvements in the text. |
warning/0507/hep-th0507107.html | ar5iv | text | # Generalized Particle Statistics in Two-Dimensions: Examples from the Theory of Free Massive Dirac Field
## 1 Introduction
The intrinsic definition of particle statistics in the approach of Algebraic Quantum Field Theory (AQFT) in a four-dimensional space-time is provided by assigning to superselection sectors equivalence classes of permutation group representations, which describe the statistics of multiparticle states. In a (3+1)-dimensional space-time, fields and particles obey Bose-Fermi alternative, exhibiting the more general bosonic or fermionic parastatistics, while in lower dimensional Minkowski space statistics are described, in general, by braid group representations. The first models leading to particles described by a one-dimensional representation of the braid group (anyons) are in , while higher dimensional representations describe plektons. In a (2+1)-dimensional space-time, strictly local quantum fields are always subject to the normal commutation rules, but particles carrying “topological charges”, created from the vacuum by the action of fields localized in cones, may exhibit intermediate statistics.
The statistics of a sector describes the interchange of identical charges. In two dimensions, DHR theory allows for two distinct statistics operators (one the inverse of the other), since the causal complement of a bounded region has two connected components. The statistics operator is a topological invariant if the pairs of spatially separated auxiliary regions can be continuously deformed from one to the other, maintaining a relative space-like distance. Therefore, the braid group enters in the description of DHR superselection charges localized in two-dimensional double cones, for intervals of the real line or in (2+1)-dimensional theories for charges localized in space-like cones.
Superselection sectors in four-dimensional theories are classified by equivalence classes of irreducible representations of the compact group of internal symmetries . However, if the superselection category in low dimensional theories is not symmetric but only braided, such a group may not exist. Indeed, some models of (1+1)-dimensional conformal fields exhibit a superselection structure which seems not to fit any representation group theory.
Why “generalized” particles statistics? Well, this is necessary since the algebraic approach to local field theories which do not fulfill Haag duality and which does not allow non inner automorphisms of the underlying field algebra does not yield a well defined notion of statistics. Thus, we need to extend it to physical theories which do not fit the prescriptions of the algebraic framework totally, as in the case of smeared-out kink operators in the context of the two-dimensional free massive Dirac field in the formalism of relativistic second quantization developed in . There, not only is Haag duality violated (Sect. 3), but a family of unitary operators implementing DHR automorphisms is not in the field algebra, forcing us to explore alternative tools. In a more general setting, field theories in (1+1)-dimensions satisfying twisted Haag duality and the split property for wedges and having an unbroken (i.e. unitarily implemented) group of inner symmetries $`G`$ give rise to a not Haag dual observable algebra $`𝒜=^G`$ . Split property for wedges has been proven recently by Buchholz and Lechner for the Bose and Fermi cases . Together with the argument in , this proves that that the observable algebra is not Haag dual when $``$ is any finite product of free massive Bose and Fermi fields and $`G`$ is non-trivial.
We remark that “free” anyons are studied in a two-dimensional space-time since no (2+1)-dimensional model of free anyons can exist<sup>1</sup><sup>1</sup>1This notion of “free” anyons refers to the on-mass-shell nature of the Fourier transform. In d=1+1 the anyon operators can live together in the same Hilbert space as the free Fermions, but they are not really free in the mass-shell sense. . In the setting of CAR algebras on the fermionic Fock space there exists a natural notion of second quantization more appropriate for a theory of relativistic particles. The theory of Fermionic gauge groups displays a wide class of unitarily implementable automorphisms on the antisymmetric Fock space.
Our choice is the natural one , i.e. that for which the winding numbers (the charges) are easily computable through index formulae , while the zero charge implementers are the well known smeared-out kink operators since they can modify the statistics of a sector . Implementable gauge groups in the one-particle Dirac theory lead to a model which exhibits strange statistics. A class of Bogoliubov automorphisms unitarily implementable in the Fock representation induce a family of localized and transportable automorphisms of the observable algebra , implemented by non local operators which are not even contained in the field algebra $`.`$
Since our investigations are strongly influenced by the violation of Haag duality and the non locality of implementers, which gives rise to non inner automorphisms, we begin with a discussion of the arguments leading to the known results, in order to emphasize the need to clarify the notion of statistics even for theories not fulfilling all the axioms of AQFT, and to better understand the developments presented in this paper. Statistics of sectors, approached first with “classical” DHR theory , depends not only on the charge (i.e. on the sector), but also on a continuous parameter which indexes a collection of unitarily implementable automorphisms which carry no charge, but modify the statistics of the composed sector .
Unfortunately, since the net of local observables does not fulfill Haag duality, some results largely exploited in AQFT are no longer true in the setting with which we are concerned here, and we are able to produce counterexamples. However, the statistics operator still possesses all of the formal properties as it does in DHR theory, since unitary intertwiners between automorphisms of the same translation equivalent class are always local elements of $`𝒜,`$ even if Haag duality is violated.
After computing the statistics operator formally, the question remains as to whether the braiding obtained in this way has a genuine meaning in terms of statistics. Actually, we cannot proceed step by step along DHR theory alone, as it deals with local objects and often exploits Haag duality as a fundamental technical assumption. We are now analyzing a theory that allows for intertwiners not lying in the algebra where endomorphisms act and where endomorphisms are not locally inner but are inner only in a asymptotical sense.
We appeal to a more recent notion of braiding , where the condition of asymptotic abelianness of intertwiners allow us to define $`\text{bi-asymptopias},`$ giving rise to a braiding for a suitable full subcategory of $`\text{End}(𝒜).`$ In higher dimensions, Roberts has shown that a DHR sector of a non-Haag dual net $`𝒜`$ extends to a DHR sector of the dual net $`𝒜^d,`$ and the latter can be studied with the usual methods . In (1+1)-dimensional massive theories this fails since $`𝒜^d,`$ satisfying Haag duality and split property for wedges, has no localized sectors as Mueger has shown in . The net $`𝒜`$ may have non-trivial localized sectors, but they necessarily become solitons when they are extended to $`𝒜^d.`$
In (1+1)-dimensional free massive Dirac field theory we exclude those braidings that do not give rise to true statistics, since they have their DHR counterpart in pseudo statistics operators constructed without remaining in the same connected component. Bi-asymptopias relative to different components are not mutually cofinal, nor path connected, due to the geometry of a two-dimensional space-time. A direct computation for each connected component yields different braidings, i.e. the category is not symmetric. Physically speaking, particles described by this theory are neither bosons nor fermions for almost all values of the solitonic parameter $`\lambda .`$
The present article is organized as follows. In Sect. 2 our assumptions are stated: relativistic second quantization, implementable gauge groups in the (1+1)-dimensional free massive Dirac field theory, index formulae for smeared-out kink operators in the formalism developed in , , .
In Sect. 3 we analyze a field theory model arising from the fermionic gauge group theory when applied to the Dirac field in two-dimensional Minkowski space. Some known results from are derived. Investigation of the statistics in the framework of DHR analysis leads to anomalous behaviour of the charge composition, and statistics is not an intrinsic property of the sector <sup>2</sup><sup>2</sup>2In the conformally invariant zero mass (or short distance) limit the situation changes radically and the standard DHR theory becomes again fully applicable. This phenomenon is inexorably linked to the emergence of new sectors in this limit (the disorder becomes charge-carrying) and has been observed in .
In Sect. 4 we prove that unitary implementers are not elements of the global field algebra $``$ for almost all values of the continuous parameter $`\lambda ,`$ thus giving a complete classification of their localization properties.
In Sect. 5 we prove asymptotic abelianness in order to exhibit a pair of disjoint (i.e not path connected) bi-asymptopias which give rise to an authentic braiding for the $`C^{}`$-subcategory of $`\text{End}(𝒜)`$ generated by our family $`\mathrm{\Delta }`$ of automorphisms. We compute the braiding explicitly and give a natural condition to be imposed in (1+1)-dimension in order to avoid those braidings with no true counterpart in the DHR setting.
## 2 Preliminaries
Since we work in the algebraic setting, we briefly list the axioms appropriate for theories where observables are defined from fields through a principle of gauge invariance.
1. The field algebra $``$ is the inductive limit of the net of von Neumann algebras $`𝒪(𝒪)`$ and its action on the Hilbert space $``$ is irreducible.
2. There exists a strongly continuous unitary representation $`L𝒰(L)`$ of the Poincaré group $`𝒫`$ on $``$ inducing automorphisms $`\alpha _L`$ of the field algebra, and the action on local algebras is geometric, i.e. $`\alpha _L((𝒪))=(L𝒪).`$ Moreover, there exists a unit vector $`\mathrm{\Omega },`$ the vacuum vector, unique up to a phase, which is left invariant by $`𝒰(L),L𝒫.`$ The vector state induced by $`\mathrm{\Omega }`$ is the vacuum state $`\omega _0`$ of $`,`$
$$\omega _0(F)=(\mathrm{\Omega },F\mathrm{\Omega }).$$
3. (Reeh-Schlieder property for double cones) The vacuum vector $`\mathrm{\Omega }`$ is cyclic and separating for every algebra $`(𝒪).`$
4. There exists a faithful representation $`g\beta _g`$ of a compact group $`G,`$ the gauge group, by automorphisms of $`.`$ $`\beta _g`$ commutes with $`\alpha _L`$ and $`\beta _g((𝒪))=(𝒪).`$ Moreover, for $`F(𝒪),`$ the correspondence $`g\beta _g(F)`$ is weakly continuous.
5. (Normal commutation relations) There exists a $`kG,`$ with $`k^2=e,`$ such that, setting
$$_\pm (𝒪)=\{F(𝒪):\beta _k(F)=\pm F\},$$
(1)
we have that $`_+(𝒪_1)`$ commutes with $`(𝒪_2)`$ and $`_{}(𝒪_1)`$ anticommutes with $`_{}(𝒪_2)`$ when $`𝒪_1𝒪_2^{}.`$ If the unitary $`\mathrm{\Gamma }(g)`$ implements the automorphism $`\beta _g,`$ we can reformulate (1) by requiring twisted locality:
$$(𝒪)^\tau (𝒪^{})^{}$$
where $`(𝒪)^\tau :=Z(𝒪)Z^{},`$ $`Z=\frac{\mathrm{𝟙}+i\mathrm{\Gamma }(k)}{1+i},`$ defines the twisted algebra.
### 2.1 Relativistic second quantization
We now fix notation and give an overview of the fundamentals of relativistic second quantization. Let $``$ be a separable complex Hilbert space with inner product $`(,).`$ The fermionic Fock space $`_a()`$ is the completion of the vector space $`𝒟_{at}:=_{n=0}^{\mathrm{}}^n`$ of antisymmetric algebraic tensors with respect to the ”natural” scalar product $`<_{n0}\xi _n|_{n0}\eta _n>=_{n0}(\xi _n,\eta _n),`$ with the standard unity vector $`\mathrm{\Omega }:=(1,0,0,\mathrm{})`$ defined as the vacuum vector. For every $`f,`$ the creation operator $`c^{}(f)`$ and its adjoint $`c(f)`$, the annihilation operator, are defined on the whole fermionic Fock space, with $`c^{}(f)=f,`$ and they satisfy the canonical anticommutation relations:
$$\{c(f),c(g)\}=\{c^{}(f),c^{}(g)\}=0$$
$$\{c(f),c^{}(g)\}=(f,g)\mathrm{𝟙}$$
for arbitrary $`f,g.`$ If $`U()`$, we define the operator $`\mathrm{\Gamma }(U)`$ on $`\overline{^n}`$ by
$$\mathrm{\Gamma }(U)c^{}(f_1)\mathrm{}c^{}(f_n)\mathrm{\Omega }=c^{}(Uf_1)\mathrm{}c^{}(Uf_n)\mathrm{\Omega },$$
(2)
which has the property $`\mathrm{\Gamma }(U)\mathrm{\Gamma }(V)=\mathrm{\Gamma }(UV).`$ In particular, if $`U𝒰():=\{A:,A\text{ unitary}\}`$, the correspondence $`c^{}(f)c^{}(Uf)`$ is an automorphism of the CAR algebra, unitarily implemented by $`\mathrm{\Gamma }(U)`$ (in the Fock representation):
$$\mathrm{\Gamma }(U)c^{}(f)\mathrm{\Gamma }(U)^{}=c^{}(Uf)$$
(3)
as follows from (2). Moreover, an arbitrary $`A()`$ induces on $`_a()`$ a sum operator $`\text{d}\mathrm{\Gamma }(A)`$ such that $`\mathrm{exp}(it\text{d}\mathrm{\Gamma }(A))=\mathrm{\Gamma }(e^{itA}),`$ which preserves the adjoint and the commutator .
In concrete cases, as that we discuss in this paper, $`=_+_{}`$, where $`_\pm `$ are copies of the same function space $`L^2H`$. Let $`P_\delta `$ ($`\delta =+,`$) be the projectors onto $`_\pm `$. If $`A`$ is an operator on $`,`$ we set $`A_{\delta \delta ^{}}:=P_\delta AP_\delta ^{},`$ $`\delta ,\delta ^{}=+,.`$ In this notation, $`A`$ is a block matrix whose entries are endomorphisms of $`L^2`$. Such a decomposition of $``$ is related to the well known fact that the free Dirac hamiltonian $`H_m`$ has spectrum $`(\mathrm{},m][m,+\mathrm{})`$, where $`m0`$ denotes the rest mass of the particle. Here, $`P_\pm `$ are the spectral projectors of $`H_m`$ onto $`[m,+\mathrm{})`$ and $`(\mathrm{},m]`$ respectively. Instead of non-relativistic second quantization $`A\text{d}\mathrm{\Gamma }(A),`$ we work with another irreducible representation of the CAR algebra, defined by
$$\stackrel{~}{c}(f):=c(P_+f)+c^{}(\overline{P_{}f}).$$
(4)
The one-particle Hilbert space is then defined by $`_1:=P_+H\overline{P_{}H}`$ while the physical Hilbert space is $`_a(_1).`$ If $`U𝒰()`$ satisfies $`[U,P_\delta ]=0`$, then the automorphism $`\stackrel{~}{c}(f)\stackrel{~}{c}(Uf)`$ is unitarily implemented by
$$\stackrel{~}{\mathrm{\Gamma }}(U):=\mathrm{\Gamma }(U_{++})\mathrm{\Gamma }(\overline{U}_{}),$$
(5)
with the compact notation $`\mathrm{\Gamma }(U_{++})\mathrm{\Gamma }(U_{++}\mathrm{𝟙}),\mathrm{\Gamma }(\overline{U}_{})\mathrm{\Gamma }(\mathrm{𝟙}\overline{U}_{})`$. Here the bar over an operator stands for the action by a fixed conjugation $`J`$ on $`,`$ i.e. $`\overline{U}=JUJ.`$ For arbitrary $`U𝒰(),`$ the Shale-Stinespring theorem states that there exists a unitary $`\stackrel{~}{\mathrm{\Gamma }}(U)`$ on the antisymmetric Fock space such that
$$\stackrel{~}{c}(Uf)=\stackrel{~}{\mathrm{\Gamma }}(U)\stackrel{~}{c}(f)\stackrel{~}{\mathrm{\Gamma }}(U)^{}f$$
if and only if the off-diagonal parts of $`U,`$ namely $`U_{\delta \delta }`$, are Hilbert-Schmidt ($`HS`$) operators. Unitaries on $``$ inducing automorphisms of the CAR algebra which are unitarily implementable on $`_a()`$ form a group, denoted $`𝒢_2`$. Let $`𝔤_2:=\{A():A_{\delta ,\delta }HS\}`$ be the complex Lie algebra of $`𝒢_2.`$ By a suitable choice for the phases of the unitary operator implementing the automorphism of the CAR algebra, it is always possible to define a one-parameter strongly continuous group
$$\stackrel{~}{\mathrm{\Gamma }}(e^{itA})=e^{it\text{d}\stackrel{~}{\mathrm{\Gamma }}(A)},A=A^{}𝔤_2$$
where $`\text{d}\stackrel{~}{\mathrm{\Gamma }}(A)`$ is the self-adjoint generator. The arbitrary additive constant in the definition of $`\text{d}\stackrel{~}{\mathrm{\Gamma }}(A)`$ is fixed by requiring that $`(\mathrm{\Omega },\text{d}\stackrel{~}{\mathrm{\Gamma }}(A)\mathrm{\Omega })=0.`$ In Sect. 3, where we deal with the Dirac field, we shall employ the more common notation $`\pi (\varphi (f))`$ rather than $`\stackrel{~}{c}(f).`$
###### Definition 1
The *charge operator* Q is the generator of the one-parameter group induced by the identity:
$$Q:=\text{d}\stackrel{~}{\mathrm{\Gamma }}(\mathrm{𝟙})=\text{d}\mathrm{\Gamma }(P_+)\text{d}\mathrm{\Gamma }(P_{}).$$
Under the action of charge operator, fermionic Fock space splits into a direct sum of charge sectors $`_a()=_n_n,`$ and
$$\stackrel{~}{\mathrm{\Gamma }}(U)_n=_{n+q(U)},U𝒢_2,$$
where $`q(U)`$ is the Fredholm index of $`U_{}.`$
The additive property $`q(U_1U_2)=q(U_1)+q(U_2)`$ also holds. In other terms, if $`q(U)=1,`$ the vacuum sector $`_0`$ can be connected to the $`n`$ charge sector by applying $`\stackrel{~}{\mathrm{\Gamma }}(U)^n,`$ while $`q(e^{itA})=0,A=A^{}𝔤_2,`$ since $`q(e^{itA})=q(\mathrm{𝟙})=0`$ by virtue of the continuity in $`t`$. Hence the charged sectors are left invariant by $`\stackrel{~}{\mathrm{\Gamma }}(e^{itA})`$. We end this overview with an identity of great relevance for computations:
$$\mathrm{\Gamma }(\mathrm{𝟙})\stackrel{~}{\mathrm{\Gamma }}(U)=(1)^{q(U)}\stackrel{~}{\mathrm{\Gamma }}(U)\mathrm{\Gamma }(\mathrm{𝟙}).$$
In view of the subsequent applications, we cite three useful propositions which establish the commutation rules between unitary implementers and their self-adjoint generators .
###### Proposition 1
For every $`A,B𝔤_2`$, on the domain $`𝒟`$ of finite particle vectors there holds
$$[\text{d}\stackrel{~}{\mathrm{\Gamma }}(A),\text{d}\stackrel{~}{\mathrm{\Gamma }}(B)]=\text{d}\stackrel{~}{\mathrm{\Gamma }}([A,B])+C(A,B)\mathrm{𝟙},$$
(6)
where $`C(A,B):=\text{Tr}(A_+B_+B_+A_+)`$ is the *Schwinger term*.
Here $`𝒟:=\{F_a():F=P_lF,\text{ for some }l\}`$ and $`P_l`$ denotes the spectral projector of the particle number operator on $`[0,l].`$
###### Proposition 2
Let $`A,B𝔤_2,A=A^{},B=B^{}\text{ and }[A,B]=0`$. We have:
$$\stackrel{~}{\mathrm{\Gamma }}(e^{iA})\stackrel{~}{\mathrm{\Gamma }}(e^{iB})=e^{C(A,B)/2}\stackrel{~}{\mathrm{\Gamma }}(e^{i(A+B)}).$$
(7)
The third proposition establishes the commutation rule between second quantization operators in the case that one of them is a charge shift (i.e. it carries a non zero charge).
###### Proposition 3
Let $`U𝒢_2,A=A^{}𝔤_2,[U,A]=0.`$ Then
$$\stackrel{~}{\mathrm{\Gamma }}(e^{iA})\stackrel{~}{\mathrm{\Gamma }}(U)=e^{i(\stackrel{~}{\mathrm{\Gamma }}(U)\mathrm{\Omega },\text{d}\stackrel{~}{\mathrm{\Gamma }}(A)\stackrel{~}{\mathrm{\Gamma }}(U)\mathrm{\Omega })}\stackrel{~}{\mathrm{\Gamma }}(U)\stackrel{~}{\mathrm{\Gamma }}(e^{iA}).$$
(8)
### 2.2 Implementable gauge groups in the one-particle Dirac theory
In the theory of (1+1)-dimensional free massive Dirac field, gauge transformations are operators of multiplication by unitary matrices on $`\stackrel{ˇ}{}L^2(,dx)^2,`$ the image of $`L^2(,dp)^2`$ by means of the Fourier transformation $`,`$ employed to diagonalize the differential operator representing the free Dirac Hamiltonian. Since we are interested in lifting these unitaries to the Fock space, we must consider only gauge transformations which define unitaries in $`𝒢_2.`$ We denote by $`H_1()`$ the Sobolev space, which consists of all absolutely continuous functions of $`L^2()`$ with derivatives in $`L^2().`$ Once we have introduced the group (under pointwise multiplication)
$$L_e\text{U}(1):=\{u\text{Map}(,\text{U}(1))|u()1H_1()\},$$
we define two faithful unitary representations $`\stackrel{ˇ}{\pi }_\pm `$ of $`L_e\text{U}(1)`$ on $`\stackrel{ˇ}{},`$ given by $`\stackrel{ˇ}{\pi }_+(u)=u(x)\mathrm{𝟙}`$ and $`\stackrel{ˇ}{\pi }_{}(u)=\mathrm{𝟙}u(x),`$ i.e. they act as multiplication by a function on one component space only. Then $`\pi _s(u)𝒢_2,`$ and two projective unitary representations $`\stackrel{~}{\mathrm{\Gamma }}(\pi _\pm )`$ on $`_a()`$ are automatically defined.
Another global gauge transformation has the form $`e^{i\phi _+}e^{i\phi _{}},\phi _\pm (0,2\pi ).`$
In the case of interest to us, i.e. rest mass of the particle $`m>0,`$ we put $`\phi _+=\phi _{},`$ otherwise the $`HS`$ condition would be violated.
We end this paragraph with a formula for $`\text{Ind}U_{}.`$ We are interested in operators of the form
$$(\stackrel{ˇ}{U}f)(x)=u(x)f(x),u(x)\text{U}(2),f\stackrel{ˇ}{},$$
where the $`2\times 2`$ matrix $`u(x)`$ is assumed to be diagonal:
$$u(x)=\left(\begin{array}{cc}u_+(x)& 0\\ 0& u_{}(x)\end{array}\right),u_\pm (x).$$
(9)
We set $`A:=^1\stackrel{ˇ}{A}`$ for $`\stackrel{ˇ}{A}`$ an operator on $`\stackrel{ˇ}{}.`$ We consider continuous multipliers of the form (9) such that, for each $`s,`$ there exists $`u_{\mathrm{}}C(\{\pm 1\},\text{U}(1))`$ satisfying
$$u_s(x)u_{\mathrm{}}\left(\frac{x}{|x|}\right)=o(1),|x|\mathrm{},s=+,.$$
These multipliers form a group, denoted by $`G_h,`$ and their Fredholm indices are easily computable in view of the following .
###### Theorem 1
Let $`UG_h.`$ Then
$$\text{ Ind }U_{}=\text{w}(u_+u_{}^1),$$
(10)
where w is the winding number which, by convention, is positive on the map $`x\frac{xi}{x+i}.`$
In order to complete the discussion of the operators we shall employ in the next section, we remark that all our unitaries induce automorphisms of the CAR algebra which are unitarily implementable on the Fock space. Indeed, if $`x_1\alpha (x_1)`$ is an odd, monotonously increasing, $`C^{\mathrm{}}`$ real valued function, which equals $`1`$ at the right of the interval $`(1,1),`$ we introduce the smeared-out kink operators
$$\stackrel{ˇ}{U}_{\lambda ,ϵ}:=e^{i\pi \lambda \alpha (/ϵ)},\lambda ,ϵ>0.$$
(11)
The off-diagonal parts of $`U_{\lambda ,ϵ}`$ are $`HS`$ for every $`\lambda ,`$ so it induces Bogoliubov transformations unitarily implementable for every $`\lambda `$ .
## 3 Strange statistics in two-dimensional free massive Dirac field theory
The theory of fermionic gauge group reveals itself as a natural setting for the construction of a model which exhibits anomalous statistics . The initial Hilbert space is $`L^2(,^2)`$. Denoting by $`𝒦`$ the set of double cones in the (1+1)–dimensional Minkowski space, let $`B_𝒪`$ be the base at time $`t`$ of $`𝒪𝒦.`$ The algebra of fields localizable in $`𝒪`$ is defined in the usual way,
$$(𝒪)=\{\pi (\varphi (e^{itH_m}f)):\text{supp}(f)B_𝒪\}^{\prime \prime },$$
and the global field algebra $``$ is the $`C^{}`$-inductive limit of the net $`\{(𝒪)\}_{𝒪𝒦}.`$ The gauge invariant parts (i.e. the subalgebras left invariant by $`\text{Ad}\stackrel{~}{\mathrm{\Gamma }}(e^{i\gamma }),`$ $`\gamma `$) form the net of observables. Let us mention that the free massive Dirac field theory in (1+1)-dimensions fulfills twisted duality $`(𝒪)^\tau =(𝒪^{})^{}`$ (the proof of this is independent of the space-time dimension; see for details), but the net of observables does not fulfil Haag duality for double cones. Indeed, if $`𝒪,𝒪_1,𝒪_2`$ are double cones, with bases at $`t=0,`$ such that $`𝒪_1,𝒪_2`$ lie in different connected components of $`𝒪^{},`$ and if $`f_i`$ are test functions with $`\text{supp}(f_i)𝒪_i(i=1,2),`$ then the observable $`\pi (\varphi (f_1))^{}\pi (\varphi (f_2))`$ is contained in $`𝒜(𝒪)^{}`$ but not in $`𝒜(𝒪^{})^{\prime \prime }`$ .
The automorphism of $``$ defined by $`\rho __Z(\pi (\varphi (f))):=\pi (\varphi (Zf))`$ is unitarily implemented in the Fock representation when $`Z`$ is one of the following unitaries of $`L^2(,^2)`$:
$$(U(n)f)(x)=(e^{i\pi n\epsilon (x)}f_1(x),f_2(x)),n$$
$$(V(\lambda )f)(x)=e^{i\pi \lambda \vartheta (x)}f(x),\lambda ,$$
which correspond in (9) to the choice $`u_+(x)=\mathrm{exp}(i\pi n\epsilon (x)),`$ $`u_{}0`$ and $`u_\pm (x)=\mathrm{exp}(i\pi \lambda \vartheta (x)),`$ respectively. Here, the functions $`\epsilon `$ and $`\vartheta `$ are characterized by the same properties of $`\alpha `$ relative to generic intervals $`(ϵ,ϵ),`$ resp. $`(\theta ,\theta ),`$ instead of $`(1,1).`$ We emphasize that this result is valid only in the massive case . The gauge group U$`(1)`$ acts on $``$ through $`e^{i\gamma }\text{Ad}\stackrel{~}{\mathrm{\Gamma }}(e^{i\gamma }),`$ and the self-adjoint unitary operator inducing the twisting is $`\stackrel{~}{\mathrm{\Gamma }}(\mathrm{𝟙}).`$
We refine Propositions 1 and 2 to suit our purpose . The computation of the statistics operator needs commutation rules between implementers when translated to mutually space-like regions. Let $`𝒪`$ be a double cone with basis $`(ϵ,ϵ)`$ at $`t=0,`$ and $`x,x^{}`$ such that $`𝒪+x^{}(𝒪+x)^{}.`$ Since the Schwinger term for the pair $`V(\lambda )_x,V(\lambda ^{})_x^{}`$ vanishes when $`𝒪+x^{}(𝒪+x)^{},`$ Proposition 2 establishes that the projective representation $`\stackrel{~}{\mathrm{\Gamma }}`$ is multiplicative in almost all cases of interest to us. Proposition 3 applies to the case $`U=U(n)_x`$ and $`e^{iA}=V(\lambda )_x^{},`$ giving
$$\stackrel{~}{\mathrm{\Gamma }}(V(\lambda )_x^{})\stackrel{~}{\mathrm{\Gamma }}(U(n)_x)=e^{i\pi n\lambda \text{sgn}(xx^{})}\stackrel{~}{\mathrm{\Gamma }}(U(n)_x)\stackrel{~}{\mathrm{\Gamma }}(V(\lambda )_x^{}),$$
$$\stackrel{~}{\mathrm{\Gamma }}(U(n)_x)\stackrel{~}{\mathrm{\Gamma }}(U(n^{})_x^{})=\stackrel{~}{\mathrm{\Gamma }}(U(n^{})_x^{})\stackrel{~}{\mathrm{\Gamma }}(U(n)_x).$$
As will be clear later, we consider only even charge $`n,`$ while the real number $`\lambda `$ is left arbitrary. Since $`𝒢_2`$ is a group, the product $`U(n)V(\lambda )`$ is unitarily implementable too. We set $`WW(n,\lambda ):=U(n)V(\lambda ).`$ Choosing the same generating function $`\vartheta =\epsilon `$ we obtain a unitary operator on $`L^2(,^2)`$ defined by:
$$W:=e^{i\pi (n+\lambda )\epsilon ()}e^{i\pi \lambda \epsilon ()}.$$
Here $`W`$ acts on both components of $`L^2(,^2)`$ as multiplication by two distinct functions. We note that $`U(n)W(n,0)`$ and $`V(\lambda )W(0,\lambda )`$. In order to determine the charge carried by the automorphism $`\rho __W`$ we must evaluate $`q(W).`$ The Fredholm index of $`W_{}`$ can be easily computed as an immediate application of Theorem 1, and it equals $`n`$. Note that $`V(\lambda )`$ belongs to the connected component of the identity, therefore $`\text{Ind}V(\lambda )=\text{Ind}(\mathrm{𝟙})=0`$, $`\lambda `$. Thus $`n`$ is the charge carried by $`W,`$ with no contribution from $`V(\lambda ).`$ (The charge is entirely transported by $`U(n)`$ while $`V(\lambda )`$ is neutral).
###### Proposition 4
Automorphisms of $``$ of the form $`\rho __W`$ induce, on restriction, automorphisms of the observable algebra $`𝒜`$.
Proof. Since any $``$-homomorphism between $`C^{}`$ algebras is continuous, the statement is an immediate consequence of the inclusion
$$\rho __W(𝒜(𝒪_1))𝒜(𝒪_1),𝒪_1𝒦.$$
In order to prove this, let us consider a unitary $`W(n,\lambda )`$ with generating function $`\epsilon `$ centred in $`𝒪𝒦.`$ (We assume that all double cones have base at $`t=0`$). We note that $`\rho __W((𝒪_1))(𝒪_1).`$ Indeed, if $`\text{supp}(f)B_{𝒪_1},`$ then $`\text{supp}(Wf)B_{𝒪_1}`$ too, hence :
$$\rho __W(\pi (\varphi (f)))\pi (\varphi (Wf))(𝒪_1).$$
Moreover, $`[\rho __W(A),\stackrel{~}{\mathrm{\Gamma }}(e^{i\gamma })]=0`$ for all $`A𝒜(𝒪_1)`$ and $`\gamma ,`$ since the adjoint actions $`\text{Ad}\stackrel{~}{\mathrm{\Gamma }}(e^{i\gamma })`$ and $`\text{Ad}\stackrel{~}{\mathrm{\Gamma }}(V)`$ commute between themselves in view of Proposition 3. The claim follows from the gauge invariance of $`A.`$
$`\mathrm{}`$
###### Proposition 5
Automorphisms of $`𝒜`$ defined as in Proposition 4 are localizable in double cones.
Proof. Let $`A:=\pi (\varphi (f)),\text{supp}(f)B_{𝒪_1},`$ where $`𝒪_1𝒪^{}.`$ Obviously, if $`x\text{supp}(f)`$ then $`(Wf)(x)=0.`$ If $`x\text{supp}(f)B_{𝒪_1},`$ then $`xB_𝒪`$ and $`\epsilon (x)=\pm 1,`$ and we have
$$(Wf)(x)=e^{\pm \pi i\lambda }f(x),$$
hence $`\rho __W(A)=\text{Ad}\stackrel{~}{\mathrm{\Gamma }}(e^{\pm \pi i\lambda })(A),A(𝒪_1).`$ Now, it is then evident that, if $`A𝒜(𝒪_1)(𝒪_1)^{\text{U}(1)},`$ then $`\rho __W(A)=A,`$ i.e. $`\rho __W`$ acts trivially on $`𝒜(𝒪^{}).`$
$`\mathrm{}`$
###### Proposition 6
For each unitary $`W`$ defined as above, $`\rho __W:=\text{Ad}\stackrel{~}{\mathrm{\Gamma }}(W)`$ defines a localizable and translatable automorphism of $`𝒜`$.
Proof. We have just proved localizability: $`\rho __W|_{𝒜(𝒪^{})}=\text{id}|_{𝒜(𝒪^{})},`$ where $`𝒪`$ is the double cone in whose base the unitary $`W`$ (i.e. its generating function $`\epsilon `$) is “centred”, and $`\rho __W(𝒜(\stackrel{~}{𝒪}))(𝒜(\stackrel{~}{𝒪}))`$ for every double cone $`\stackrel{~}{𝒪}𝒪.`$ In order to prove translatability, let us observe that denoting the translates by $`W_x:=T(x)WT(x)`$, the automorphism $`\rho _{_{W_x}}`$ is localized in $`𝒪+x`$. Moreover, the unitary $`\stackrel{~}{\mathrm{\Gamma }}(W_xW^{})`$ intertwines $`\rho __W:=\text{Ad}\stackrel{~}{\mathrm{\Gamma }}(W)`$ and $`\rho _{_{W_x}}:=\text{Ad}\stackrel{~}{\mathrm{\Gamma }}(W_x),`$ and induces an equivalence between them since it is a local observable. Indeed, $`\stackrel{~}{\mathrm{\Gamma }}(W_xW^{})𝒜(\stackrel{~}{𝒪}),`$ with $`\stackrel{~}{𝒪}𝒪𝒪_x`$. The gauge invariance comes from the commutation relation of Proposition 3 between $`\stackrel{~}{\mathrm{\Gamma }}(V)`$ and $`\stackrel{~}{\mathrm{\Gamma }}(e^{iA})`$, where the inner product (8) now vanishes. Indeed, since $`q(W_xW^{})=q(W_x)+q(W^{})q(W)q(W)=0,`$ it follows that the intertwiner $`\stackrel{~}{\mathrm{\Gamma }}(W_xW^{})`$ preserves the charge and so $`\stackrel{~}{\mathrm{\Gamma }}(W_xW^{})\mathrm{\Omega }\mathrm{\Omega }.`$ We thus have:
$$(\stackrel{~}{\mathrm{\Gamma }}(W_xW^{})\mathrm{\Omega },\text{d}\stackrel{~}{\mathrm{\Gamma }}(\gamma \mathrm{𝟙})\stackrel{~}{\mathrm{\Gamma }}(W_xW^{})\mathrm{\Omega })=(\mathrm{\Omega },\text{d}\stackrel{~}{\mathrm{\Gamma }}(\gamma \mathrm{𝟙})\mathrm{\Omega })=0.$$
In order to prove that $`\stackrel{~}{\mathrm{\Gamma }}(W_xW^{})(\stackrel{~}{𝒪^{}})^{},`$ let us consider $`𝒪_1\stackrel{~}{𝒪}^{}`$ and $`\text{supp}(f)B_{𝒪_1}`$. In order to evaluate the expression
$$\stackrel{~}{\mathrm{\Gamma }}(W_xW^{})\pi (\varphi (f))\stackrel{~}{\mathrm{\Gamma }}(W_xW^{})^{}=\rho _{_{W_x}}\rho __W^{}(\pi (\varphi (f))),$$
(12)
we notice that $`\rho __W^{}\rho __W^1`$ is still localized in $`𝒪,`$ and then, since $`W(n,\lambda )^{}=W(n,\lambda ),`$
$$\rho __W^{}(\pi (\varphi (f)))=\pi (\varphi (e^{i\pi \lambda }f)).$$
Analogously
$$\rho _{_{W_x}}(\pi (\varphi (f)))=\pi (\varphi (e^{i\pi \lambda }f)),$$
so the right side of (12) reduces to $`\pi (\varphi (f))`$ and the result follows from twisted duality.
$`\mathrm{}`$
In the previous proof we have incidentally established that unitaries $`\stackrel{~}{\mathrm{\Gamma }}(W)`$ are gauge invariant if and only if $`q(W)=0.`$
We are now in position to perform the computation of the statistics operator $`\epsilon _{\rho __W}.`$ For simplicity, we start with an automorphism $`\rho __W`$ localized in a double cone $`𝒪`$ centred at the origin.
###### Proposition 7
If the automorphism $`\rho __W`$ is localized in a double cone centred at the origin, its statistics operator is
$$\epsilon _{\rho __W}=e^{\pm 2\pi in\lambda }\mathrm{𝟙},$$
according to the connected component of $`𝒪^{}.`$
Proof. Since we work at $`t=0,`$ we omit the component-subscript and consider $`x.`$ The automorphism $`\rho _{_{W_x}}`$ is localized in $`𝒪+x`$ and is unitarily equivalent to $`\rho __W`$ through the intertwiner $`\stackrel{~}{\mathrm{\Gamma }}(W_xW^{})`$. Then,
$$\epsilon _{\rho __W}=\stackrel{~}{\mathrm{\Gamma }}(W_xW^{})^{}\rho __W(\stackrel{~}{\mathrm{\Gamma }}(W_xW^{}))=\stackrel{~}{\mathrm{\Gamma }}(W)\stackrel{~}{\mathrm{\Gamma }}(W_x)^{}\stackrel{~}{\mathrm{\Gamma }}(W)\stackrel{~}{\mathrm{\Gamma }}(W_x)\stackrel{~}{\mathrm{\Gamma }}(W^{})^2.$$
(Here, and in the sequel, we omit all cocyles since they are always coupled with their conjugate). With this convention the previous expression yields
$$\stackrel{~}{\mathrm{\Gamma }}(W)\stackrel{~}{\mathrm{\Gamma }}(W_x^{})\stackrel{~}{\mathrm{\Gamma }}(U)\stackrel{~}{\mathrm{\Gamma }}(V)\stackrel{~}{\mathrm{\Gamma }}(U_x)\stackrel{~}{\mathrm{\Gamma }}(V_x)\stackrel{~}{\mathrm{\Gamma }}(V)^{}\stackrel{~}{\mathrm{\Gamma }}(U)^{}\stackrel{~}{\mathrm{\Gamma }}(W)^{}.$$
By Proposition 3 and our remarks on the specific cases discussed in ,
$`\stackrel{~}{\mathrm{\Gamma }}(V)\stackrel{~}{\mathrm{\Gamma }}(U_x)\stackrel{~}{\mathrm{\Gamma }}(e^{iX(\lambda )})\stackrel{~}{\mathrm{\Gamma }}(U(n)_x)`$ $`=`$ $`e^{i(\stackrel{~}{\mathrm{\Gamma }}(U_x)\mathrm{\Omega },\text{d}\stackrel{~}{\mathrm{\Gamma }}(X(\lambda ))\stackrel{~}{\mathrm{\Gamma }}(U_x)\mathrm{\Omega })}\stackrel{~}{\mathrm{\Gamma }}(U_x)\stackrel{~}{\mathrm{\Gamma }}(e^{iX(\lambda )})`$
$`=`$ $`e^{i\pi n\lambda \text{sgn}(x)}\stackrel{~}{\mathrm{\Gamma }}(U_x)\stackrel{~}{\mathrm{\Gamma }}(V),`$
while $`\stackrel{~}{\mathrm{\Gamma }}(V)`$ commutes with $`\stackrel{~}{\mathrm{\Gamma }}(V_x),`$ and $`\stackrel{~}{\mathrm{\Gamma }}(U_x)^{}`$ with $`\stackrel{~}{\mathrm{\Gamma }}(U).`$ We then obtain
$$e^{i\pi n\lambda \text{sgn}(x)}\stackrel{~}{\mathrm{\Gamma }}(W)\stackrel{~}{\mathrm{\Gamma }}(V_x^{})\stackrel{~}{\mathrm{\Gamma }}(U)\stackrel{~}{\mathrm{\Gamma }}(V_x)\stackrel{~}{\mathrm{\Gamma }}(U)^{}\stackrel{~}{\mathrm{\Gamma }}(W)^{}.$$
Repeating the same arguments for the two central terms, one has
$$\epsilon _{\rho __W}=e^{2\pi in\lambda \text{sgn}(x)}\mathrm{𝟙}.$$
$`\mathrm{}`$
If the automorphism is translated to a double cone $`𝒪+x,`$ for arbitrary $`x,`$ it assumes the form $`\rho _{_{W_x}}`$, with $`W`$ localized around the origin. With a proof identical to that of Prop. 7, one easily obtains the following.
###### Proposition 8
For every $`x`$,
$$\epsilon _{\rho _{_{W_x}}}=e^{2\pi in\lambda \text{sgn}(xy)}\mathrm{𝟙},$$
with $`𝒪+y`$ the auxiliary double cone, spatially separated from $`𝒪+x`$, used in the construction of the statistics operator.
Remark. In view of the decomposition $`\rho __W=\rho __U\rho __V,`$ an alternative method of evaluating $`\epsilon _{\rho __W}`$ is based on the identity
$$\epsilon _{\rho __U\rho __V}=\rho __U(\epsilon (\rho __U,\rho __V))\epsilon _{\rho __U}\rho _{_U}^{}{}_{}{}^{2}(\epsilon _{\rho __V})\rho __U(\epsilon (\rho __V,\rho __U)).$$
(13)
A straightforward computation reduces (13) to $`\epsilon _{\rho __W}=\rho __U(\epsilon _M(\rho __U,\rho __V)),`$ where the monodromy operator is simply $`\epsilon _M(\rho __U,\rho __V)=e^{2\pi in\lambda \text{sgn}(yx)}\mathrm{𝟙}.`$
We also observe that we could have determined the latter by exploiting the low-dimensional quantum field theory as formulated in , where only the statistics phases are involved. Indeed, since $`\kappa __W=e^{2\pi in\lambda \text{sgn}(yx)}`$ and $`\kappa __U=\kappa __V=1,`$ the claim follows from
$$\epsilon (\rho __V,\rho __U)\epsilon (\rho __U,\rho __V)=\frac{\kappa __W}{\kappa __U\kappa __V}\mathrm{𝟙}$$
\[13, Lemma 3.3\]. Therefore, these results are still consistent with the general theory of local quantum fields in low dimension.
We have incidentally noticed that in a (1+1)-dimensional massive QFT the statistics of a product may not coincide with the product of statistics, i.e. composition of DHR morphisms with ordinary statistics may generate braid statistics. This possibility is excluded by (3+1)-dimensional QFT with Haag duality \[8, pag. 179\], where $`\epsilon _{\xi _1}\epsilon _{\xi _2}=\epsilon _{\xi _1\xi _2}`$ for two arbitrary superselection sectors $`\xi _1,\xi _2,`$ (i.e. equivalence classes of Poincaré covariant localized automorphisms). The factorization property of statistics is no longer true in theories where non ordinary statistics can occur. In our model this property is equivalent to the triviality of monodromy, i.e. if and only if the automorphism carries ordinary statistics. Since unitary intertwiners between translation equivalent automorphisms are local observables, the violation of the multiplicative property of statistics cannot be attributed to the violation of Haag duality but to the geometry of space-time.
It is easily seen that the statistics operator depends on the translation equivalent class of automorphisms but not on its representative. The following Corollary is then evident.
###### Corollary 1
The statistics operator $`\epsilon _{\rho __W}`$ gives rise to a one-dimensional representation of the braid group if and only if $`\mathrm{\hspace{0.17em}2}n\lambda .`$
We end this section with an expression for the statistics operator $`\epsilon (\rho __W,\rho __W^{})`$ when the unitaries $`W`$ and $`W^{}`$ are centred in the same interval. If the induced automorphisms are localized in $`𝒪,`$ let $`x,y^2`$ be such that $`𝒪_x`$ and $`𝒪_y`$ lie in the right component of $`𝒪^{},`$ with $`𝒪_x𝒪_y,`$ where by $`𝒪_x𝒪_y`$ we mean that $`𝒪_x`$ lies in the right component of the space-like complement of $`𝒪_y`$. Let $`\rho _{_{W_x}}`$ be localized in $`𝒪_x`$ and equivalent to $`\rho __W.`$ Analogously, let $`\rho _{_{W_y^{}}}`$ be localized in $`𝒪_y`$ and equivalent to $`\rho __W^{}.`$ One then has
$`\epsilon (\rho __W,\rho __W^{})`$ $`=`$ $`\stackrel{~}{\mathrm{\Gamma }}(W^{}W_{y}^{}{}_{}{}^{})\times \stackrel{~}{\mathrm{\Gamma }}(WW_x^{})\stackrel{~}{\mathrm{\Gamma }}(W_xW^{})\times \stackrel{~}{\mathrm{\Gamma }}(W_y^{}W^{})`$
$`=`$ $`\stackrel{~}{\mathrm{\Gamma }}(W^{}W_{y}^{}{}_{}{}^{})\stackrel{~}{\mathrm{\Gamma }}(W_y^{})\stackrel{~}{\mathrm{\Gamma }}(WW_{x}^{}{}_{}{}^{})\stackrel{~}{\mathrm{\Gamma }}(W_y^{})^{}\stackrel{~}{\mathrm{\Gamma }}(W_xW^{})\rho __W(\stackrel{~}{\mathrm{\Gamma }}(W_y^{}W^{}))`$
$`=`$ $`\stackrel{~}{\mathrm{\Gamma }}(W^{}W_{y}^{}{}_{}{}^{})\stackrel{~}{\mathrm{\Gamma }}(W_y^{})\stackrel{~}{\mathrm{\Gamma }}(WW_{x}^{}{}_{}{}^{})\stackrel{~}{\mathrm{\Gamma }}(W_y^{})^{}\stackrel{~}{\mathrm{\Gamma }}(W_y^{}W^{})\stackrel{~}{\mathrm{\Gamma }}(W_xW^{})`$
$`=`$ $`e^{i\pi (n\lambda ^{}+n^{}\lambda )}\stackrel{~}{\mathrm{\Gamma }}(W^{})\stackrel{~}{\mathrm{\Gamma }}(W)\stackrel{~}{\mathrm{\Gamma }}(W^{})^{}\stackrel{~}{\mathrm{\Gamma }}(W)^{}`$
$`=`$ $`e^{\pi i(n\lambda ^{}+n^{}\lambda )}\mathrm{𝟙}.`$
On the other hand, if we choose $`𝒪_x𝒪_y,`$ the exponent in the last member changes sign. Therefore
$$\epsilon (\rho __W,\rho __W^{})=e^{i\pi (n\lambda ^{}+n^{}\lambda )\text{sgn}(yx)}\mathrm{𝟙}.$$
(14)
## 4 Charge implementers are not quasilocal
As already stated, charge implementers $`\stackrel{~}{\mathrm{\Gamma }}(W)`$ do not belong to the observable algebra. In this section we will show that they are not even in the field algebra $``$ when $`\lambda 0mod2.`$ As a first step we observe that if $`\stackrel{~}{\mathrm{\Gamma }}(W),`$ then its translates $`\stackrel{~}{\mathrm{\Gamma }}(W)_x`$ too, through $`\alpha _x(\pi (\varphi (f)))=\pi (\varphi (\tau _xf)),`$ where $`(\tau _xf)(\xi ):=f(\xi x).`$ Once we have determined the statistics operator, we can exclude the trivial case $`\lambda =0`$ (no kinks present), since it yields ordinary statistics. Obviously, this is not the unique value of $`\lambda `$ for which the statistics reduce to the ordinary one. The other values which realize this possibility depend on the charge, since they are given by $`2n\lambda ,`$ and are trivially taken into account.
For $`x`$ such that $`𝒪+x𝒪^{},`$ we have
$$\stackrel{~}{\mathrm{\Gamma }}(W)\stackrel{~}{\mathrm{\Gamma }}(W_x)=e^{2\pi iq\lambda \text{ sgn}(x)}\stackrel{~}{\mathrm{\Gamma }}(W_x)\stackrel{~}{\mathrm{\Gamma }}(W),$$
(15)
where $`qq(W)=q(W_x)`$ is the charge carried by both $`\rho __W`$ and $`\rho _{_{W_x}}.`$ For a general field $`F`$ we have
$$F^\tau =F_+iF_{}\stackrel{~}{\mathrm{\Gamma }}(\mathrm{𝟙}),$$
(16)
where $`F_\pm `$ denote the bosonic, resp. fermionic, part of $`F.`$ This can be easily seen from the explicit form of the twisted field in the general case:
$$F^\tau =\text{Ad}Z(F),Z=\frac{\mathrm{𝟙}+i\stackrel{~}{\mathrm{\Gamma }}(\mathrm{𝟙})}{1+i}$$
(the symbol $`\mathrm{𝟙}`$ always denotes the identity operator on the corresponding Hilbert space, as is clear from the context). We are interested in the case $`2q\lambda ,`$ since this leads to a contradiction – the commutators are non vanishing, according to (15). Let $`\stackrel{~}{\mathrm{\Gamma }}(W).`$ There exists a sequence of local fields $`\{F_n\}_n`$ norm convergent to $`\stackrel{~}{\mathrm{\Gamma }}(W),`$ with $`F_n(𝒪_n),n.`$ Since $`\stackrel{~}{\mathrm{\Gamma }}(W)^\tau =\stackrel{~}{\mathrm{\Gamma }}(W),`$ we have
$$F_n^\tau F_n2\stackrel{~}{\mathrm{\Gamma }}(W)F_n.$$
(17)
On the other hand, according to (16), we have $`F_n^\tau F_n=(1+i)F_n^{}Z,`$ and thus $`F_n^\tau F_n=\sqrt{2}F_n^{}.`$ By virtue of (17) we then have
$$F_n^{}\sqrt{2}\stackrel{~}{\mathrm{\Gamma }}(W)F_n.$$
(18)
Let now $`M>0`$ be such that $`F_n<M.`$ An $`\epsilon /3`$ argument yields
$$[\stackrel{~}{\mathrm{\Gamma }}(W),\stackrel{~}{\mathrm{\Gamma }}(W)_x]<(M+1)\stackrel{~}{\mathrm{\Gamma }}(W)F_n+F_nF_{n,x}\stackrel{~}{\mathrm{\Gamma }}(W)_x\stackrel{~}{\mathrm{\Gamma }}(W),$$
(19)
where we have used the fact that the translations, being implemented by unitary operators, are norm-preserving. Here $`F_{n,x}`$ denotes the translate of $`F_n`$ by $`x,`$ for the generic $`n.`$ The second term on the right hand side of (19) is dominated by
$$F_nF_{n,x}F_{n,x}F_n+F_{n,x}F_n\stackrel{~}{\mathrm{\Gamma }}(W)_x\stackrel{~}{\mathrm{\Gamma }}(W).$$
Since $`F_n`$ and $`F_{n,x}`$ are local fields, $`[F_n,F_{n,x}]=2F_n^{}F_{n,x}^{}`$ when $`𝒪_n+x𝒪_n^{}.`$ In this case we obtain, using (18):
$$F_nF_{n,x}F_{n,x}F_n2F_n^{}F_{n,x}^{}{}_{}{}^{}4\stackrel{~}{\mathrm{\Gamma }}(W)F_n^2.$$
In the last step we have exploited the commutativity between the actions $`\alpha _^2`$ and $`\beta _{\text{U}(1)}`$ in order to yield $`F_{x}^{}{}_{}{}^{}=F_{}^{}{}_{x}{}^{}`$ for each quasilocal field $`F.`$
Let now $`\epsilon >0,`$ and let $`n_\epsilon `$ be a positive integer such that
$$\stackrel{~}{\mathrm{\Gamma }}(W)F_n<\frac{\epsilon }{3(M+1)},F_{n,x}F_n\stackrel{~}{\mathrm{\Gamma }}(W)_x\stackrel{~}{\mathrm{\Gamma }}(W)<\frac{\epsilon }{3}$$
for all $`nn_\epsilon .`$ The integer $`n_\epsilon `$ is independent of $`x`$. For $`n=n_\epsilon ,`$ let $`x>0`$ be such that $`𝒪_{n_\epsilon }+x𝒪_{n_\epsilon }^{}{}_{}{}^{}.`$ Then, the left hand side of (19) can be made arbitrarily small, contradicting (15).
This proof, though intuitive, does not work if $`2q\lambda ,`$ since no contradiction is obtained in the latter case. In particular, our arguments exclude the case $`q=0.`$ We will give an alternative proof which overcomes this impediment, based on twisted duality for the field algebra.
###### Theorem 2
For every $`\lambda 0mod2`$ the unitary implementers $`\stackrel{~}{\mathrm{\Gamma }}(W)`$ are not elements of the field algebra $`.`$ If $`\lambda =0mod2,`$ they are local elements of the field algebra.
Proof. If $`\stackrel{~}{\mathrm{\Gamma }}(W),`$ let $`\{F_n\}_n`$ be a sequence of local elements of $``$ norm convergent to $`\stackrel{~}{\mathrm{\Gamma }}(W),`$ with $`F_n(𝒪_n)`$ for suitable double cones $`𝒪_n.`$ Since $`F_n^\tau (𝒪_n)^\tau =(𝒪_{n}^{}{}_{}{}^{})^{},`$ $`[F_n^\tau ,A]=0,A(𝒪_{n}^{}{}_{}{}^{}),n.`$ If $`\widehat{𝒪}_n𝒪_{n}^{}{}_{}{}^{}`$, we have $`[F_n^\tau ,\pi (\varphi (f))]=0,\text{ supp}(f)B_{\widehat{𝒪}_n},n.`$ With this notation it follows that
$$\stackrel{~}{\mathrm{\Gamma }}(W)^\tau \pi (\varphi (f))\pi (\varphi (f))\stackrel{~}{\mathrm{\Gamma }}(W)^\tau $$
$$(\stackrel{~}{\mathrm{\Gamma }}(W)^\tau F_n^\tau )\pi (\varphi (f))+F_n^\tau \pi (\varphi (f))\pi (\varphi (f))F_n^\tau +$$
$$+\pi (\varphi (f))(F_n^\tau \stackrel{~}{\mathrm{\Gamma }}(W)^\tau )$$
$$2\pi (\varphi (f))\stackrel{~}{\mathrm{\Gamma }}(W)^\tau F_n^\tau +F_n^\tau \pi (\varphi (f))\pi (\varphi (f))F_n^\tau .$$
(20)
Without lost of generality, we consider only normalized functions. Then, since the correspondence $`f\pi _P(\varphi (f))`$ is isometric, $`\pi (\varphi (f))=1.`$ Fixing an arbitrary $`ϵ>0,`$ let $`n_ϵ`$ be such that $`\stackrel{~}{\mathrm{\Gamma }}(W)^\tau F_n^\tau <ϵ/2`$ for every $`nn_ϵ.`$ For $`n=n_ϵ`$ and $`\text{supp}(f)B_{\widehat{𝒪}_{n_ϵ}},`$ where $`\widehat{𝒪}_{n_ϵ}𝒪_{n_ϵ}^{}{}_{}{}^{},`$ the expression in the last line of (20) is less than $`ϵ.`$ Therefore, for each fixed $`q`$ and $`\lambda `$, $`\stackrel{~}{\mathrm{\Gamma }}(W)^\tau \pi (\varphi (f))\pi (\varphi (f))\stackrel{~}{\mathrm{\Gamma }}(W)^\tau `$ can be made arbitrary small by choosing a suitable function.
On the other hand, being $`\stackrel{~}{\mathrm{\Gamma }}(W)^\tau =\stackrel{~}{\mathrm{\Gamma }}(W),`$ the expression in the first line of (20) assumes the simpler form
$$e^{i\pi \lambda }\pi (\varphi (f))\stackrel{~}{\mathrm{\Gamma }}(W)\pi (\varphi (f))\stackrel{~}{\mathrm{\Gamma }}(W)=|e^{i\pi \lambda }1|,$$
(21)
in view of the CAR relations and the unitarity of $`\stackrel{~}{\mathrm{\Gamma }}(W).`$ Thus, we have a contradiction between (20) and (21), since the latter establishes that, for $`\lambda 0mod2,`$ the norm is a strictly positive constant. Finally, if $`\lambda =0mod2`$ the unitary implementers are local elements of $`.`$ Indeed, if $`W`$ is centred in $`𝒪,`$
$$\stackrel{~}{\mathrm{\Gamma }}(W)\pi (\varphi (f))\stackrel{~}{\mathrm{\Gamma }}(W)^{}=\pi (\varphi (f))$$
for every $`f`$ with support spatially separated from $`𝒪.`$ Therefore, by twisted duality, $`\stackrel{~}{\mathrm{\Gamma }}(W)(𝒪^{})^{}=(𝒪)^\tau .`$ On the other hand, since the twisting is involutive on any local field algebra, i.e. $`(𝒪)^{\tau \tau }=(𝒪),`$ $`\stackrel{~}{\mathrm{\Gamma }}(W)(𝒪)`$ follows from $`\stackrel{~}{\mathrm{\Gamma }}(W)^\tau \stackrel{~}{\mathrm{\Gamma }}(W).`$
$`\mathrm{}`$
This result shows that the strange behaviour of statistics for this model appears only when the implementers are not elements of the field algebra, confirming that there is no contradiction between what we expected from the general theory of superselections sectors and the peculiarities arising from this model. When the “solitonic” parameter $`\lambda `$ vanishes, the statistics is again trivial, i.e. conventional Bose-Einstein or Fermi-Dirac. Since only the zero charge $`\stackrel{~}{\mathrm{\Gamma }}(W)`$ are gauge invariant, the unique cases in which the implementers are observables (and local) are when $`\lambda 2`$ and $`q=0.`$ The classification of the localizability property of our implementers is thus complete.
## 5 Braiding Structure and Asymptotic Abelianness
In order to approach the study of strange statistics with a more general tool, we observe that the well known AQFT as formulated in , and cannot be applied here in its entirety, since Haag duality is violated by (1+1)-dimensional free massive Dirac field theory. A more appropriate setting seems to be that proposed in in order to construct symmetric tensor $`C^{}`$-categories in QED, since it extends to theories where intertwiners are not contained in the algebra where the endomorphisms act. Moreover, the endomorphisms are not necessary locally inner <sup>3</sup><sup>3</sup>3We recall that $`\stackrel{~}{\mathrm{\Gamma }}(W)𝒜`$ in almost all cases., but only in an asymptotic sense. Asymptotic abelianness yields a tensor $`C^{}`$-category starting directly from representations, without exploiting Haag duality. We set up notation and state definitions.
Let us introduce two sets of nets $`\rho __W𝒰_{\rho __W}`$ and $`\rho __W𝒱_{\rho __W}`$ for every object $`\rho __W`$. Each net consists of unitary intertwiners in $`(\rho __W,\rho _{_{W_{x_m}}})`$, where $`\rho _{_{W_{x_m}}}`$ tends pointwise in norm to the identity morphism on $`𝒜`$ for suitable sequences $`\{x_m\}_m.`$
###### Definition 2 (Asymptotic abelianness)
A field theory model satisfies asymptotic abelianness if, given intertwiners $`R(\rho __W,\rho __W^{}),`$$`S(\rho _{_{\stackrel{~}{W}}},\rho _{_{\stackrel{~}{W}^{}}})`$ and nets $`U_m𝒰_{\rho __W},U_m^{}^{}𝒰_{\rho __W^{}},V_n𝒱_{\rho _{_{\stackrel{~}{W}}}},`$ $`V_n^{}^{}𝒱_{\rho _{_{\stackrel{~}{W}^{}}}}`$,
$$U_m^{}^{}RU_m^{}\times V_n^{}^{}SV_n^{}V_n^{}^{}SV_n^{}\times U_m^{}^{}RU_m^{}0$$
(22)
in norm as $`m,m^{},n,n^{}\mathrm{}`$.
Finally, the sets of nets must be compatible with products: for each pair $`\rho __W,\rho __W^{}\mathrm{\Delta },`$ there exist $`U_m𝒰_{\rho __W}\text{ and }U_m^{}𝒰_{\rho __W^{}}\text{ such that }U_m\times U_m^{}𝒰_{\rho __W\rho __W^{}},`$ and similarly for $`𝒱.`$ Here $`\mathrm{\Delta }`$ denotes a semigroup of endomorphism of $`𝒜.`$ If the nets satisfy all these conditions, they give rise to a bi-asymptopia.
###### Theorem 3
** If $`\rho __W,\rho _{_{\stackrel{~}{W}}}\mathrm{\Delta }`$, then:
$$\epsilon (\rho __W,\rho _{_{\stackrel{~}{W}}}):=\underset{m,n\mathrm{}}{lim}V_n^{}\times U_m^{}U_m\times V_n$$
exists, is independent of $`U_m𝒰_{\rho __W}`$ and $`V_n𝒱_{\rho _{_{\stackrel{~}{W}}}},`$ and lies in $`(\rho __W\rho _{_{\stackrel{~}{W}}},\rho _{_{\stackrel{~}{W}}}\rho __W).`$ Moreover, if $`R(\rho __W,\rho __W^{}),`$$`S(\rho _{_{\stackrel{~}{W}}},\rho _{_{\stackrel{~}{W}^{}}})`$ and $`\rho _{_{\widehat{W}}}\mathrm{\Delta },`$ then:
$$\epsilon (\rho __W^{},\rho _{_{\stackrel{~}{W}^{}}})R\times S=S\times R\epsilon (\rho __W,\rho _{_{\stackrel{~}{W}}}),$$
$$\epsilon (\rho __W\rho _{_{\stackrel{~}{W}}},\rho _{_{\widehat{W}}})=\epsilon (\rho __W,\rho _{_{\widehat{W}}})\times \mathrm{𝟙}_{\rho _{_{\stackrel{~}{W}}}}\mathrm{𝟙}_{\rho __W}\times \epsilon (\rho _{_{\stackrel{~}{W}}},\rho _{_{\widehat{W}}}),$$
$$\epsilon (\rho __W,\rho _{_{\stackrel{~}{W}}}\rho _{_{\widehat{W}}})=\mathrm{𝟙}_{\rho _{_{\stackrel{~}{W}}}}\times \epsilon (\rho __W,\rho _{_{\widehat{W}}})\epsilon (\rho __W,\rho _{_{\stackrel{~}{W}}})\times \mathrm{𝟙}_{\rho _{_{\widehat{W}}}}.$$
We apply this method to our class $`\mathrm{\Delta }`$ of automorphism in the setting of the (1+1)-dimensional free massive Dirac field. In light of the computation performed in the previous section, it remains to verify the condition of asymptotic abelianness in order to have a bi-asymptopia and consequently a braiding for our category. Starting from (2+1)-dimensional models, where a space-like cone $`𝒞`$ and its opposite $`𝒞`$ are usually chosen as asymptotic localization regions in the definition of the families $`𝒰`$ and $`𝒱,`$ we extends the method to a (1+1)-dimensional space-time by choosing the two standard wedges $`W_\pm .`$ Since in (1+1)-dimensions there is a natural notion of right and left, hence of $`+\mathrm{}`$ and $`\mathrm{}`$, we set
$$𝒰_{\rho __W}:=\{(U_a)_a(\rho __W,\rho _{_{W_a}}),a+\mathrm{}\},$$
$$𝒱_{\rho __W}:=\{(V_b)_b(\rho __W,\rho _{_{W_b}}),b\mathrm{}\},$$
where the two families of nets are contained in $`W_+`$, $`W_{}`$ respectively. (This condition implies that the nets are contained in the causal complement of every bounded region for large values of the indexes). Let $`R:\rho __W\rho __W^{},`$ $`S:\rho _{_{\stackrel{~}{W}}}\rho _{_{\stackrel{~}{W}^{}}},`$ and $`z,\zeta `$ be defined by
$$R=z\stackrel{~}{\mathrm{\Gamma }}(W^{}W^{}),S=\zeta \stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}^{}\stackrel{~}{W}^{}).$$
As intertwiners $`U_m𝒰_{\rho __W},U_m^{}^{}𝒰_{\rho __W^{}},V_n𝒱_{\rho _{_{\stackrel{~}{W}}}},V_n^{}^{}𝒱_{\rho _{_{\stackrel{~}{W}^{}}}},`$ we set:
$`U_m`$ $`=`$ $`\lambda _m\stackrel{~}{\mathrm{\Gamma }}(W_{x_m}W^{}),x_m\stackrel{W_+}{}+\mathrm{}`$
$`U_m^{}^{}`$ $`=`$ $`\lambda _m^{}^{}\stackrel{~}{\mathrm{\Gamma }}(W_{x_m^{}^{}}^{}W^{}),x_m^{}^{}\stackrel{W_+}{}+\mathrm{}`$
$`V_n`$ $`=`$ $`\mu _n\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n}\stackrel{~}{W}^{}),y_n\stackrel{W_{}}{}\mathrm{}`$
$`V_n^{}^{}`$ $`=`$ $`\mu _n^{}^{}\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n^{}^{}}^{}\stackrel{~}{W}^{}),y_n^{}^{}\stackrel{W_{}}{}\mathrm{}`$
where $`\lambda _m,\lambda _m^{}^{},\mu _n,\mu _n^{}^{}`$ are defined as the scalar $`z,`$ for every $`m,m^{},n,n^{}.`$ To simplify notation, in the sequel $`x_m^{}`$ stands for $`x_m^{}^{}`$. The term of the sequence to which we refer will be clear from the context. Analogous simplifications will be adopted for $`y_n^{}^{},\lambda _m^{}^{}`$ $`\mu _n^{}^{}.`$ We now perform the computation of the limit in (22), which now assumes the form:
$$\lambda _m^{}\stackrel{~}{\mathrm{\Gamma }}(W_{x_m^{}}^{}W^{})z\stackrel{~}{\mathrm{\Gamma }}(W^{}W^{})\overline{\lambda }_m\stackrel{~}{\mathrm{\Gamma }}(W_{x_m}W^{})^{}\times $$
$$\times \mu _n^{}\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n^{}}^{}\stackrel{~}{W}^{})\zeta \stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}^{}\stackrel{~}{W}^{})\overline{\mu }_n\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n}\stackrel{~}{W}^{})^{}$$
$$\mu _n^{}\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n^{}}^{}\stackrel{~}{W}^{})\zeta \stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}^{}\stackrel{~}{W}^{})\overline{\mu }_n\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n}\stackrel{~}{W}^{})^{}\times $$
$$\times \lambda _m^{}\stackrel{~}{\mathrm{\Gamma }}(W_{x_m^{}}^{}W^{})z\stackrel{~}{\mathrm{\Gamma }}(W^{}W^{})\overline{\lambda }_m\stackrel{~}{\mathrm{\Gamma }}(W_{x_m}W^{})^{}=$$
$$=\lambda _m^{}\overline{\lambda }_mzD(W_{x_m^{}}^{}W^{},W^{}W^{})D(W_{x_m^{}}^{}W^{},WW_{x_m}^{})\stackrel{~}{\mathrm{\Gamma }}(W_{x_m^{}}^{}W_{x_m}^{})\times $$
$$\times \mu _n^{}\overline{\mu }_n\zeta D(\stackrel{~}{W}_{y_n^{}}^{}\stackrel{~}{W}^{},\stackrel{~}{W}^{}\stackrel{~}{W}^{})D(\stackrel{~}{W}_{y_n^{}}^{}\stackrel{~}{W}^{},\stackrel{~}{W}\stackrel{~}{W}_{y_n}^{})\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n^{}}^{}\stackrel{~}{W}_{y_n}^{})(),$$
(23)
where ($``$) denotes the cross product of the same terms in the inverse order. Here $`D(A,B)`$ is the cocycle of the projective representation $`\stackrel{~}{\mathrm{\Gamma }}`$, i.e. $`\stackrel{~}{\mathrm{\Gamma }}(W_1)\stackrel{~}{\mathrm{\Gamma }}(W_2)=D(W_1,W_2)\stackrel{~}{\mathrm{\Gamma }}(W_1W_2).`$ Collecting all scalars in a factor $`q`$, (23) reduces to
$$q\left(\stackrel{~}{\mathrm{\Gamma }}(W_{x_m^{}}^{}W_{x_m}^{})\times \stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n^{}}^{}\stackrel{~}{W}_{y_n}^{})\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n^{}}^{}\stackrel{~}{W}_{y_n}^{})\times \stackrel{~}{\mathrm{\Gamma }}(W_{x_m^{}}^{}W_{x_m}^{})\right)$$
(24)
so, to evaluate the asymptotic behaviour of (22) it suffices to compute the limit of the expression in parentheses. Since
$$\stackrel{~}{\mathrm{\Gamma }}(W_{x_m^{}}^{}W_{x_m}^{}):\rho _{_{W_{x_m}}}\rho _{_{W_{x_m^{}}^{}}},\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n^{}}^{}\stackrel{~}{W}_{}^{}{}_{y_n}{}^{}):\rho _{_{\stackrel{~}{W}_{y_n}}}\rho _{_{\stackrel{~}{W}_{y_n^{}}^{}}}$$
the expression in (24) becomes:
$$q(\stackrel{~}{\mathrm{\Gamma }}(W_{x_m^{}}^{}W_{x_m}^{})\stackrel{~}{\mathrm{\Gamma }}(W_{x_m})\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n^{}}^{}\stackrel{~}{W}_{y_n}^{})\stackrel{~}{\mathrm{\Gamma }}(W_{x_m})^{}$$
$$\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n^{}}^{}\stackrel{~}{W}_{y_n}^{})\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n})\stackrel{~}{\mathrm{\Gamma }}(W_{x_m^{}}^{}W_{x_m}^{})\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n})^{})=$$
$$=q^{}\left(\stackrel{~}{\mathrm{\Gamma }}(W_{x_m^{}}^{})\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n^{}}^{})\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n}^{})\stackrel{~}{\mathrm{\Gamma }}(W_{x_m})^{}\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n^{}}^{})\stackrel{~}{\mathrm{\Gamma }}(W_{x_m^{}}^{})\stackrel{~}{\mathrm{\Gamma }}(W_{x_m})^{}\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n}^{})\right),$$
where
$$q^{}:=q\overline{D(W_{x_m^{}}^{},W_{x_m}^{})D(\stackrel{~}{W}_{y_n^{}}^{},\stackrel{~}{W}_{y_n}^{})}.$$
By construction $`y_n^{}<x_m^{},y_n<x_m`$ for $`m,n`$ sufficiently large, then unitaries $`W_{x_m^{}}^{}`$ and $`\stackrel{~}{W}_{y_n^{}}^{}`$ are centred in disjoint intervals for $`m,n`$ sufficiently large. Equivalently, the corresponding implementers $`\stackrel{~}{\mathrm{\Gamma }}(W_{x_m^{}}^{})`$ and $`\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n^{}})`$ are localized in causally disjoint regions. We are then in a position to exploit the commutation rules between second quantization operators we have established in Sect. 2. For example,
$$\stackrel{~}{\mathrm{\Gamma }}(W_{x_m^{}}^{})\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n^{}}^{})=e^{i\pi (n\lambda ^{}+n^{}\lambda )}\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n^{}}^{})\stackrel{~}{\mathrm{\Gamma }}(W_{x_m^{}}^{}),$$
and analogously for $`W_{x_m}`$. Performing the substitutions, it turns out that (22) vanishes for large values of the indexes and thus, a fortiori, tends to zero. We have thus shown the following
###### Proposition 9
The subcategory of $`\text{End}(𝒜)`$ generated from $`\mathrm{\Delta }`$ admits a braiding structure $`\epsilon `$.
This model shows that the method of asymptotic abelianess, in the form stated in , cannot be applied to (1+1)-dimensional massive theories in order to obtain a generalized statistics operator, since the definition of bi-asymptopias is not consistent with the geometric peculiarity of space-time which may give rise to two distinct statistics operators. More precisely, we compare the braiding $`\epsilon `$ of Theorem 3 with the statistics operator $`\epsilon (\rho __W,\rho __W^{})`$ computed in the purely algebraic setting as in (14). If $`\{x_m\}_m`$ and $`\{y_n\}_n`$ are such that $`𝒪_{x_m}𝒪_{y_n}`$ for sufficiently large $`m`$ and $`n`$, one has
$$V_n^{}\times U_m^{}U_m\times V_n=$$
$$=\overline{\mu }_n\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}\stackrel{~}{W}_{y_n}^{})\times \overline{\lambda }_m\stackrel{~}{\mathrm{\Gamma }}(WW_{x_m}^{})\lambda _m\stackrel{~}{\mathrm{\Gamma }}(W_{x_m}W^{})\times \mu _n\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n}\stackrel{~}{W}^{})=$$
$$\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}\stackrel{~}{W}_{y_n}^{})\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n})\stackrel{~}{\mathrm{\Gamma }}(WW_{x_m}^{})\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n}^{})\stackrel{~}{\mathrm{\Gamma }}(W_{x_m}W^{})\stackrel{~}{\mathrm{\Gamma }}(W)\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{W}_{y_n}\stackrel{~}{W}^{})\stackrel{~}{\mathrm{\Gamma }}(W)^{},$$
which coincides with the expression of $`\epsilon (\rho __W,\rho _{_{\stackrel{~}{W}}})`$ gained by transporting $`\rho __W`$ and $`\rho _{_{\stackrel{~}{W}}}`$ resp. to $`\rho _{_{W_{x_m}}}`$ and $`\rho _{_{\stackrel{~}{W}_{y_n}}}.`$ Since the two nets of double cones are contained in distinct components of $`𝒪^{},`$ this is incompatible with the basic prescriptions of AQFT, since this procedure would be equivalent to not remaining in a fixed connected component! We remark that asymptotic abelianness requires that the two nets of double cones $`𝒪_m^𝒰`$ and $`𝒪_n^𝒱,`$ which appear in $`𝒰_{\rho __W}`$ and $`𝒱_{\rho __W},`$ do lie in distinct components of $`𝒪^{}`$, since only this configuration guarantees that $`𝒪_m^𝒰𝒪_n^𝒱`$ tends space-like to infinity as $`m,n\mathrm{}.`$ (There are no alternatives in (1+1)-dimension with no additional constraint). So, the theory of bi-asymptopias may not lead to true statistics when the Minkowski space is not at least (2+1)-dimensional and braid statistics occurs. As in the AQFT setting, we can always collect path connected bi-asymptopias into equivalence classes, but without additional prescription on the double limit in Definition 2, we could include objects which have no physical meaning, since they correspond in (1+1)-dimensions to working with both components of $`𝒪^{}`$ at the same time. A natural way to generalize this approach to bidimensional theories is to reformulate some definitions, giving a restricted notion of asymptotic abelianness appropriate to all dimensions. Instead of performing the double limit as in Definition 2 and Theorem 3, we choose a particular “direction”, for example the diagonal one, i.e. $`m=n,`$
$$\epsilon (\rho __W,\rho _{_{\stackrel{~}{W}}}):=\underset{n\mathrm{}}{lim}V_n^{}\times U_n^{}U_n\times V_n,$$
provided $`𝒪_{x_n}`$ and $`𝒪_{y_n}`$ are space-like separated for $`m`$ and $`n`$ sufficiently large. (Analogously for the definitions of asymptotic abelianness and bi-asymptopias). All properties continue to be valid, since the true reason for sending to infinity the two nets of double cones is to exploit morphisms which commute. For example, with the definition $`𝒪_n^𝒱:=𝒪_n^𝒰\pm \widehat{e}__1n,`$ all conditions are satisfied and each pair of nets lies in the wedge $`W_\pm `$ for large values of the indexes. (Here $`\widehat{e}__1`$ denotes the unit vector in the $`x_1`$-direction). In this way, the braiding structure arising from asymptotic abelianness of intertwiners coincides with the braiding computed for a Haag-dual net whose morphisms are all inner, and it gives rise to a true braided tensor $`C^{}`$-category in all other cases, e.g. the present one, (where the statistics operator now has a genuine, not simply formal, meaning of statistics). In other terms, we have excluded all cases in which intertwiners satisfy asymptotic abelianness, but the limits in Theorem 3 do not give rise to a statistics operator. The braiding induced by these bi-asymptopias is really non symmetric, since bi-asymptopias $`\{𝒰,𝒱\},`$ $`\{𝒱,𝒰\}`$ are not path connected. In (3+1) dimensions all particles exhibit ordinary statistics since we can choose $`𝒪_n^𝒱=𝒪_n^𝒰`$ when we deal with strictly localized morphisms. This ensures the possibility of changing continuously from $`\{𝒰,𝒱\}`$ to $`\{𝒱,𝒰\}`$ along a chain of double cones. In (2+1) dimensions, where cones give the better notion of localization, one can choose $`\rho _a`$ in such a way that $`a`$ tends to space-like infinity remaining in a space-like cone $`𝒞,`$ resp. $`𝒞,`$ for $`𝒰_\rho `$, $`𝒱_\rho ,`$ and it is always possible to interchange the two cones by a sequences of allowed moves. This is not possible in (1+1) dimensions, since $`𝒪^{}`$ is not connected and thus, a fortiori, is not arcwise connected. We remark that a distinction between $`\epsilon (\rho ,\sigma )`$ and $`\epsilon (\sigma ,\rho )^{}`$ for a generic pair of DHR morphisms cannot be achieved by interchanging the roles of $`𝒰`$ and $`𝒱,`$ since both nets of double cones are in the same connected component of $`𝒪^{}.`$ Hence, in two dimensions to each morphism (object) we must associate two bi-asymptopias, one for each side of $`𝒪^{}`$. A direct computation of the limits in Theorem 3 for each connected component separately, gives two distinct values, $`e^{\pm 2\pi in\lambda },`$ which coincide with those found before.
Remark. Although two arbitrary automorphisms of the form $`\rho __W`$ are always connected by a similarity transformation (i.e., $`\rho _{_{W_2}}=\text{Ad}\stackrel{~}{\mathrm{\Gamma }}(W_2W_1^{})\rho _{_{W_1}}`$), this does not imply that they are unitarily equivalent through a local element of the observable algebra, since the unitary intertwiner is not necessarily in $`𝒜.`$ More precisely, if $`(n_1,\lambda _1)(n_2,\lambda _2)`$, then $`\stackrel{~}{\mathrm{\Gamma }}(W_1W_2^{})𝒜`$ as a consequence of the previous observation and of group relations for unitaries $`W(,;\epsilon ),`$ i.e. $`W(n,\lambda )W(n^{},\lambda ^{})=W(n+n^{},\lambda +\lambda ^{}),W(n,\lambda )^{}=W(n,\lambda ).`$ On the other hand, products $`W_xW^{}`$ have a different behaviour, since
$$W(n,\lambda ;\epsilon )_xW(n,\lambda ;\epsilon )^{}=W(n,\lambda ;\tau _x\epsilon \epsilon )$$
(25)
and, as already stated, $`\stackrel{~}{\mathrm{\Gamma }}(W_xW^{})𝒜(\stackrel{~}{𝒪}),\text{ where }\stackrel{~}{𝒪}𝒪𝒪_x.`$ We emphasize that the notation in (25) may give rise to ambiguities, since the operator on the right hand side carries no charge even if $`n0,`$ due to the particular form of the generating function.
## Conclusions
In the setting of AQFT we have shown that a family of localized and transportable automorphisms of the observable algebra $`𝒜`$ exhibits non ordinary statistics. Inside each sector one has different braiding structures labelled by a solitonic parameter $`\lambda `$ which reflects the action of smeared-out kink operators carrying no charge.
Owing to non locality of charge implementers, statistics is not an invariant of the sector, as already known in some two-dimensional particle theories or in solitonic theories. On the underlying ordinary structure, smeared-out kink operators give rise to a continuous family of braided tensor categories in the sense of the theory of bi-asymptopias.
The results are consistent with AQFT, which must be handled carefully here when tackling problems arising from the non locality of unitary implementers, the violation of Haag duality and the topological peculiarity of (1+1)-dimensional space-time. Owing to the latter, some results of local field theory are no longer valid in a two-dimensional world, giving rise to a range of intermediate situations and strengthening the concept that for massive theories in (1+1) dimensions statistics is not an intrinsic characteristic of sectors a priori . In the present case, since Haag duality can be overcome by peculiarities of the model, strangeness of statistics has its origin in the fact that implementers do not lie in the field algebra.
The interpretation of the braiding structure of this model extends to the CAR algebra the constructive method exploited for the Weyl algebra. Since not all the braidings obtained in this way give rise to a notion of statistics compatible with the DHR analysis, but only those constructed from pairs of sets of nets which tend to the same space-like infinity, the method of bi-asymptopias can be carried over to (1+1)-dimensional space-time only if we add a compatibility condition. This kind of selection criterion reflects the “initial condition” which determines uniquely the statistics operator in the standard algebraic approach, i.e. trivialization of $`\epsilon (\rho __W,\rho __W^{})`$ for $`\rho __W^{}\rho __W`$ . The particles described by this model are “statistical schizons”, since the same sector allows “pseudo” statistical descriptions and they exist in the same Hilbert space either as bosons or as fermions or as proper anyons , i.e. in two-dimensional massive theories not only the spin but also the statistics is a convention.
### Acknowledgement
I am greatly indebted to the supervisor of my PhD thesis, S. Doplicher, for many helpful discussions and much encouragement. I would like to thank C. D’Antoni, J. Roberts and M. Gabriel for a careful reading of the manuscript, and B. Schroer for useful correspondence. It is a pleasure to thank A. Silva, who was the coordinator of the PhD program at the Department of Mathematics of “La Sapienza”, University of Rome, in the period when this work has been carried out. |
warning/0507/hep-ph0507127.html | ar5iv | text | # 1 Introduction
## 1 Introduction
We now have overwhelming evidence that ordinary matter accounts for a minute portion of what constitutes the Universe at large. Most impressive is the confirmation from the very recent WMAP data. Very interestingly, many extensions of the standard model whose primary aim was related to the Higgs sector and the mechanism of symmetry breaking do provide a good dark matter candidate. Very soon, with the energy frontier that will open up at the LHC, intense searches for this new physics with its associated dark matter candidate will be pursued in earnest.
Meanwhile, many astroparticle experiments are going on, and will be improved by the time the LHC runs, to detect dark matter particles. The problem, either for direct or indirect detection of dark matter outside the colliders, is that we do not control many astrophysical parameters. For indirect detection which is the result of the annihilation of a dark matter pair in, say, the galactic halo of our galaxy, the photon signal is cleaner than that of the charged positron and antiproton that are considered as sources of exotic cosmic rays. The photon will point back to the source while the antiproton flux suffers from uncertainties due to the propagation. Of course in both cases one still has to rely on a modelling of the dark matter profile since one needs to know the number density of the annihilating dark matter particles. A very distinctive signal though would be that of a “direct” annihilation into a monochromatic photon. In this case the spectrum will reveal a peak at an energy corresponding to the mass of the annihilating particles since the latter move at essentially zero relative velocity $`v`$. In the galactic halo $`v/c10^3`$. Therefore, provided one has a detector with good energy resolution, the flux from the “direct” annihilation will clearly stand out above the (astrophysical) background or the diffuse contribution. The latter is due essentially to annihilation into quarks and $`W`$ which subsequently fragment and radiate/decay into photons. This contribution has a continuous featureless energy distribution which is only cut-off at a maximum energy corresponding to the mass of the dark matter particle. There are, and there will be, many powerful detectors to search for such photon signals, covering a wide range in energy from MeV to TeV. These are either ground based, like the atmospheric Cerenkov telescopes, ACT,(Cangaroo, HESS, MAGIC, VERITAS,..) or space borne telescopes (EGRET, AMS, the upcoming GLAST,..). Many see in some of the present data an excess that might be a sign for New Physics and dark matter annihilation but we should probably be cautious and await confirmation from other more precise detectors covering the same energy range. One should also improve on the theoretical predictions and a better understanding of the background and the astrophysical component that enter the calculation of the photon yield.
Our aim in this paper is to revisit the calculation of the “direct” self-annihilation into $`\gamma \gamma `$, $`Z\gamma `$ and $`gg`$ of the lightest supersymmetric particle (LSP). This is a neutralino that we will denote as $`\stackrel{~}{\chi }_1^0`$. There have been a few attempts of calculating the one-loop induced $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$ before two complete calculations settled the issue. These calculations have been made in the limit $`v=0`$ as is appropriate for dark matter annihilation in the halo. A very recent calculation has also been made for this mode. Their results for $`v=0`$ agree in their most important features (higgsino limit, for example) with those of Refs , but as far as we are aware no systematic comparison has been performed. Much more important however is that there is, at the moment, only one calculation of $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$ (performed at $`v=0`$) despite the fact that new features appear in this computation. These features, as we will see, can not be a mere generalisation of the $`2\gamma `$ final state. We will in this paper calculate both $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$and $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$for any velocity and make a tuned comparison with the existing codes for $`v=0`$, DarkSUSY for $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$and $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$and PLATONdml for $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$. In $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$we have identified a new contribution not taken into account in. We will also show some results for $`v=0.5`$ for both processes.
As a by-product we will also compute the self-annihilation into gluons: $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0gg`$. This can be derived from $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$by only keeping the coloured particles and dealing properly with the colour structure. This process could contribute to, for example, the antiproton signal. We will see that our results for $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0gg`$completely agree with those of DarkSUSY and PLATONdmg for $`v=0`$ and with PLATONgrel for $`v=0.5`$.
As is known, the largest contributions to $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$and $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$, especially for large neutralino masses, occur when the neutralino is a wino or a higgsino. As first pointed out in the cross section times the relative velocity, $`\sigma v`$, for both modes, tends to an asymptotic constant value that scales as $`1/M_W^2`$ for $`v=0`$ and large LSP mass. This result which breaks unitarity is due to the one-loop treatment of a “threshold” singularity that is nonetheless regulated by $`M_W`$. It has very recently been, admirably, shown how to include the higher order corrections through a non-relativistic non-perturbative approach. The latter reveals the formation of bound states with zero binding energy that show up as sharp resonances that dramatically enhance the cross section for particular masses. We have therefore thought it worthwhile to study the one-loop derivation in these scenarios and see how one can match the non-perturbative regime. The reason we do this and the main reason we carry the calculation of $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$and $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$is that one needs a reliable code for the photon flux from self-annihilating neutralino LSP’s. These cross sections have been missing from micrOMEGAs that we have been developing for a very accurate derivation of the relic density in supersymmetry and which is currently being adapted to direct and indirect detection. The present paper only deals with the cross section calculation, leaving aside the astrophysical issues to the implementation and exploitation within micrOMEGAs. See however Ref. for a preliminary use of micrOMEGAs to indirect detection using some of the results of this paper.
The results presented in this paper constitute some of the first applications of a code for the automatic computation of one-loop processes in supersymmetry relevant both for the colliders and astrophysics, such as the problem at hand. Most crucial for the latter is a careful treatment of the loop integrals since for these applications the use of general libraries is not appropriate leading to division by zero because of the appearance of vanishing Gram determinants in the reduction of the tensor integrals. We will show how to easily circumvent this problem.
## 2 Set-up of the automatic calculation
Even in the standard model, one-loop calculations of $`22`$ processes involve hundreds of diagrams and a hand calculation is practically impracticable. Efficient automatic codes for any generic $`22`$ processes, that have now been exploited for many $`23`$ and even some $`24`$ processes, are almost unavoidable for such calculations. For the electroweak theory these are the GRACE-loop code and the package FormCalc based on FeynArts and LoopTools.
With its much larger particle content, far greater number of parameters and more complex structure, the need for an automatic code at one-loop for the minimal supersymmetric standard model is even more of a must. A few parts that are needed for such a code have been developed based on the package FeynArtsusy but, as far as we know, no complete code exists or is, at least publicly, available. Grace-susy is now also being developed at one-loop and many results exist. One of the main difficulties that has to be tackled is the implementation of the model file, since this requires that one enters the thousands of vertices that define the Feynman rules. On the theory side a proper renormalisation scheme needs to be set up, which then means extending many of these rules to include counterterms. When this is done one can just use, or hope to use, the machinery developed for the SM, in particular the symbolic manipulation part and most importantly the loop integral routines and tensor reduction algorithms.
The calculations that we are reporting here are based on a new automatic tool that uses and adapts modules, many of which, but not all, are part of other codes. We will report on this approach elsewhere. Here we will be brief. In this application we combine LANHEP (originally part of the package COMPHEP) with the FormCalc package but with an extended and adapted LoopTools. LANHEP is a very powerful routine that automatically generates all the sets of Feynman rules of a given model, the latter being defined in a simple and compact format very similar to the canonical coordinate representation. Use of multiplets and the superpotential is built-in to minimize human error. The ghost Lagrangian is derived directly from the BRST transformations. The LANHEP module also allows to shift fields and parameters and thus generates counterterms most efficiently. Understandably the LANHEP output file must be in the format of the model file of the code it is interfaced with. In the case of FeynArts both the generic (Lorentz structure) and classes (particle content) files had to be given. Moreover because we use a non-linear gauge fixing condition, the FeynArts default generic file had to be extended.
This brings us to the issue of the gauge-fixing. We use a generalised non-linear gauge adapted to the minimal supersymmetric model. The gauge fixing writes
$`_{GF}`$ $`=`$ $`{\displaystyle \frac{1}{\xi _W}}|(_\mu ie\stackrel{~}{\alpha }\gamma _\mu igc_W\stackrel{~}{\beta }Z_\mu )W^{\mu +}+\xi _W{\displaystyle \frac{g}{2}}(v+\stackrel{~}{\delta }h+\stackrel{~}{\omega }H+i\stackrel{~}{\kappa }\chi _3)\chi ^+|^2`$
$`{\displaystyle \frac{1}{2\xi _Z}}(.Z+\xi _Z{\displaystyle \frac{g}{2c_W}}(v+\stackrel{~}{ϵ}h+\stackrel{~}{\gamma }H)\chi _3)^2{\displaystyle \frac{1}{2\xi _\gamma }}(.\gamma )^2.`$
$`h`$ and $`H`$ are the CP-even physical Higgses, with $`h`$ denoting the lightest. $`\gamma `$ is the photon field and the masses of the charged and neutral weak bosons are related through $`M_W=M_Zc_W`$. The $`\chi `$’s are the Goldstone fields. The non-linear gauge fixing parameters are $`\stackrel{~}{\alpha }`$, $`\stackrel{~}{\beta }`$, $`\stackrel{~}{\delta }`$, $`\stackrel{~}{\omega }`$, $`\stackrel{~}{\kappa }`$, $`\stackrel{~}{ϵ}`$ and $`\stackrel{~}{\gamma }`$. The $`\xi `$ are the usual Feynman parameters. In our implementation the latter are set to $`\xi =1`$ not only to avoid very large expressions due to the “longitudinal” modes of the gauge bosons but most importantly so that high rank tensors for the loop integrals are not needed. Gauge parameter independence which is a non trivial check on the result of the calculation can be made through the non-linear gauge fixing terms. In many instances a particular choice of the non-linear gauge parameter may prove much more judicious than another. For the case at hand, $`\stackrel{~}{\alpha }=1`$, preserves $`U(1)_{\mathrm{em}}`$ gauge invariance which explains the vanishing of the $`W^+\chi ^{}\gamma `$ vertex. This will prove crucial for the calculation of $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$.
This brings us to the implementation of the loop integrals and their use in the most general application to radiative corrections in SUSY both for the colliders, for indirect detection and improvement of the relic density calculation beyond tree-level. In LoopTools for example, the tensor loop integrals are reduced recursively to a set of scalar integrals by essentially following the Passarino-Veltman procedure. The reduction involves solving a set of equations that brings in the inverse Gram determinant<sup>1</sup><sup>1</sup>1For a recent overview of the problem with the Gram determinant, see . We will however present, for the $`22`$ processes, a simple solution.. Although for applications to the colliders the latter only vanishes for exceptional points in phase space, for the indirect detection calculation of tensor integrals involving annihilating LSP’s with small relative velocity $`v`$, the Gram determinant is of order $`v^2`$. Therefore it vanishes exactly for $`v=0`$ or can get extremely small slightly away from this value rendering the calculation highly instable. In the Appendix we show how we dealt with this problem in an automatic implementation. In a nut-shell, we have used a segmentation procedure based on the fact that when some momenta are dependent like what occurs with $`v=0`$, a $`N`$-point function writes as a sum of $`N1`$ point functions. This also applies to the tensorial structures. This observation is not new (see for example) and has been used mostly in hand calculations. Some aspects of it may be remotely related to. The scheme also allows an expansion away from exactly vanishing Gram determinants.
A selection of diagrams that contribute to both $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$and $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$is shown in Fig. 1 (see also ). Diagrams of type a) in Fig. 1 are particularly important in the large wino and higgsino limit. In this limit the LSP and the internal chargino are of almost equal mass. If the $`W`$ mass can be neglected this leads to a threshold singularity.
Moving from $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$to $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$brings mixing effects in the loops. Most diagrams can be derived from $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$. There is however an important class that is only present for $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$as shown in Fig.2. This class of diagrams is missing in. It corresponds to the insertion of the $`\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_1^0\gamma `$ transition<sup>2</sup><sup>2</sup>2The radiative neutralino decay $`\stackrel{~}{\chi }_j^0\stackrel{~}{\chi }_i^0\gamma `$ is calculated in.. In a general gauge, the virtual transition would be gauge dependent and not ultraviolet finite. To remedy both these problems requires the $`\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_1^0\gamma `$ counterterm which is generated starting from the (tree-level) $`\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_1^0Z`$ vertex through a $`Z\gamma `$ one-loop transition. It is well known that the latter is gauge dependent, see for example. The counterterm requires the field normalisation $`\delta Z_{Z\gamma }^{1/2}`$. This field renormalisation constant in fact also induces, like in the standard model, $`(H,h)Z\gamma `$ vertices not present at tree-level. This induced vertices are also needed for the class of diagrams shown in Fig. 1, in particular those with (H,h) exchange. The full set of counterterms for $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$is in fact obtained from the tree-level $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0ZZ`$, replacing a $`Z`$ by a photon and inserting the $`\delta Z_{Z\gamma }^{1/2}`$ renormalisation constant. However, it is known that this renormalisation constant vanishes for $`\stackrel{~}{\alpha }=1`$. We have checked this property explicitly with our code. After this check has been made, $`\stackrel{~}{\alpha }=1`$ was set, since it considerably reduces the number of diagrams and most importantly allows to discard all the counterterm contributions. Further gauge parameter independence of the result for $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$was checked by varying the other non-linear gauge parameters that enter the calculation, namely $`\stackrel{~}{\beta },\stackrel{~}{\delta }`$ and $`\stackrel{~}{\omega }`$. When discussing our results for $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$we will weight the effect of the new class of diagrams shown in Fig. 2 against that of the full contribution, in doing so we will specialise to $`\stackrel{~}{\alpha }=1`$. As we will see these diagrams give a non negligible contribution especially for the Higgsino case.
The application to $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0gg`$is rather straightforward. This process confirms that our code handles the colour summation correctly. Keeping only one flavour of quark with charge $`Q_f`$, $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0gg`$can be derived from $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$through $`(N_cQ_f\alpha )^22\alpha _s^2`$ ($`N_c=3`$ is the number of colours, $`\alpha `$ is the electromagnetic fine structure constant and $`\alpha _s`$ the QCD equivalent).
## 3 Results and comparisons
We first check that our results are ultraviolet finite by changing the numerical value of the parameter $`1/ϵ`$ that controls a possible ultraviolet divergence. $`ϵ=4n`$, where $`n`$ is the dimensionality of space. We also check for gauge parameter dependence by varying the non-linear gauge parameters, namely $`\stackrel{~}{\alpha },\stackrel{~}{\delta }`$ and $`\stackrel{~}{\omega }`$ for $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$and $`\stackrel{~}{\beta },\stackrel{~}{\delta }`$ and $`\stackrel{~}{\omega }`$ with $`\stackrel{~}{\alpha }=1`$ for $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$. These checks are carried in double precision and show that, for all points we have studied, the results are consistent up to $`13`$ digits. It is important to maintain the relation $`M_W=c_WM_Z`$. If these parameters are taken as independent the gauge parameter independence is lost. We first discuss our results for $`6`$ representative scenarios that we think are good checks on different parts of the calculations and also because they reveal the most important characteristics of these cross sections. Moreover these scenarios also serve to perform tuned comparisons against codes that are publicly available. To achieve this it is best to feed the codes the same parameters. Comparisons that use as input high-scale values for some SUSY parameters that are run down through some Renormalisation Group Equation, RGE, package often need to specify an interface. Moreover most often the RGE codes are updated and one does not always have access to the same version to perform a tuned comparison. For all these scenarios the input parameters are defined at the electroweak scale and are: $`M_1`$ the $`U(1)`$ gaugino mass, $`M_2`$ the $`SU(2)`$ counterpart, $`\mu `$ the Higgsino “mass”, $`M_A`$ the pseudoscalar mass and $`m_{\stackrel{~}{f}}`$ the common sfermion mass. $`\mathrm{tan}\beta `$ is set to $`10`$. The sfermion trilinear parameter $`A_f`$ is set to zero for all sfermions but the stop, depending on the mass of the latter. Our Higgs masses here are tree-level Higgs masses, so we avoided points too close to any Higgs resonance and the issue of the implementation of the width. When our code for $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$,$`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$and $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0gg`$will be fully incorporated within micrOMEGAs, corrected Higgs masses and mixing angles will be properly implemented in a gauge invariant manner following. This improved implementation is of relevance only for $`v0`$ since the CP-even Higgses do not contribute when the neutralinos are at rest. When we refer to cross sections this would in fact refer to the cross section times the relative velocity $`v`$, $`\sigma v`$ expressed in $`\mathrm{cm}^3/s`$. In terms of $`v`$, the total invariant mass of the neutralinos is $`s=4m_{\stackrel{~}{\chi }_1^0}^2/(1v^2/4)`$.
The six scenarios have been chosen so as to represent different properties, like the gaugino/higgsino content for different masses of the neutralino. We made no attempt whatsoever to pick up points that lead to a good relic abundance in accord with WMAP. This said, for each point, we give the corresponding values of the relic density, extracted from micrOMEGAs. One should however keep in mind that different unconventional histories of the Universe<sup>3</sup><sup>3</sup>3A few possibilities are described in. could alter the usual thermal prediction.
* Scenario 1: “Sugra”. This reproduces a typical output from so-called mSUGRA scenarios, although the latter would not produce a common sfermion mass. The neutralino is mostly bino with mass around $`200`$GeV. The lightest chargino is a wino.
* Scenario 2: “nSugra”. The neutralino is quite light, about $`100`$GeV and it is essentially bino. Here the mSugra relation does not hold, rather $`M_2=4M_1`$.
* Scenario 3: “higgsino 1”. The neutralino is a light higgsino of about $`200`$GeV. The lightest chargino has a mass about $`6`$GeV away.
* Scenario 4: “higgsino 2”. The neutralino is a heavy higgsino of about $`4`$TeV. It is quite degenerate with the lightest higgsino-like chargino. The mass diference is about $`0.1`$GeV.
* Scenario 5: “wino 1”. The neutralino and lightest chargino are light, about $`200`$GeV. The mass difference on the other hand is extremely small $`0.01`$GeV.
* Scenario 6: “wino 2”. This is like the previous example but for TeV masses. The LSP is a wino of mass $`4`$TeV completely degenerate with the chargino.
Table 1 shows that the nature of the LSP and its mass critically determine its self-annihilation cross section to $`\gamma \gamma `$ and $`Z\gamma `$. The results for the different scenarios vary by $`6`$ orders of magnitude, especially for $`v=0`$. The bino-like LSP gives far too small cross sections that are unlikely to be observed as a $`\gamma `$-ray line. The largest cross sections $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$and $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$for a LSP mass up to $`4`$TeV are found for the wino-like LSP. Moreover in this case the signal is almost an order of magnitude stronger in $`Z\gamma `$ than in $`\gamma \gamma `$, however the two lines even for $`M_{\stackrel{~}{\chi }_1^0}200`$GeV are only $`10`$GeV away, even before any smearing is taken into account. For the wino case the contribution of the $`\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_1^0\gamma `$ transition to $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$is negligible for the two scenarios we display in Table 1. This is not the case of the higgsino scenarios (nor the bino-like for that matter) where this contribution could amount to a correction of more than $`30\%`$. It is also interesting to note, see later, that for the very heavy wino scenario the cross section drops very quickly as we increase the velocity. We will study the wino and higgsino case in more detail below.
For $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0gg`$, the LSP composition does not show dramatic differences in the cross sections. The largest are however found for a light wino and a light higgsino of $`200`$GeV.
Let us now turn to the comparisons, essentially for $`v=0`$, with the codes PLATON and DarkSUSY. For this fine-tuned comparison we have taken the same masses for the quarks and $`M_Z`$ as in DarkSUSY as well as for the electromagnetic and strong coupling. On the other hand we imposed $`M_W=M_Zc_W`$. Taking for example, $`M_W=80.33`$GeV, with all other parameters as in Table 1 not only gives gauge parameter dependent results, but in the wino case the LSP would turn out to be the chargino.
Table 1 shows that our results (for $`v=0`$) agree perfectly with those of PLATONdml as concerns $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$as well as with PLATONdmg as concerns $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0gg`$. PLATONgrel also perfectly confirms our results for $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0gg`$for $`v=0.5`$. Excellent agreement with DarkSUSY is also observed (at $`v=0`$) for $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$and $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0gg`$. To compare with the results of DarkSUSY for $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$we do not consider the contribution from the $`\stackrel{~}{\chi }_i\stackrel{~}{\chi }_1^0\gamma `$ “insertion”. With this restriction we find exactly the same results as DarkSUSY in the case of the wino and higgsino but not in the case of the bino, where our results are about $`30\%`$ higher. However as pointed out earlier, the cross sections in the bino case are tiny.
The effect of the new contribution from $`\stackrel{~}{\chi }_i\stackrel{~}{\chi }_1^0\gamma `$ for both $`v=0`$ and $`v=0.5`$ is not noticeable when the neutralino is a pure wino, but it can be important in the higgsino case where the total contribution is also large. The new contribution brings in a relatively small correction in the bino case, where the cross sections are tiny anyhow.
Results for $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$and $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$for $`v=0.5`$ have been computed here for the first time and can be relevant for the computation of the relic density in some regions of the parameter space.
## 4 The wino and higgsino limits
The results of Table 1 make it clear that most interesting scenarios for the monochromatic $`\gamma `$ ray line signals are of a wino and higgsino type even when the LSP has a mass of about $`2M_Z`$. Fig. 3 shows the dependence of the $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$and $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$cross sections at $`v=0`$ as a function of the LSP mass in the case of a wino and a higgsino LSP. The mass of the LSP is in the range $`70`$GeV to $`100`$TeV. In fact masses below $`100`$GeV may be excluded by LEP2 but it is interesting to see how the cross sections grow past the $`100`$GeV mass to stabilise around a plateau. The masses of the other supersymmetric particles are taken extremely heavy here. Note that in the higgsino case $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$shows much more structure. The peak cross section is much more pronounced before the cross section decreases and reaches a plateau of the same order as $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$. The wino cross section for $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$is the largest of all and is almost an order of magnitude larger than $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$and two orders of magnitude larger than in the higgsino case.
It is also interesting to see how the plateau is reached for a fixed mass of the wino and higgsino LSP, or rather a fixed value of $`M_2`$ and $`\mu `$, depending on the composition of the LSP. We therefore fix these two values and vary the other supersymmetric parameters of the neutralino sector. The behaviour of the cross section as we vary these parameters is shown in Fig. 4. For the wino case one can see that once $`M_1,\mu >M_2`$, and therefore the LSP is mostly wino, the asymptotic values are reached abruptly especially for the case of a wino of $`10`$TeV and large $`\mathrm{tan}\beta `$. Below this transition, the cross sections have a smooth behaviour. In the higgsino limit, a fast transition occurs once $`M_1,M_2>\mu `$ but past this threshold there is still a smooth and slow increase of the cross sections before the asymptotic values are reached.
Most of this behaviour can, in fact, be recovered through simple analytical expressions that serve also as a further check on our results and the accuracy of the calculation in these extreme scenarios. It had been observed that when the LSP is heavy, much heavier that the $`W`$-boson, the cross sections (times velocity) for $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$and $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$tend to an asymptotic value that can be computed from the dominant contribution, that of diagram a) of Fig. 1. The limiting behaviour can be easily understood from the fact that in the heavy mass limit, the annihilating LSP neutralino and the chargino are degenerate with a mass much larger than the weak boson mass. This develops a threshold singularity like what one finds in QED, although here $`M_W`$ acts as a regulator. Another important factor that measures how the asymptotic values are reached is the deviation from exact degeneracy between the LSP neutralino and the lightest chargino given by their mass difference, $`\delta m`$, $`\delta m=m_{\stackrel{~}{\chi }_1^+}m_{\stackrel{~}{\chi }_1^0}`$. This scenario has been revisited in a series of excellent papers where it has been shown how the one-loop calculation in these cases need to be improved through a non-relativistic treatment. Our aim here, in the rest of this section, is to see how our one-loop results can be made to match with the non-perturbative treatment. This paves the way to an implementation in a code for indirect detection that can be used in all generality like what we have started to do in micrOMEGAs.
First, we will show how our one-loop results effectively capture the behaviour of the cross sections in these scenarios and how the asymptotic value in the case of a wino is reached dramatically fast.
In the higgsino limit, $`\mu M_1,M_2`$. We will also take $`M_2=2M_1=2M_S`$ and consider also the large $`\mathrm{tan}\beta `$ (in fact $`\mathrm{tan}\beta >2`$ suffices). $`\mu `$ will be taken positive.
In the wino limit, $`M_2\mu ,M_1`$. We will also take $`\mu =M_1=M_S`$ and large $`\mathrm{tan}\beta `$. The (tree-level) mass difference in the higgsino, $`\delta m^{\stackrel{~}{h}}`$, and wino limit, $`\delta m^{\stackrel{~}{w}}`$, write
$`\delta m^{\stackrel{~}{h}}`$ $``$ $`{\displaystyle \frac{m_Z^2}{2M_2}}c_W^2(1\mathrm{sin}2\beta )+{\displaystyle \frac{m_Z^2}{2M_1}}s_W^2(1+\mathrm{sin}2\beta ){\displaystyle \frac{5}{16}}{\displaystyle \frac{M_Z^2}{M_S}},`$
$`\delta m^{\stackrel{~}{w}}`$ $``$ $`{\displaystyle \frac{m_Z^4}{M_1\mu ^2}}s_W^2c_W^2\mathrm{sin}^22\beta {\displaystyle \frac{M_Z^2M_W^2}{M_S^3}}{\displaystyle \frac{1}{\mathrm{tan}\beta ^2}}.`$ (1)
We see that in the wino case the mass difference scales like $`1/M_S^3`$<sup>4</sup><sup>4</sup>4It is important to note that it is essential to have $`M_W=M_Zc_W`$, otherwise we could get a mass difference $`1/M_S`$. compared to the $`1/M_S`$ in the higgsino case. In these configurations the cross sections are well approximated by $`\stackrel{~}{\sigma }^{V\gamma ,\stackrel{~}{h}}v`$ in the higgsino case and $`\stackrel{~}{\sigma }^{V\gamma ,\stackrel{~}{w}}v`$ in the wino case ($`V=Z,\gamma `$) which are the results of the dominant diagrams of Fig. 1-a),
$`\stackrel{~}{\sigma }^{V\gamma ,\stackrel{~}{h}}v`$ $`=`$ $`\sigma _{\mathrm{}}^{V\gamma ,\stackrel{~}{h}}v\left(1+\sqrt{{\displaystyle \frac{2m_{\stackrel{~}{\chi }_1^0}\delta m}{M_W^2}}}\right)^2=\sigma _{\mathrm{}}^{V\gamma ,\stackrel{~}{h}}v\left(1+\sqrt{{\displaystyle \frac{5\mu }{6M_S}}}\right)^2,`$
$`\stackrel{~}{\sigma }^{V\gamma ,\stackrel{~}{w}}v`$ $`=`$ $`\sigma _{\mathrm{}}^{V\gamma ,\stackrel{~}{w}}v\left(1+\sqrt{{\displaystyle \frac{2M_Z^2M_2}{M_S^3\mathrm{tan}\beta ^2}}}\right)^2;V=Z,\gamma .`$ (2)
where the asymptotic values ($`\delta m=0`$) are given by
$`\sigma _{\mathrm{}}^{\gamma \gamma ,\stackrel{~}{h}}v`$ $`=`$ $`{\displaystyle \frac{\alpha ^4\pi }{4M_W^2s_W^4}}10^{28}cm^3/s,`$
$`\sigma _{\mathrm{}}^{\gamma \gamma ,\stackrel{~}{w}}v`$ $`=`$ $`16\sigma _{\mathrm{}}^{\gamma \gamma ,\stackrel{~}{h}}v\mathrm{1.6\hspace{0.33em}10}^{27}cm^3/s,`$ (3)
$`\sigma _{\mathrm{}}^{Z\gamma ,\stackrel{~}{h}}v`$ $`=`$ $`2{\displaystyle \frac{(1/2s_W^2)^2}{s_W^2c_W^2}}\sigma _{\mathrm{}}^{\gamma \gamma ,\stackrel{~}{h}}v\mathrm{0.8\hspace{0.33em}10}^{28}cm^3/s,`$ (4)
$`\sigma _{\mathrm{}}^{Z\gamma ,\stackrel{~}{w}}v`$ $`=`$ $`2{\displaystyle \frac{c_W^2}{s_W^2}}\sigma _{\mathrm{}}^{\gamma \gamma ,\stackrel{~}{w}}v10^{26}cm^3/s.`$ (5)
We have verified that our code including the complete contributions agrees extremely well with these approximation for the cross sections, Eq. 4, and that moreover in the wino case the asymptotic values, Eq. 4, are reached very fast due to very small $`\delta m`$. This is also exemplified in Fig. 4.
As demonstrated in the one-loop treatment of the threshold singularity that is responsible for the behaviour of these cross sections in the higgsino and wino regime at high LSP mass is not adequate an breaks unitarity. The non-perturbative non-relativistic approach of Ref. not only improves on the calculation but it also unravels the formation of bound states that drastically enhance the annihilation cross sections for specific combinations of masses. Fig. 5 shows the effect of such resonances and the departure of the cross section $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$from a full one-loop treatment as the mass of the higgsino LSP increases. The “resonance” curve is based on the use of a fitting function as given in <sup>5</sup><sup>5</sup>5We thank J. Hisano and M. Nojiri for confirming that Eq. 61 of hep-ph/0412403 should be squared and that the entry $`i=0,j=1`$ in Table 1 of that paper is 10 times smaller.. All other particles are taken heavy apart from the higgsino mass parameter $`\mu `$ and $`M_1=M_S=M_2/2=25`$TeV. For the whole range of $`M_{\stackrel{~}{\chi }_1^0}\mu `$ we have $`\delta m=0.1`$GeV. The figure also shows the value of the approximate one-loop result as given for the higgsino limit in Eq. 4, together with the full one-loop treatment based on our calculation. The resonance formation, here around $`6`$TeV, brings an enhancement factor of more than $`4`$ orders of magnitude. On the other hand departure from the full one-loop calculation is of relevance only for $`M_{\stackrel{~}{\chi }_1^0}`$ masses around $`1`$TeV. The insert in Fig. 5 shows in more detail the comparison of the full calculation compared to the approximate result for the smaller higgsino LSP masses, well before the resonance effects settle in. For $`M_{\stackrel{~}{\chi }_1^0}>600`$GeV the approximation is very good, only around $`M_{\stackrel{~}{\chi }_1^0}200`$GeV, the full calculation captures the effect of other contributions like those of Fig. 1b,e). For this particular case it looks like a good matching between the full one-loop result and the non perturbative one should be made at around $`M_{\stackrel{~}{\chi }_1^0}=400`$GeV. A possible strategy for choosing this matching point would have to rely on the knowledge of both the full one-loop result, the approximate one-loop result as given Eq. 4 and the non-perturbative results based on the fitting functions of Ref. . This would, of course, only be carried out in the limit of almost pure higgsino or wino. We would then have to compare the three results. To revert to the non-perturbative regime means that the full one-loop and the approximate one-loop agree fairly well and are quite different from the non-perturbative regime. If, on the other hand, these two one-loop results differ sensibly this means that one is not quite in the asymptotic region and that we might be missing some one-loop contributions. If this is the case one should also expect the higher order effects computed for the threshold region to be small so that the non-perturbative result and the approximate one-loop are very similar. Of course, as shown in the example of Fig. 5 these differences in the higgsino region, compared to taking the perturbative parameterisation, are rather small compared to the uncertainty that is inherent in the astrophysics part of the prediction of the gamma ray line. For the wino, as we saw, the transition to the asymptotic value is rather drastic especially for TeV LSP’s, therefore one should quickly capture the non-perturbative regime. Especially in this case one should also revert to a one-loop use of the chargino-LSP mass difference.
## 5 Conclusions
There has been a flurry of activity in the last few years in the search of dark matter with, among other strategies, many experiments dedicated to the indirect searches of dark matter in particular the gamma ray signal. The mono-energetic gamma ray line signal constitutes a clean signature. The improvement in coverage and accuracy of the measurements should be matched by precise theoretical calculations that should be publicly available through general purpose codes. In this paper we have provided a new calculation for $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$and $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$for the annihilation of the supersymmetric dark matter candidate and rederived as a bonus the $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0gg`$rate. These calculations have been made both for small (zero) relative velocity of the neutralinos as adequate for annihilation in the halo of our galaxy for example, but also for velocities that would be needed for the contribution of these channels in a precise derivation of the relic density. For $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$and $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$at $`v0`$ as would be needed for an improved relic density prediction, these results are new. For $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$three full one-loop calculations performed for $`v=0`$ have already been performed. We have performed a tuned comparison with the results of DarkSUSY and PLATON and have found perfect agreement. The calculation of $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$is trickier and can not just be deduced from $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$. Until now there has been only one calculation of this process. The latter has missed some contributions that may not always be negligible. Comparisons of our results with the previous ones without these contributions are quite good for scenarios where the cross section is not small. In this paper we have not made an attempt to fold in with the astrophysics part that involves, for example, the halo profile but concentrated on the particle physics part which must be an unambiguous prediction. To pave the way for an implementation in micrOMEGAs we felt it was important to critically review the large mass higgsino and wino LSP scenarios especially that the latter give large cross sections. As shown in one needs to go beyond the one-loop treatment in this regime. We have argued how one could match the full one-loop treatment with the non-perturbative result.
Another important aspect of this paper is the way all these calculations have been performed in a unified manner and the techniques that we used. These processes are the first application of a code for the calculation of one-loop processes in supersymmetry both at the colliders and for astrophysics/relic density calculations that require also a new way of dealing with the reduction of the tensor integrals. The calculations are performed with the help of an automatised code that allows gauge parameter dependence checks to be performed. The use of a generalised non-linear gauge is crucial.
Acknowledgment
We acknowledge discussions with our other colleagues of the micrOMEGAs Project in particular G. Bélanger who made some early comparisons with DarkSUSY and A. Pukhov who made some test runs with the $`\gamma \gamma `$ and $`Z\gamma `$ annihilations routines within micrOMEGAs. P. Brun and S. Rosier-Lees also made some test runs for the gamma ray predictions using a private version of micrOMEGAs for indirect detection and we gratefully thank them here. F.B. would like to thank the members of the Minami-Tateya group for important discussions and a fruitful collaboration. He also acknowledges discussions on the Gram determinant issue with T. Binoth, S. Dittmaier, J.P. Guillet and E. Pilon. F. Renard has been, as always, of invaluable help. We also thank J. Hisano and M. Nojiri for confirming that Eq. 61 of hep-ph/0412403 should be squared and that the entry $`i=0,j=1`$ in Table 1 of that paper is 10 times smaller. P. Gondolo has promptly provided us with the results of DarkSUSY reported in Table 1 and explained how to enter the parameters within this code for our tuned comparison. Fig. 1 and Fig. 2 are, in part, generated with JaxoDraw. This work is supported in part by GDRI-ACPP of the CNRS (France). The work of A.S is supported by grants of the Russian Federal Agency of Science NS-1685.2003.2 and RFBR 04-02-17448 and that of D.T is supported, in part, by a postdoctoral grant from the Spanish Ministerio de Educacion y Ciencia.
## A Appendix: Segmentation of loop integrals
The tensor integral of rank $`M`$ corresponding to a $`N`$-point graph, $`\{M,N\}`$, that we encounter in the general non-linear gauge but with Feynman parameters $`\xi =1`$ are such that $`MN`$. For the evaluation of $`22`$ processes $`N=4`$ is the maximum value and corresponds to the box. The general tensor integral writes as
$`T_{\underset{M}{\underset{}{\mu \nu \mathrm{}\rho }}}^{(N)}={\displaystyle \frac{d^nl}{(2\pi )^n}\frac{l_\mu l_\nu \mathrm{}l_\rho }{D_0D_1\mathrm{}D_{N1}}},MN,`$ (A. 1)
where
$`D_i=(l+s_i)^2M_i^2,s_i={\displaystyle \underset{j=1}{\overset{i}{}}}p_j,s_0=0.`$ (A. 2)
$`M_i`$ are the internal masses, $`p_i`$ the incoming momenta and $`l`$ the loop momentum.
The $`N`$-point scalar integrals correspond to $`M=0`$. All higher rank tensors for a $`N`$-point function, $`M1`$, can be deduced recursively from the knowledge of the $`N`$-point (and lower) scalar integrals. In Looptools this is based on the Passarino-Veltman algorithm. In Grace-loop the implementation is outlined in Ref. . The tensor reduction involves solving, recursively, a system of equations which explicitly requires the evaluation of the Gram determinant: $`DetG(p_1,p_2,p_3)=DetG_{123}=Detp_ip_j`$. For special kinematics the latter vanishes or can get very small, leading to severe numerical instability. This special kinematics for the general $`22`$ process one encounters in high energy occurs for exceptional points in phase space, for instance in extremely forward regions and most generally the weight of this contribution may be dismissed. For the case at hand, when the two neutralinos are at rest, or with extremely low relative velocity, the Gram determinant vanishes for all points because the incoming neutralinos have the same momentum and can not be considered independent. This is exactly what occurs in our case. Indeed here, the box diagrams for $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0Z\gamma `$have the Gram determinant
$`DetG=M_{\stackrel{~}{\chi }_1^0}^6v^2{\displaystyle \frac{\mathrm{sin}^2\theta }{(1v^2/4)^3}}(1z^2),z^2={\displaystyle \frac{M_Z^2}{4M_{\stackrel{~}{\chi }_1^0}^2}}(1v^2/4),`$ (A. 3)
$`v`$ is the relative velocity and $`\theta `$ is the scattering angle. For $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\gamma \gamma `$, $`z=0`$. In our application the sub-determinant with the incoming LSP neutralinos is responsible for the vanishing of the box Gram determinant for all angles:
$`DetG(p_1,p_2)=M_{\stackrel{~}{\chi }_1^0}^4v^2{\displaystyle \frac{1}{(1v^2/4)^2}}.`$ (A. 4)
This also means that the reduction of the tensor integrals for triangles with the two LSP as external legs, needs special treatment. Such triangles are of the type Fig. 1.e) (but not Fig.1 d)).
We therefore need to implement a routine for cases when the determinant vanishes due to the fact that the momenta $`s_i,i0`$ are not independent. There are a few ways of dealing with the tensor integrals when the Gram determinant is exactly zero. Sometimes, the problem is even avoided by the grouping of terms so that this spurious inverse determinant cancels out. Our aim was to find, at least for $`22`$ process, an efficient method that can, not only be easily automated, but also calls most of the existing routines that are present for a general purpose algorithm designed for non zero Gram determinants. Take the box for example. Observe that (in any $`n`$ dimension), in most generality, we may write for any given pair of constants $`\alpha ,\beta `$
$`{\displaystyle \frac{1}{D_0D_1D_2D_3}}`$ $`=`$ $`({\displaystyle \frac{1}{D_0D_1D_2}}\alpha {\displaystyle \frac{1}{D_0D_2D_3}}\beta {\displaystyle \frac{1}{D_0D_1D_3}}+(\alpha +\beta 1){\displaystyle \frac{1}{D_1D_2D_3}})\times `$
$`{\displaystyle \frac{1}{A+2l.(s_3\alpha s_1\beta s_2)}}`$
$`A`$ $`=`$ $`(s_3^2M_3^2)\alpha (s_1^2M_1^2)\beta (s_2^2M_2^2)(\alpha +\beta 1)M_0^2.`$ (A. 5)
Obviously if $`s_3=\alpha s_1+\beta s_2`$ and hence the momenta are linearly dependent, the box splits into a sum of triangles. We will refer to this as segmentation. This segmentation is independent of the tensor structure. This means that the reduction of the tensor box with zero Gram determinant amounts to a tensor reduction for triangle with a non-zero Gram determinant for which one uses the usual procedure and hence uses the general library. Observe that if $`\alpha =0`$ or $`\beta =0`$, there are three segments instead of four. The missing segment, triangle integral, does in fact have a zero Gram determinant. Therefore when one approaches the zero of the Gram determinant of the box in these specific cases, $`\alpha 0`$ for example, $`\alpha `$ will be numerically very small but non zero. A numerical instability could still develop due now to the Gram determinant of the associated triangle. These “algebraic” zeros could be missed at the numerical level. This can again lead to (less severe) numerical instabilities due to the reduction of the associated tensor. In this case, these triangles are segmented even further into two-point functions, following the same recipe. This way their contribution is negligible even at the numerical, automatic, level. In any case, since we also encounter (tensor) triangle diagrams (Fig. 1.e)) that have a vanishing Gram determinant, we have also included a segmentation that also applies to the tensor triangles using the same trick as for the boxes.
There is an important observation to make. The segmentation of a tensor of rank $`M`$ for the $`N`$-point function, $`\{M,N\}`$, amounts to applying the tensor reduction on $`\{M,N1\}`$. If $`M=N`$, after segmentation one would need a library for $`\{N,N1\}`$. These are not supplied by default in the general libraries. These libraries would then need to be extended. Reduction of $`N1`$-point function tensor integrals are much more compact and easier than for $`N`$-point functions. This said for the case at hand, and for that matter any relic density calculation where the LSP is a neutralino, these highest rank tensors are not needed. It is easy to show that the highest rank tensors for $`22`$ only occurs when all the external particles are bosons. In our case, for the box, one has $`M_{\mathrm{max}}=3`$. For $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0f\overline{f}`$, $`M_{\mathrm{max}}=2`$.
The choice of the momenta circulating in the loop, $`s_1,s_2,s_3`$, is not unique and depends on the particular graph. If $`DetG(s_1,s_2,s_3)=0`$ we first form all three sub-determinants $`DetG(s_i,s_j)`$ and take the couple $`s_i,s_j`$ that corresponds to $`Max|Det(s_i,s_j)|`$. Then the third $`s_k`$ is distributed in the basis $`s_i,s_j`$ and the corresponding $`\alpha `$ and $`\beta `$ are read. In fact, suppose $`DetG(s_1,s_2)0`$, then it is revealing to always write
$`s_3`$ $`=`$ $`\alpha s_1+\beta s_2+\epsilon _T\mathrm{with}`$
$`s_1.\epsilon _T`$ $`=`$ $`s_2.\epsilon _T=0\mathrm{meaning}`$
$`\epsilon _{T,\mu }`$ $`=`$ $`ϵ_{\mu \alpha \beta \delta }s_1^\alpha s_2^\beta t^\delta .`$ (A. 6)
$`\epsilon _T`$ is a vector that is orthogonal to both $`s_1`$ and $`s_2`$ that is easily reconstructed once $`\alpha `$ and $`\beta `$ are
$`\alpha ={\displaystyle \frac{s_2^2s_3.s_1s_1.s_2s_2.s_3}{DetG(s_1,s_2)}},\beta =\alpha (s_1s_2).`$ (A. 7)
This construction makes it clear that
$`DetG(s_1,s_2,s_3)=\epsilon _T^2DetG(s_1,s_2).`$ (A. 8)
This shows, in a most transparent manner, that the determinant vanishes when $`\epsilon _T^2=0`$. This can occur when the components of this vector vanish, $`\epsilon _T^\mu 0`$, and therefore $`s_3`$ is not an independent vector as the case of this paper for $`v=0`$. It also occurs, a point which is often overlooked, when $`\epsilon _T`$ is light-like. However the orthogonality constraint means that the other vectors $`s_1,s_2`$ are space-like. Therefore this case does not occur for real particles that are time-like, and hence does not occur for our $`22`$ process. It will be shown, in a separate publication, that when $`\epsilon _T`$ is light-like a segmentation is still possible. The algorithm can also be improved by expanding around $`DetG(s_1,s_2,s_3)=0`$. We will come back to the details of this issue in a future publication. Note that for $`v=0.5`$ and for all three processes studied in this paper the standard reduction formalism, without segmentation, was used. |
warning/0507/cond-mat0507349.html | ar5iv | text | # Self-learning Kinetic Monte-Carlo method: application to Cu(111)
## I Introduction
The past decade has witnessed a surge in research activities which aim at bridging the gap in length and time scales at which a range of interesting phenomena take place. Some examples of such activities pertain to studies of epitaxial growth, and nanostructuring of materials. The aim in such work is to utilize information obtained at the microscopic level to predict behavior at macroscopic scales. There are thus several key tasks to be undertaken, each of which is a challenge in itself. The first of these is an accurate determination of the energetics and dynamics of the system at the microscopic level. For selected systems this may be achieved through ab initio electronic structure calculations yu96 which are becoming increasingly feasible for complex systems, even though they remain computationally intensive. A reasonable alternative, albeit not as reliable or accurate, is the application of one of several genres of many body interatomic potentials erc88 . With these interatomic potentials it has been possible to carry out computational and theoretical studies of a range of surface phenomena using techniques like molecular statics and molecular dynamics. Molecular dynamics simulations in particular are capable of revealing the essential details of microscopic phenomena as they unfold as a function of temperature, pressure and other global variables but the application is limited in time and length scales. Since most thermally activated atomistic processes occur in the range of picoseconds, they are best captured with time steps in femtoseconds which limits total simulation time to few microseconds. These times are many orders of magnitude smaller than processes happening in the laboratory. For example, epitaxial growth and surface morphological changes take place in minutes and hours and are controlled by atomic processes which are infrequent compared to atomic vibrational times of picoseconds. The challenge in molecular dynamics simulations is to find reliable ways which capture infrequent processes and extend to longer time scales with reasonable computational resources.
An alternative to molecular dynamics simulations for examining surface phenomena is offered by the kinetic Monte-Carlo (KMC) technique in which the rates of various eligible atomic processes are provided as input bor75 ; gil76 ; vot86 . If this input is accurate and complete, KMC simulations are in good position to mimic experiments. Since the task of accumulating a complete set of atomic processes is non-trivial, standard KMC simulations are typically performed with a set of most obvious simple atom or concerted processes as input, and all others either ignored or included in approximate ways (e-g bond counting models) or added in an ad hoc manner to fit experimental data. \[NEW\]With a reduced set of barriers, activation energies become effective values rather than actual values which may be compared with those obtained from experimental data but may not reveal the intervening microscopic processes. This is obviously problematic. Furthermore, it has been shown that novel multiple atom processes may play important role in providing mass transport on surfaces such as Cu(100) rowshear ; clustdif , and Ir(111) diff\_exp2 ; diff\_exp3 . Any realistic simulation should have a provision for uncovering such processes and including their energetics in evolution of the system.
To overcome these limitations of the two most common approaches for simulating temperature dependent morphological evolutions of surfaces and interfaces, several accelerated schemes have been presented in recent times vot00 ; hen01 ; hen03 . In a set of studies, Voter et al, vot02 have concentrated on enhancing the time scales achievable in MD simulations through three different strategies: parallel-replica, temperature-accelerated dynamics and hyperdynamics. Fichthorn and co-workers, in related work, apply the bond boost method mir03 to extend the time scales in their simulations. The basis principle in these methods is to make the system evolve faster, sampling a larger phase space, either through smartly connected parallel processors, or application of a boost so that the system can overcome energy barriers with relative ease, or by raising the temperature of the system. At the very least, novel and infrequent processes may be revealed through such accelerated schemes. The main issue is the assurance of one-to-one correspondence between the temporal evolution of the accelerated and non-accelerated systems and whether the approach actually leads to a large speed-up for a particular system of interest. The reader is referred to the original papers for further details and suitability of the techniques, to specific cases.
Another promising scheme has focussed on the completeness issue of KMC by allowing the system to evolve according to single and multiple atom processes of its choice. The key to the method is the generation of saddle points in the potential energy surface and benefits from the advances that Jonsson and co-workers neb ; dimer have made in procedures for extracting diffusion paths and energy barriers using efficient search procedures. Once a large (sufficient) number of saddle points have been identified, the expectation is that the system will evolve naturally according to its inherent mechanisms. The method we propose here is in principle related to the latter approach, with a very important difference. We employ a pattern recognition scheme which allows efficient storage and subsequent retrieval of information from a database of diffusion processes, their paths and their activation energy barriers. The procedure presented here is thus efficient and reliable. The removal of redundancies and repetitions in the calculations of energetics of system dynamics speeds up the simulations by several orders of magnitude making it feasible for a range of applications. Since the generation of the database and its future usage through recognition patterns is akin to the simulation procedure learning from itself, we call the technique proposed here self learning KMC (SLKMC). While the proposed technique can be applied to any surface systems, our interest is in the examination of atomistic phenomena as related to growth on fcc(111) surfaces. This is a challenging surface since the lack of surface corrugation makes the energy landscape relatively flat with a number of diffusion processes which are equally competitive. Some such atomistic processes may include those with multiatoms which are typically ignored in standard KMC techniques. In this paper we focus our attention on some characteristics of the proposed technique and its application to homo-epitaxy on fcc(111) surfaces through considerations to of the diffusion and coalescence of two dimensional Cu adatom islands on Cu(111). The structure of the paper is as follows. In the next section we present some essentials of the self-learning KMC framework. This is followed in section III with results of the applications of the method to examine morphological evolution of two dimensional Cu islands on Cu(111). Section IV contains our conclusions.
## II Essentials of Self-Learning Kinetic Monte Carlo Method
Although the principle of the proposed technique is generally applicable, we need a specific surface geometry to illustrate its details. For reasons mentioned above our interest is homoepitaxy on fcc(111) surfaces. We provide in this section some details of the model system, together with an outline of the standard kinetic Monte Carlo method for completeness. This is followed by a summary of the pattern recognition and labeling scheme that we invoke to obtain a self-learning KMC methodology.
### II.1 MODEL SYSTEM
To mimic fcc(111) surface we consider a 2-layer substrate, with periodic boundary conditions in the XY plane (which is parallel to the surface), to uniquely identify the fcc and hcp hollow sites on the surface. The system of interest (such as an adatom island, vacancy island, or any other nanostructure whose morphological evolution or diffusion is to be determined) is placed on top of the substrate. In this initial study only occupancy of fcc sites (i.e. hollow sites with no atom in the layer below) on the substrate is allowed. \[NEW\]While there is experimental justification for assuming fcc-site occupancy for Cu adatoms on Cu(111)gie03 , we are aware that on Ir(111) atoms may also occupy hcp sites (hollow sites with an atom in the layer below) wan90 . Infact, even for homoepitaxial growth on Cu(111) under certain other experimental conditions hcp-site occupancy has been reported cam00 . Furthermore, adatoms, dimers, and other smaller clusters may use the hcp site as an intermediate repp2003 one during its motion. The method we are proposing can easily be generalized to include hcp occupancy. We are also assuming that the diffusion is via hopping and that during the simulation the substrate atoms are kept fixed. This restriction can be removed in future work. For the moment our interest is in the in-plane (2D) motion of adatoms, vacancies and their clusters on Cu(111).
### II.2 SOME INGREDIENTS OF KINETIC MONTE CARLO
The goal of kinetic Monte Carlo (KMC) is to mimic real experiments through sophisticated simulations. For these simulations to be realistic, one has to implement increasingly complex scenarios requiring intensive use of state-of-the art software and hardware. At the heart of a KMC simulation of the time evolution of a given system lie the mechanisms that are responsible for determining the microscopic processes to be performed at any given time. To illustrate the point, consider a system containing N particles at a given time with $`N_e`$ possible types of processes. Let us also associate with each process-type (i), $`n_i`$ the number of particles in the system that are candidates for this process-type, the activation energy barrier $`\mathrm{\Delta }E_i`$ and a pre-factor $`\nu _i`$. The microscopic rate associated with process (i), within Transition State Theory (TST) tst , is then,
$$r_i=\nu _iexp\mathrm{\Delta }E_i/kT,$$
(1)
where k is the Boltzman constant, and T the surface temperature. The total rate R of the system is further given by:
$$R=\underset{i=1}{\overset{N_e}{}}R_i,$$
(2)
where $`R_i=n_ir_i`$ is the macroscopic rate associated with process-type i.
In KMC simulations, the acceptance of a chosen process is always set to one. However, the choice of a given process is dictated by the rates. First, a process-type is chosen according to its probability $`p_i=R_i/R`$, and then a particle is randomly chosen from the set $`n_i`$ to perform this process.
\[NEW\]The essential elements of the KMC method are thus the processes ’i’ and their activation energy barriers $`\mathrm{\Delta }E_i`$ whose determination requires that availability of reliable interatomic interaction which may be obtained from first principles or from model potentials. In this paper, all activation energies are determined using the embedded atom method (EAM). This is a semi-empirical, many-body interaction potential eam8693 . Although the EAM potentials neglect the large gradient in the charge densities near the surface and use atomic charge density for solids, for the six fcc metals Ag, Au, Cu, Ni, Pd, and Pt, and their alloys, it has done a successful job of reproducing many of the characteristics of the bulk and the surface systemseam8693 .
To get back to the issue of the determination of diffusion processes, their paths and their activation energy barriers, we should note that several interesting and appealing approaches have been proposed in the past few years. These methods include the nudged elastic band (NEB) method neb , the step and slide mothod mir01 eigenvector following eigen , and temperature accelerated MD tempac . Each of these methods has its own computational demand and measure of accuracy whose balance dictates the choice of the approach. For the studies presented in this paper, we find the simple ’drag’ method to be adequate, as we shall see. This is, of course, a rudimentary method in which the moving entity is dragged in very small steps towards probable (aimed) final state. The dragged atoms is constrained in the direction towards the aimed position while the other two degrees of freedom (perpendicular to this direction) and all degrees of freedom of the rest of the atoms in the system are allowed to relax. The other atoms are thus free to participate in the move, thereby activating many-particle processes (when neighbor adatoms start to follow central leading atom). In connection with SLKMC, the central atom is always dragged towards one of its vacant fcc site. A more general way to map out the potential energy surface is to use the grid method which has been successful in finding non-trivial diffusion paths and saddle points karim01 .
### II.3 SELF LEARNING KINETIC MONTE CARLO METHOD
As we have already mentioned, the limitation of standard KMC is its reliance on an ad hoc choice of processes and hence lack of completeness. For these reasons and also because of experimental observations of complex and unforeseen processes, the predictive power of KMC is in question. A rethinking of the way we perform KMC has become a necessity. Simulations with an a priori chosen catalogue of processes need to be replaced by a continuous identification of possible processes as the environment changes. For these innovations in the KMC procedure, the local environment is the key issue and its complexities need to be exploited. With this in mind we are proposing a methodology in which the base ingredient is the collection of local environments of undercoordinated atoms found automatically during the simulations and labeled and stored for subsequent usage in the simulation. As a concrete example of our approach we have choosen the fcc(111) surface with a six-fold symmetry. For simplicity, we assume that any process in this system will involve a central (undercoordinated) atom and atoms in the next 3-shells as illustrated in Fig. 1. The motif in Fig. 1 is to serve as a ’cookie cutter’ and is placed on all active atoms in the system to define their local environment. We further assume, without loss of generality that any process may be described in terms of the central atom moving to a neighboring vacancy accompanied by the motion of any other atom or atoms in the 3 surrounding shells. The labeling of the surrounding atoms is done in binary and a base ten number is then associated with the first shell configuration. The same procedure is followed for atoms in the second and third shell. Hence, for an atom in the system to be active (i.e. central atom for a given process), it should have a vacancy in its first shell (or an occupancy number less than 63 for the cookie cutter); as illustrated in Fig.1.b.
Once the atoms are classified as active and non-active and encrypted within the 3-shell scheme, we proceed by determining all possible processes associated with every active atom. Next the determination of the activation energy and pre-factor is performed for all processes. Examples of how processes are labeled and stored in the database are given in Fig.2. In this figure, full circles represent occupied sites and open circles vacancy sites. Fig. 2a illustrates the ”diffusion along a step” process where the central atom labeled 1 moves to the vacant site 2 along the step formed by atoms numbered 30, 15, 6, 7, 19 and 37 in the cookie cutter. The initial configuration for this process is recorded in base 10 as (48,3968,261120) in the database and shown with the base 2 label in the figure. The move in Fig. 2a is recorded as atoms 1 going to position 2 (1,2) and the activation energy barrier for the process in Fig. 2a is found to be 0.31 eV. Similarly, for the multi-atom process illustrated in Fig. 2b, the initial configuration in base 10 is recorded along with the sequence of motion of atoms involved in the process which in this case is 1 going to 4, 6 to 1 and 15 to 15, which is recorded as (1,4;6,1;15,5). This multi-atom process was found during the coalescence of two islands and will be discussed later. Its activation energy barrier of 0.595 eV is also recorded with the label.
The bottleneck for the simulation is the determination of the activation energy and the prefactor for all possible processes. Even when we make the widely-used assumption that all the processes have the same pre-factor, the calculation of the activation energy is very expensive if one needs accurate values. Note that since the activation energy is in the exponential, any small variation in the activation energy results in a substantial change in the relative probabilities and hence the outcome of the whole simulation. In standard KMC these energy barriers are provided as input. If, however, as we and others hen01 are proposing that these barriers be calculated on-the-fly, the process will be sped-up if provisions are made to avoid recalculations. In the method proposed here this is achieved through the storage of activation energy barriers tagged to specific atomic processes in the database. This is the basis of our KMC in which ”self-learning” is achieved by the system through the ability to: (1) calculate activation energies on the fly; (2) store them in a database; and (3) recognize and retrieve them using the labeling described above. Step 1 is not new. It was already proposed by Jonsson et al hen01 and Voter and Montalenti vot02 . Step 2 and 3 are unique to our approach and help remove redundancies in the calculations. At any given time, after all the processes have been sorted out, a search for the activation energies in the database is launched. If a new process is encountered, the actual calculation is performed and this process with its activation energy is added to the database. Once the processes are classified and macroscopic rates are calculated, we proceed to perform one Monte Carlo step in which a randomly selected process is executed. The entire simulation process is summarized in the flowchart (Fig. 3). At later times in the simulation, when the system encounters environments for which some of the possible processes have been met earlier, a retrieval process of the activation energy from the database substitutes the actual calculation. This gives a tremendous gain in the execution times as evident in our application to the diffusion of 2D Cu clusters on Cu(111). With modest computational resources, it was possible to carry out the simulation for a number of MC steps, large enough to provide good statistics. The exact number of steps may vary from problem to problem.
In the next section, we discuss some key features of the database collected during an extended simulation along with the results obtained from applying the SLKMC to post-deposition analysis of homoepitaxy on Cu(111).
## III Application of SLKMC to Morphological Evolution of 2D islands on Cu(111)
Since the devil is generally in the details, we present below results of the application of SLKMC to study Cu cluster diffusion and coalescence on Cu(111). After giving some specifics of the model system, we present an analysis of the database which includes an evaluation of the accuracy of the calculated energy barriers and other factors affecting the simulation (cpu) time. We also comment of the presence and importance of multiatom processes. This is followed by the results and discussion of the diffusion and coalescence of 2D clusters on Cu(111).
### III.1 Model Systems
\[NEW\] In the first example, i.e. the study of the diffusion of 2D adatom islands of Cu(111) we have chosen four specific sizes (19, 26, 38 and 100 Cu atoms) for which we already have results for comparison with a KMC simulation using a fixed data base of logical processes involving single atom periphery diffusion mrs05. For the second application to the process of cluster coalescence, our model system consists of 2 adatom islands, one consisting of 78 atoms and the other 498 atoms placed on top of the 2-layer substrate.
### III.2 Examination of the collected database
\[NEW\] Since it is important to check the reliability of the data in the created database, we have compared in Table 1 the energy barriers that we obtained for some typical diffusion processes presented in Fig. 4, using both the drag and the NEB methods. We also include in the Table values available in the literature. The comparison in the Table attests to the reliability of the drag method as compared to the more time-consuming NEB procedure. For example, with the drag method we were able to achieve speed-up of atleast an order of magnitude in the cpu time for the calculation of the energy barriers, as compared to one in which we applied the spherical repulsion methode tru04 to obtain the final states for a given initial state followed by application of NEB method for the calculation of the activation energies.
As an illustration of the richness of the database that we collect, we plot in Fig. 5 the energy distribution of about 5000 diffusion processes which have been accumulated during a simulation containing several hundreds of millions of Monte Carlo steps. Note from Fig. 5 that the distribution is very wide covering activation energies as small as few tens of a meV to about 1 eV. \[NEW\] Unlike studies in which activation energy barriers are calculated for perfect and periodic systems kar95 , energy barriers cannot be classified into groups. Note that in the calculations of the energy barriers differences are introduced when the effect of next nearest meighbors of the local environment is included in the calculation, as we have done. Note also that the accumulation of the database does not proceed uniformly with time, as reflected in in the insert of Fig. 6. The SLKMC starts, in this case, by accumulating about 400 different processes during the very first MC step, after which the database is ”quasi-saturated” for a certain period of CPU time. This is followed by an other phase of accumulation of about 600 processes, and so on. It is clear from the slope in Fig. 6 that when the simulation runs with a ”quasi-saturated” database, the number of KMC steps/CPU time increases dramatically. During a heavy build up of the database, the yield is about 80 KMC steps per second and can go up to several thousands of of KMC steps per seconds as the database saturates. \[NEW\]The onset new events in the database after certain duration of simulation does raise the issue of measures that would assure that the data base is complete. So far we have found the database to saturate after runs of about 100 - 500 million MC steps. Actually, for the systems under study we have rarely found new processes to set in after 10 million time steps \[OLEG, ALTAF, PLEASE CHECK\].
One of the most important features of the method, as we have seen, is its ability to treat many-particle processes, the so called ”concerted atomic motion”. The recent version of the code allows inclusion of simultaneous displacements for atoms up to the third shell. from our simulations of several types of local environments (straight steps with kinks, compact islands, fractal like islands) we found that in some cases many particle processes play important role in providing atomic transport mrs . They are especially important in the case of low coordination systems, like fractal islands. In such cases atoms are weakly bound and prefer to perform concerted motion rather than single atomic jumps. Furthermore, their importance increases with decreasing size of the cluster. In fact molecular dynamics simulations simulations of a 10-atom Cu island on Cu(111) at 700K and 900K, show that the island moves by concerted displacement rather than through single atoms motion ahlam05 . We next move move onto examination of the results for two specific applications of SLKMC.
### III.3 Morphological Evolution
#### III.3.1 Diffusion of 2D islands
As a first application of the SLKMC method, we present results for the diffusion of 2D Cu islands on Cu(111) of four sizes: 19, 26, 38 and 100 atoms. These simulations were performed using 10 million \[IN THE MRS PAPER WE SAY 500 MILLION, WHICH IS IT\] MC-steps at 300K and 500K. During the simulation, the position of the center of mass was recorded at each MC step along with the performed process. After 10 million MC steps, the islands have moved far enough that their diffusion coefficient may be extracted from the mean square displacement of the center of mass. In Fig. 7 we show the trace of the position of the center of mass on the (x,y) plane for both 19 and 38 atoms clusters at 300K. Note the dark spots for both cases indicating a stick-slip type motion of the center of mass. The corresponding mean square displacement, for these two atoms, as a function of time show a linear behavior (within statistical errors) and is shown in Fig.8. The extracted slope from the mean square displacement plot gives the coefficient diffusion. In Table 2, we report the diffusion coefficient for the four cluster sizes at 300K and 500K. Note that the diffusion coefficient increases exponentially with temperature. The decrease of the diffusion coefficient with the cluster size follows a power law ($`D=N^{1.57}`$ at 300K and $`D=N^{1.64}`$ at 500K), which is in good agreement with previous results dif\_th3 . The virtue of our calculation is that the atomic processes leading to cluster diffusion were picked by the system itself during the simulation. The frequencies of the contributing processes vary with cluster size and, more importantly, with surface temperature. \[NEW\] Although these results have already appeared in a conference proceeding mrs , we present in Table 3 the frequencies of the various contributing atomistic processes. Included in the table are also the frequencies of events that were found in a parallel simulation carried out using the standard KMC simulation with a predetermined catalogue of 294 (49x6) processes involving single-atom peripheral diffusion /citegho03. It is interesting to note that for the later type of simulation the 19, 38, and 100 atom cluster hardly moved at 300 K, while the diffusion coefficient for the 26 atom cluster was found to be 0.48 $`\AA ^2`$/sec, to be compared with 0.17$`\AA ^2`$/sec. At 500 K, the diffusion coefficients (in units of $`\AA ^2`$/sec) for the 19, 26, 38, and 100 atom clusters are found by this standard KMC simulation to be 1.29x$`10^4`$, 2.24x$`10^4`$, 2.22x$`10^4`$, and 0.32x$`10^4`$, respectively. For 19 atom cluster the diffusion coefficients by KMC and SLKMC differ by an order of magnitude: single atom periphery diffusion underestimates the mobility of this perfect hexagon. For the other sizes the differences are also noticeable. Some insights for these differences can be obtained from Table 3 in which we have summarized the frequencies of the processes that are executed in the two types of simulations. While the most dominant mechanism is that of a single adatom diffusion along the (100)-microfacetted step edge (A) in both SLKMC and standard KMC, there are differences in the frequencies which the various processes are executed in the two types of simulations, and also in the appearance of several new processes like detachment from the corners engulfed by the A and B-type ((111)-microfacetted) step edges mrs . In Table 3 the activation energy barriers obtained from the drag method for each of the processes are also compared with those from the nudged elastic band (NEB) method. Further details of the processes listed in Table 3, including their nomenclature, local configuration, etc. can be found at (http://www.phys.ksu.edu/ rahman). Note that even though multiple atom processes are not very frequent, they do occur, and in some cases they may be the rate limiting step for the diffusion of cluster of a particular size.
#### III.3.2 Island coalescence
As a second example of application of the SLKMC method, we present here results of simulation of the coalescence process in which two adatom islands join together to form a larger island with an equilibrium shape on Cu(111). This simulation was performed at 300K using a small island containing 78 atoms and with an arbitrary shape, put close to a larger island containing 498 atoms with a circular shape. Successive snapshots of the system during the SLKMC simulation are shown in Fig.9, for a total number of forty millions KMC steps. From this figure, one notes that a neck between the two islands forms during the first hundred thousand KMC steps, corresponding to a physical time of 0.25s. After this time, the neck grows until the two islands form an elongated single island after about 10 seconds. Finally, The shape of the island evolves to a quasi-triangle with mostly (111)-steps (A-type), which is a result of the assymetry in the activation energy barriers associated with A and B-type steps (see Table 1). In order to get an insight in the mechanisms involved in the neck formation, we have analyzed the frequency distribution of key processes during the first and second hundred thousand KMC steps. \[NEW\]Three types of processes appear prominent in the coalescence of these two clusters: kink detachment on an A-type step (2a in Fig.4), the reverse of 2a (Rev. 2a in Table 2) also called kink incorporation, and diffusion along A-type step (4a in Figure 4). Listed in Table 4 are the frequencies for these processes. We note from Table 4 that during the formation of the neck, kink detachment and kink incorporation count for about 15% of the performed processes, another 70% involve diffusion along A-type step and other single and multiple atom processes including kink-rounding and two atoms diffusion along steps constitutes the remaining 15%. For the second hundred thousand KMC steps, the simulation is mostly dominated by diffusion along the A-type step (about 96%), with about 4% counting for various mechanisms. The important fact to note here is that kink-detachment and kink-incorporation contributions drop to almost zero after the neck has been formed. Detailed analysis of similar simulations involving islands of various sizes and shapes are actually in the processes of being performed and will be published elsewhere. A similar process for our simulations of cluster island coalescence are in qualitative agreement with the observations made by M. Giesen et al giesen using scanning tunneling microscopy.
## IV Conclusion
We have addressed the issue of completeness of KMC simulations by proposing a method in which the system finds, calculates, and collects the energetics of all possible diffusion processes that the moving entities are capable of performing. What separates our technique from others recently proposed is the provision for storing and retrieving the environment dependent activation energy barriers from a data base. Examination of the data base shows that the simulation proceeds much faster when the set of processes is ”quasi-saturated’ and that after sampling such regions the system has the ability to trigger the participation of new diffusion processes requiring enhanced CPU’s for the calculation of new activation energy barriers. The system eventually settles down, the number of MC steps needed to do this depends on the system and the number of entries already in the data base (about $`10^710^8`$ steps). With the use of the pattern recognition scheme we are able to identify and calculate the frequency of occurance of individual single and multiple atom diffusion processes that actually participate in the evolution of a particular entity. The microscopic details of the processes involved in surface morphological evolution can thus be documented for a system that have the freedom to evolve on their own accord. We show this through applicationto the diffusion and coalescence of 2D adatom islands on Cu(111) for which the simulation began with an empty data base. Once a substantial accumulation has occured, the simulation time speeds up by orders of magnitudes and allows the calculation of system dynamics for time scales relevant to those phenomena happening in the laboratory. Interestingly, the two simple examples that we have presented here show that only a few dozen diffusion processes are in the end vital for a diffusion event. The question of course is: which ones? Our approach answers this question. As we have already alluded to, the task of calculating diffusion prefactors is still ahead of us. This is particularly important since we find many competing processes to differ only slightly in energy and differences in their vibrational entropy contributions to the prefactors can make a difference in the ultimate evolution of the film morphology. Another important result from our simulations with the open data base is that dynamical evolution of the system with prejudged diffusion processes may yield erroneous results. Also, the pattern recognition schemes to be a prudent way to develop data base of diffusion processes and their energetics. It does involve a lot of work in the beginning but once the data base is compiled, it can be used for any type of simulation of the system. Of course, for realistic simulations of thin films we need to incorporate exchange and other processes which involve motion in 3D. \[NEW\] We have already alluded to the importance of inclusion of hops and occupancies of the hcp site. Such efforts are currently underway. Infact, preliminary results have already been obtained for the diffusion of small clusters (2-10 atoms) in which the SLKMC code with the both fcc and hcp occupancy. However we leave this for presentation elsewhere. We have also confined ourselves so far to homoepitaxy. Our future work will also include application to heteroepitaxy for which only minor changes need to be introduced in the methodology.
## V Aknowledgements
We thank James Evans, Chandana Ghosh and Ahlam Alrawi for helpful discussions. This work was supported by NSF-CRDF RU-P1-2600-YA-04, NSF-ERC 0085604 and NSF-ITR 0428826.
Figure Captions
Fig. 1: a) The three-shell indexing around the central atom labelled 1; b) Signature of a particular 2D cluster configuration in base 2 and base 10.
Fig. 2: Sample (a) single atom and (b) multiple atom processes involved in the diffusion of 2D clusters presented with their specific labels for our database.
Fig. 3: Flowchart for SLKMC simulation
Fig. 4: Selected single atom processes on the two types of steps, A (100-microfacetted) and B (111-microfacetted), on fcc(111) surface. Process 1 is kink-detachement-rounding, 2 is kink-detachement-along step, 3 is adatom detachment from step, and 4 is adatom diffusion along step. The labels ’a’ and ’b’ refer to step A and B, respectively.
Fig. 5: Distribution (percentage) of activation energies of stored processes in the database during a SLKMC simulation.
Fig. 6: Variations in the number of KMC steps per CPU time ( i.e. performance) and the build-up of the databse as a function of the number of KMC steps (insert).
Fig. 7: Trace of center of mass of 19 and 38 atoms Cu clusters on Cu(111) at 300K as obtained from SLKMC simulations ($`10^7`$ steps).
Fig. 8: Mean square displacement (MSD) for 19 and 38 atoms Cu clusters on Cu(111), as function of time at 300K.
Fig. 9: Coalescence of a small Cu cluster (78 atoms) with a larger one (498 atoms) on Cu(111) at 300K, using SLKMC ($`10^7`$ steps).
Table 1: Diffusion energy barriers for selected mechanisms as shown in Fig. 4.
| Process | Drag Method (eV) | NEB Method (eV) | Ref. kar95 (eV) |
| --- | --- | --- | --- |
| 1a | 0.68 | 0.66 | - |
| 2a | 0.53 | 0.52 | - |
| 3a | - | 0.65 | - |
| 4a | 0.25 | 0.25 | - |
| 1b | 0.60 | 0.59 | 0.59 |
| 2b | 0.58 | 0.56 | 0.54 |
| 3b | 0.68 | 0.67 | 0.67 |
| 4b | 0.32 | 0.30 | 0.29 |
Table 2: Diffusion coefficient for 2D Cu islands on Cu(111) ($`\AA ^2`$/sec)
| cluster Size | 300K | 500K |
| --- | --- | --- |
| 19 | 0.196 | 1.67x$`10^5`$ |
| 26 | 0.170 | 8.05x$`10^4`$ |
| 19 | 0.117 | 4.27x$`10^4`$ |
| 19 | 0.016 | 1.02x$`10^4`$ |
Table 4: Frequency of selected processes during the coalescence of 2 islands
| Process | Barrier (eV) | Frequency ($`01x10^5`$ steps) | Frequency ($`12x10^5`$ steps) |
| --- | --- | --- | --- |
| 2a | 0.530 | 7.41 | 0.03 |
| Rev. 2a | 0.220 | 8.43 | 0.04 |
| 4a | 0.25 | 69.66 | 95. 88 |
| others | - | 14.50 | 4.05 | |
warning/0507/cond-mat0507465.html | ar5iv | text | # Long-Range Coulomb Effect on the Antiferromagnetism in Electron-doped Cuprates
\[
## Abstract
Using mean-field theory, we illustrate the long-range Coulomb effect on the antiferromagnetism in the electron-doped cuprates. Because of the Coulomb exchange effect, the magnitude of the effective next nearest neighbor hopping parameter increases appreciably with increasing the electron doping concentration, raising the frustration to the antiferromagnetic ordering. The Fermi surface evolution in the electron-doped cuprate Nd<sub>2-x</sub>Ce<sub>x</sub>CuO<sub>4</sub> and the doping dependence of the onset temperature of the antiferromagnetic pseudogap can be reasonably explained by the present consideration.
\]
The Fermi surface (FS) evolution with electron doping in Nd<sub>2-x</sub>Ce<sub>x</sub>CuO<sub>4</sub> has been observed by angle-resolved photoemission spectroscopy (ARPES) experiments. Recently, much attention has been paid to understanding the physics of this phenomenon. At low temperatures, the electron-doped cuprates are in the states of antiferromagnetic (AF) phase within a wide doping region. The FS evolution is therefore closely relevant to the antiferromagnetic correlation. From the framework of the Hubbard model, the doped electrons occupy the upper Hubbard band. On the other hand, according to the experimental observation, with increasing electron doping, the energy gap should decrease and eventually close up at the optimal doping, $`x0.14`$, where the AF phase terminates. Therefore, to interpret the experimental results, one needs to assume the on-site interaction $`U`$ in the Hubbard model to be dramatically decreasing with increasing electron doping concentration. This is a puzzle within the Hubbard model with constant on-site interaction. Some investigators have treated the doping dependence of $`U`$ by considering some kind of screening.
In this work, we explore the long-range Coulomb effect (LRCE) on the AF ordering. Due to the Coulomb exchange effect, the LRCE results in excess electron hopping. As a whole, the excess hopping tends to frustrate the AF ordering. With increasing the electron doping, this exchange effect becomes significant, leading to the decreasing of the energy gap. We will see this gives a reasonable explanation of the FS evolution and the envelope of the onset temperature of the antiferromagnetic pseudogap in electron-doped Nd<sub>2-x</sub>Ce<sub>x</sub>CuO<sub>4</sub>.
We start with the following two-dimensional square lattice model
$$H=\underset{ij,\sigma }{}t_{ij}c_{i\sigma }^{}c_{j\sigma }+U\underset{i}{}n_in_i+\frac{1}{2}\underset{ij}{}v_{ij}n_in_j$$
(1)
where $`t_{ij}`$ denotes the hopping energy of an electron between the lattice sites $`i`$ and $`j`$, $`c_{i\sigma }^{}`$ ($`c_{i\sigma }`$) represents the electron creation (annihilation) operator of spin-$`\sigma (=\pm 1`$ for up and down spins, respectively) at site $`i`$, $`n_{i\sigma }=c_{i\sigma }^{}c_{i\sigma }`$, and $`n_i=n_i+n_i`$ is the electron density operator. Besides the on-site Hubbard repulsion $`U`$, we take into account the long-range Coulomb interaction in the third term. Such a similar type of Hamiltonian has been used to investigate the LRCE on the $`d`$-wave pairing for the hole-doped case.
In the hopping term, besides the nearest neighbor (n.n) hopping, the other hopping processes within a range need to be included. The essential role of the next n.n hopping for the validity of the single-band Hubbard model for describing the cuprates has been investigated by comparing it with the two-band and three-band Hubbard models. The particle-hole asymmetry in the cuprates can be understood by taking into account the next n.n hopping . By including the next and third n.n hopping parameters in the types of $`tJ`$ and Hubbard models, numerous studies have been carried out to explore the properties of the electron-doped cuprates. We here take two more additional parameters for the fourth and fifth hopping, each of which is smaller than the formers. In terms of the n.n hopping parameter, $`t_{(1,0)}t`$, the values of other 4 parameters are given as $`t_{(1,1)}=0.325t`$, $`t_{(2,0)}=0.17t`$, $`t_{(2,1)}=0.121t`$, and $`t_{(2,2)}=0.07t`$, where the subscripts instead of $`ij`$ denote the components of a vector of length $`|𝐢𝐣|`$. The values of the next and third n.n hopping parameters are approximately the same as that in the literatures. From the hopping term, the single-particle dispersion $`ϵ_k^0`$ is given by
$$\begin{array}{cc}\hfill ϵ_k^0=& 2t(c_x+c_y)4t_{(1,1)}c_xc_y2t_{(2,0)}(c_{2x}+c_{2y})\hfill \\ & \\ & 4t_{(2,1)}(c_{2x}c_y+c_{2y}c_x)4t_{(2,2)}c_{2x}c_{2y}\hfill \end{array}$$
with $`c_{lx}=\mathrm{cos}(lk_x)`$.
For the long-range Coulomb interaction $`v_{ij}`$, we take the following form
$$v_{ij}=V_1\mathrm{exp}(|𝐢𝐣|/d)/|𝐢𝐣|,$$
(2)
where the length scale is in unit of the lattice constant. Here we set $`V_1=1.0t`$ and $`d`$ = 4 for the present calculation. For the on-site $`U`$, we use $`U=4.8t`$, which is within the range of the doping-dependent interaction adopted in the existing calculations. The strength of $`V_1`$ with $`V_1/U0.21`$ is a typical one.
By the mean-field (MF) approximation, from the long-range Coulomb term, we get the exchange part of the self-energy as
$$\mathrm{\Sigma }_k^x=\frac{1}{N}\underset{k^{}}{}v(|𝐤𝐤^{}|)c_k^{}^{}c_k^{}$$
(3)
where $`N`$ is the number of lattice sites, $`v(|𝐤𝐤^{}|)`$ is the Fourier transform of $`v_{ij}`$, and $`\mathrm{}`$ denotes an average. This exchange self-energy is equivalent to a hopping energy in real space. The corresponding hopping integral is given as $`t_{ij}^x=v_{ij}c_i^{}c_j`$. The significance of this exchange effect and its doping dependence will be discussed later. With the contribution from the exchange self-energy, the effective single-particle dispersion is given by $`ϵ_k=ϵ_k^0+\mathrm{\Sigma }_k^x`$.
At low temperatures, the strong on-site Coulomb interaction leads to the AF ordering. In the case of the electron doping under consideration, the AF ordering is a commensurate spin-density wave. The order parameter is given by $`\mathrm{\Delta }=Un_{}(Q)n_{}(Q)/2N`$, where $`n_\sigma (Q)`$ is the Fourier transform of the electron density of spin-$`\sigma `$ at wave vector $`Q=(\pi ,\pi )`$. By the MF approximation, the energy spectrum of the quasiparticle is then given by
$$E_k^\pm =(ϵ_k+ϵ_{k+Q}\pm E_k)/2,$$
(4)
where the + (-) sign refers to the upper (lower) Hubbard band, and $`E_k=\sqrt{(ϵ_kϵ_{k+Q})^2+4\mathrm{\Delta }^2}`$. The parameter $`\mathrm{\Delta }`$ and the chemical potential $`\mu `$ are self-consistently determined by
$`{\displaystyle \frac{U}{N}}{\displaystyle \underset{k}{}}[f(E_k^{})f(E_k^+)]/E(k)=1,`$ (5)
(6)
$`{\displaystyle \frac{1}{N}}{\displaystyle \underset{k}{}}[f(E_k^{})+f(E_k^+)]=n,`$ (7)
where $`f`$ is the Fermi distribution function, and $`n`$ is the electron density. In terms of doping concentration $`x`$, $`n=1+x`$.
To analyze the ARPES observations on the FS evolution, one usually considers the spectral density occupied by the electrons. By the MF theory, this spectral density is given by
$$A^<(k,\omega )=f(\omega )[u^2(k)\delta (\omega E_k^+)+v^2(k)\delta (\omega E_k^{})]$$
(8)
with $`u^2(k)=1v^2(k)=[1+(ϵ_kϵ_{k+Q})/E_k]/2`$. Fig. 1 shows the Fermi-level electron distributions in momentum space obtained by integrating $`A^<(k,\omega )`$ in a energy window $`|\omega \mu |<0.15t`$ at temperature $`T=0.1t`$. At low doping, the FS (in the first quadrant of the Brillouin zone) appears as two small pockets centered at $`(\pi ,0)`$ and $`(0,\pi )`$. With increasing doping, these pockets extend to larger squares. Such an evolution of the pockets reflects the occupation of the doped electrons at the upper Hubbard band. On the other hand, a small pocket begins to form around $`(\pi /2,\pi /2)`$ at $`x0.1`$ and grows up with increasing doping. The formation of this pocket stems from the contribution of the lower Hubbard band; with increasing doping, the energy gap decreases and meanwhile the lower Hubbard band shifts up toward to the FS. At high doping, the energy gap closes up, therefore the FS is a single curve centered at $`(\pi ,\pi )`$ ($`x`$ = 0.18). Clearly, the present calculation reproduces the results of the existing calculations based on the doping-dependent $`U`$ Hubbard model and is in agreement with the experimental observations. By the $`tt^{}t^{\prime \prime }U`$ model, a constant $`U`$ actually results in a nearly constant energy gap in a wide region of the doping concentration, and hence cannot explain the FS evolution. Though $`U`$ is constant in the present calculation, but because of the LRCE, the energy gap decreases remarkably with increasing doping. We will discuss this problem later again.
In Fig. 2, we depict the MF Néel temperature $`T_N`$ (solid line with circles) as a function of the doping concentration $`x`$, and compare it with the result of the extended Hubbard model (EHM, in which the LRCE is not included) and the experimental data. It is seen that the difference between the present calculation with LRCE and the EHM is remarkable. By the EHM, the dependence of $`T_N`$ on $`x`$ is very weak in the region $`0<x<0.2`$. In contrast to the EHM, $`T_N`$ given by the present calculation with LRCE decreases distinctly with increasing doping. In particular, it drops sharply at $`x0.141`$. This doping concentration corresponds to the quantum critical point at which the zero-temperature AF transition terminates and is consistent with the experimental observation. Also, the envelope of $`T_N`$ with LRCE is in fairly good agreement with the experimental results for the pseudogap $`\mathrm{\Delta }_{\mathrm{pg}}`$ and its onset temperature $`T^{}`$. The MF Néel transition is associated with the appearance of the local magnetization, and hence the comparison between the MF $`T_N`$ and $`T^{}`$ or $`\mathrm{\Delta }_{\mathrm{pg}}`$ is meaningful. The MF order parameter is a measure of the local order of the real system with long-wavelength fluctuation, similar to the case of superconducting ordering.
Why does the LRCE lead to $`T_N`$ decreasing with increasing doping? To answer this question, we here consider the excess hopping parameter $`t_{ij}^x`$ due to the Coulomb exchange. This quantity can be written as
$$t_{ij}^xt_l^x=\frac{v_{ij}}{2N}\underset{k}{}c_k^{}c_kS_l(k)$$
(9)
where $`S_l(k)=\mathrm{cos}(l_xk_x)\mathrm{cos}(l_yk_y)+\mathrm{cos}(l_yk_x)\mathrm{cos}(l_xk_y)`$ and again $`l=(l_x,l_y)`$ represents a vector of length $`|𝐢𝐣|`$. The doping dependence of $`t_{ij}^x`$ originates from the Fermi distribution $`c_k^{}c_k`$. At low temperature, the contribution to the integral in Eq. (9) comes from the Fermi area. The Fermi surface varies with changing doping concentration. Especially, with increasing doping, the Fermi area increases considerably in the regions near the points $`(\pi ,0)`$ and $`(0,\pi )`$ since at which the upper Hubbard band reaches its minimum and where the energy dispersion is nearly flat. Therefore, the variation of $`t_l^x`$ stems predominantly from the integral in these regions. On the other hand, with the change of the Fermi surface, the factor $`S_l(k)`$ varies slowly for small $`l`$, but rapidly for large $`l`$. It is then easy to see the variation of $`t_l^x`$ with increasing doping at low temperature: for example, $`t_{(1,0)}^x`$ is almost unchanged for $`S_{(1,0)}(\pi ,0)=0`$; $`t_{(1,1)}^x`$ has a large variation because of $`S_{(1,1)}(\pi ,0)=2`$. Since the original hopping parameter $`t_{(1,1)}`$ is negative, the magnitude of the effective hopping parameter $`t_{(1,1)}+t_{(1,1)}^x`$ should be enhanced with increasing doping. For large $`l`$, because of the cancellation from the destructive factor $`S_l(k)`$, not only the variation, but also the magnitude of $`t_l^x`$ are negligible small. At high temperature, since the integral in Eq. (9) is taken over a spreading area wider than the Fermi area, the behavior of $`t_l^x`$ is not so intuitive. Fig. 3 exhibits the variations of the magnitudes of the effective hopping parameters $`t_l+t_l^x`$ as functions of the doping concentration $`x`$ at temperature $`T=0.53t`$. The quantity $`\delta t_l(x)`$ is defined as
$$\delta t_l(x)=[t_l^x(x)t_l^x(0)]\mathrm{sgn}[t_l+t_l^x(0)].$$
(10)
From Fig. 3, it is seen that except for $`\delta t_{(1,0)}`$ that is nearly constant, the other parameters vary approximately linearly with $`x`$. Especially, among those effective hopping parameters only the magnitude of $`t_{(1,1)}+t_{(1,1)}^xt^{}`$ increases appreciably with increasing $`x`$. This is consistent with the above analysis. It is know that a large ratio $`t^{}/t`$ can destroy the AF instability at weak $`U`$. This is also true for strong on-site interactions. For large $`t^{}/U`$, the next n.n. AF coupling constant is $`J^{}=4t^2/U`$, which leads to the frustration of AF ordering in the square lattice. It is therefore clear that the AF order parameter and the transition temperature decrease with increasing doping concentration.
Shown in Fig. 4 is the parameter $`t_l^x`$ as a function of distance at the MF $`T_N=0.39t`$ at $`x=0.1`$. It is seen that only the magnitudes of the first 4 parameters are appreciable; from the fifth to the one at distance 4, all of them are very small; after that all other parameters are negligible small. In a wide region on the temperature-doping concentration phase diagram, the parameter $`t_l^x`$ behaves almost the same as shown in Fig. 4, with only a visual change in the next n.n hopping parameter. From the behavior of the parameter $`t_l^x`$, the range of the predominant Coulomb exchange effect seems to be shorter than 3. Of course, these results are obtained for the particular sets of the input parameters ($`t_l`$’s, $`U`$, $`V_1`$ and $`d`$). But these parameters are reasonable. In the present calculation, all the input parameters were not fine-tuned.
In summary, we have investigated the long-range Coulomb effect on the antiferromagnetism in the electron-doped cuprates using the mean-field theory. Due to the Coulomb exchange, the magnitude of the effective next nearest neighbor hopping parameter in especial increases appreciably with increasing the electron doping concentration. This leads to stronger frustration to the AF ordering at higher doping concentration. Therefore the transition temperature decreases with increasing doping. Consequently, the AF phase terminates at a doping concentration $`x0.14`$ in consistent with experiments. The present calculation gives a reasonable explanation to the doping dependence of the onset temperature of the AF pseudogap as well as to the FS evolution in the electron-doped cuprate Nd<sub>2-x</sub>Ce<sub>x</sub>CuO<sub>4</sub>.
The author thanks Prof. C. S. Ting for useful discussions on the related problems. This work was supported by Natural Science Foundation of China under grant number 10174092 and by Department of Science and Technology of China under grant number G1999064509. |
warning/0507/quant-ph0507149.html | ar5iv | text | # On local-hidden-variable no-go theorems
## 1 Introduction
In 1935, Einstein, Podolsky and Rosen (EPR) wrote their famous paper “Can quantum-mechanical description of physical reality be considered complete?” , wherein the authors answered the question in the negative. Completeness was then considered an important criterion regarding the validity of any physical theory of Nature, and would still be for any proponent of determinism and a classical view of the universe. Thus, their paper was a great threat to the validity of quantum mechanics (QM).
Completeness is defined in the EPR paper as “every element of the physical reality must have a counterpart in the physical theory”, where physical reality should be interpreted as “If, without in any way disturbing a system, we can predict with certainty (i.e., with probability equal to unity) the value of a physical quantity, then there exists an element of physical reality corresponding to this physical quantity.” EPR used the correlations obtained from bipartite measurements of an entangled state to claim that position and momentum can, do, and in fact must have simultaneous realities. More precisely, they used a bipartite entangled state and separated the two subsystems into space-like-separated regions. In their particular case, if one was to perform a position measurement on the first subsystem, from the result we could determine the exact position of the second subsystem. The same could be done for momentum. From their definition of physical reality and since, according to relativity, the space-like separation of the two subsystems prevents any interaction between them, EPR concluded that position and momentum must have simultaneous realities. From the fact the there are no interactions between the subsystems, they also concluded that these realities existed all along, from the time of the separation of the subsystems. Since the formalism of QM precludes such a description of reality, they were forced to conclude that QM cannot be considered a complete theory of Nature. For EPR, Nature possesses local hidden variables (LHVs), that can or cannot be known, which determine the behavior of a system under any measurement. In this picture, a measurement only reveals a pre-existing property of Nature, while the usual consensus in quantum mechanics is that a measurement forces a property of the system into existence. Einstein then spent the rest of his life in search of such a theory.
The LHV paradigm is the straightforward mathematical representation of local realism. In a LHV model, the hidden variables are set according to some probability distribution at the creation of a state. They can only be accessed experimentally though measurement, which perturbs the state and possibly alters the hidden variables. If we were to know the actual values of the hidden variables, we could fully predict the behaviour of the system under any measurement of an element of reality, but this lack of knowledge forces to average over the possible values of the LHVs. Thus the Heisenberg uncertainty principle is not violated and the probabilistic structure of QM is preserved. The locality criterion means that the outcome of a measurement on space-like separated systems cannot be correlated in any way other than through the original hidden variables, which are fixed once the state is created and change only according to local operations.
Even though the EPR paper contained logical errors , it was still the firmest attack against a quantum description of the physical world. Even the mightiest and hardiest defender of QM, Bohr , was not able to firmly put the EPR argument to rest . It took the better part of three decades for a complete refutation of the EPR argument to be put forth: in , Bell laid down the most general LHV model for a particular measurement setup. He then showed that the expectation value of the measurement operator could not, in *any* LHV model, come near the actual value predicted by QM. Experiments have confirmed the correctness of the predictions of QM .
The work of Bell has been called “the most profound discovery of science” . At least, it could easily be argued that Bell’s theorem has changed our view of Nature in the same way as Newton’s classical mechanics and Einstein’s relativity has. From his work, we now know that conjugate operators do not have simultaneous existence and only a measurement can force a state which is not an eigenstate of the operator to bring a value to the physical quantity attributed to the operator into existence.
In this paper, we present different forms of no-go theorems concerning LHV theories. To the knowledge of the author, only three forms of refutation exist, and a comprehensive comparative study has not yet been published. It should be noted that this paper does not have the goal of being an exhaustive survey of all the LHV no-go theorems. In Section 2, we will give formal definitions to the three forms of no-go theorems, namely Bell theorems, Bell theorems without inequalities<sup>1</sup><sup>1</sup>1Also called Bell inequalities without probabilities, Bell inequalities without inequalities or all-versus-nothing refutation of EPR. and pseudo-telepathy. Section 3 will then contain a discussion of the similarities and differences of the three forms.
## 2 The three forms
### 2.1 Bell theorems
The first proof that the physical world could not be described by any LHV theory came from Bell in 1964 in the form of an inequality . Bell bounded the absolute value of the expectation value of a specific operator in *any* LHV model and showed that quantum mechanics violates this bound. As such, we define a Bell theorem as follows.
###### Definition 1 (Bell theorem).
A Bell theorem is a set of multipartite measurements on an entangled state where the correlations obtained from the measurements cannot be reproduced by any local classical model where no communication between the participants is allowed.
#### 2.1.1 An example of a Bell theorem
The example we are to discuss here is often thought as the generic Bell theorem and was put forth as an experimental proposal by Clauser, Horne, Shimony and Holt . Suppose $`A_1`$ and $`A_2`$ are measurement settings on Alice’s apparatus and $`B_1`$ and $`B_2`$ are measurement settings on Bob’s. Given that the value of the outcome to each measurement lies between $`1`$ and 1, we can easily bound the value of the operator $`A_1B_1+A_1B_2+A_2B_1A_2B_2`$ in any LHV theory. Since the equation is linear in every variable and since every variable must be independent of one another (locality constraint) the maximum will be reached with every variable taking an extremal value, 1 or $`1`$. Therefore, we have $`|A_1B_1+A_1B_2+A_2B_1A_2B_2|2`$. In quantum mechanics, we can find appropriate measurements on a singlet state, $`(|+|+)/\sqrt{2}`$, to yield $`|A_1B_1+A_1B_2+A_2B_1A_2B_2|=2\sqrt{2}`$. The details are left as an exercise to the reader. From the fact that no LHV model can obtain a higher expectation value than 2, while QM can, is a clear proof that quantum correlations cannot be simulated by a LHV theory. If we now take the predictions of QM to be correct, we have to forsake the search for a local realistic description of reality.
### 2.2 Bell theorems without inequalities
A Bell theorem appeals to a statistical argument and, as such, does not look attractive to many. Thus, a more direct rebuttal of LHVs was sought for. The first such proof came from Heywood and Redhead , where the authors aimed to propose an experimental verification of the Kochen-Specker theorem to reject locality<sup>2</sup><sup>2</sup>2The first example is often wrongly attributed to Greenberger, Horne and Zeilinger .. It is important to note that Bell inequalities without inequalities are not experimental proposals which would rule out any LHV model in only one run. As Asher Peres once said : “The list of authors \[who has made this mistake\] is too long to give explicitly, and it would be unfair to give only a partial list.” We will discuss this in more detail in Section 3.1.
###### Definition 2 (Bell theorems without inequalities).
A Bell theorem without inequalities is a set of multipartite measurements on an entangled state where any local classical model, which is to attempt to simulate the probability distribution of the outputs given by quantum mechanics, will attribute a non zero probability to a measurement outcome that is forbidden by quantum mechanics or will never produce certain outcomes which are predicted with a non zero probability in quantum mechanics.
#### 2.2.1 An example of a Bell theorem without inequalities
For this example, we will give Brassard’s rendition of Hardy’s proof . Let us start with the state $`(|+|++|+)/\sqrt{3}`$ along the $`z`$ axis. Let say that Alice and Bob are now given the choice of performing either the $`\sigma _z`$ or $`\sigma _x`$ measurement. According to QM, if Alice and Bob are to measure $`\sigma _x\sigma _x`$, then they will receive the output $``$ with probability $`1/12`$. Let us now assume that the state has LHVs that will produce a $``$ output on a $`\sigma _x\sigma _x`$ measurement. From the criteria of locality and realism, we now have that any local $`\sigma _x`$ measurement on this particular state will produce the output $``$. Let us now see what happens if Alice and Bob are to measure $`\sigma _x\sigma _z`$ or $`\sigma _z\sigma _x`$. From the predictions of QM, the state should never be allowed to produce $``$ as output. Once again according to the criteria of locality and realism, we are forced to conclude that upon a local $`\sigma _z`$ measurement, the state will produce $`+`$ as output. Therefore, a $`\sigma _z\sigma _z`$ measurement on this particular instance is bound to output $`++`$, which is forbidden by QM. So in order for the LHV theory to output $``$ on a $`\sigma _x\sigma _x`$ measurement with a non zero probability, it will also output $`++`$ on a $`\sigma _z\sigma _z`$ measurement with non zero probability.
### 2.3 Pseudo-telepathy
Pseudo-telepathy was first defined, although not yet termed as such, in a paper by Brassard, Cleve and Tapp, where they turned the Deutsch-Jozsa algorithm into a distributed form to show that an exponential number of bits of communication was needed to simulate the correlations of a certain number of singlet states . For a complete survey, please refer to .
###### Definition 3 (Game).
A *bipartite game* $`G=(X,Y,R)`$ is a set of inputs $`X=X^{(A)}\times X^{(B)}`$, a set of outputs $`Y=Y^{(A)}\times Y^{(B)}`$ and a relation $`RX^{(A)}\times X^{(B)}\times Y^{(A)}\times Y^{(B)}`$.
###### Definition 4 (Winning Strategy).
A *winning strategy* for a bipartite game $`G=(X,Y,R)`$ is a strategy according to which for every $`x^{(A)}X^{(A)}`$ and $`x^{(B)}X^{(B)}`$, Alice and Bob output $`y^{(A)}`$ and $`y^{(B)}`$ respectively such that $`(x^{(A)},x^{(B)},y^{(A)},y^{(B)})R`$.
###### Definition 5 (Pseudo-telepathy).
We say that a bipartite game $`G`$ exhibits *pseudo-telepathy* if bipartite measurements of an entangled quantum state can yield a winning strategy, whereas no classical strategy that does not involve communication is a winning strategy.
The extension to the multipartite case is trivial. The generalization of Definition 5 can be translated into a set of multipartite measurements on an entangled state where any local classical model which is to attempt to produce outputs that are not forbidden by quantum mechanics will fail. We refer to this phenomenon by the term pseudo-telepathy because of its shortness and its illustrativeness. To someone who has no knowledge of quantum mechanics, and thus still believes in local realism, these correlations between the measurement outcomes can only be explained by some communication between the sub-systems. Thus if we make these measurements in space-like separated regions the sub-systems act as if they could telepathically communicate instantaneously.
#### 2.3.1 An example of a pseudo-telepathy game
We present here the pseudo-telepathy game generally known as the Magic Square game . The participants, namely Alice and Bob, are each presented with a question: a random trit $`x^{(A)}\{0,1,2\}`$ and $`x^{(B)}\{0,1,2\}`$ respectively. They must produce three bits each, $`y_1^{(A)}`$, $`y_2^{(A)}`$, $`y_3^{(A)}`$ and $`y_1^{(B)}`$, $`y_2^{(B)}`$, $`y_3^{(B)}`$ respectively. In order for them to win, $`y_1^{(A)}+y_2^{(A)}+y_3^{(A)}`$ must be even, $`y_1^{(B)}+y_2^{(B)}+y_3^{(B)}`$ must be odd and $`y_{x^{(B)}}^{(A)}`$ must equal $`y_{x^{(A)}}^{(B)}`$. In order for them to have a winning strategy without resorting to QM or communicating, Alice and Bob must share a $`3\times 3`$ table of 0s and 1s such that the sum of the elements in each row is even and the sum of the elements in each column is odd. A simple parity argument shows that such a table is impossible. On the other hand, if Alice and Bob are allowed to share entanglement, they can find a strategy which, while convoluted, can be explicitely constructed. It should be noted that they cannot construct the $`3\times 3`$ table required classically, but they can use the correlations of the quantum state to give three bits each respecting the above conditions.
## 3 Discussion of the similarities and differences
### 3.1 Similarities
The obvious similarity between the three forms of no-go theorems is that they all reject a local realistic description of Nature. Therefore, the physical world cannot be described by any LHV theory. As mentioned in Section 1, all three require many runs of the same experiment, with settings chosen at random on each run, to rule out any LHV model. A LHV might be lucky and answer correctly for many runs. Since a LHV theory can only succeed with a marginal probability of success, but still can succeed, we can only collect overwhelming evidence against LHVs. The only rejection we can make in one run is of quantum mechanics itself. If we consider perfect apparatus, then the production of a forbidden output by quantum mechanics, in a Bell theorem without inequalities or pseudo-telepathy experiment, would invalidate the correctness of quantum mechanical predictions.
It is also interesting to remark that every pseudo-telepathy game is a Bell theorem without inequalities and every Bell inequality without inequalities is a Bell theorem. These facts follows from Definition 1, 2 and 5. In pseudo-telepathy, the fact that no LHV strategy can always output within the relation $`R`$ can be seen as giving a non zero probability to an output that is forbidden in quantum mechanics—quantum mechanics being able to never give this output. And from Definition 2 and Definition 5, we clearly have a set of measurements on an entangled state where the correlations obtained from the measurements cannot be reproduced by any LHV model, hence falling into Definition 1.
### 3.2 Differences
Bell theorems and pseudo-telepathy are in a sense only quantitatively different. Pseudo-telepathy is simply an inequality where the quantum violation of the inequality reaches the maximal algebraic value. We saw in Section 2.3.1 an example where we appeared not to be concerned with reaching the maximal value of an inequality, but even this example can be converted as an inequality. The exercise is left to the reader.
It might be tempting to think that Bell theorems without inequalities are the same as pseudo-telepathy. In fact, the list of people who have made this error might also be too long to include here. The reason is that in both of these paradigms, we are concerned with proving a no-go theorem by a contradiction. However, it was proven recently that no pseudo-telepathy game can exist where the participants share only on entangled qubit, a $`2\times 2`$ system . For pseudo-telepathic correlations to arise we need at least a $`3\times 3`$ system or a $`2\times 2\times 2`$ system . Therefore, Hardy’s state cannot yield correlations strong enough for pseudo-telepathy. As a consequence, pseudo-telepathy and Bell theorems without inequalities are different paradigms. The difference is buried subtly in Definition 2 and Definition 5. In the first case, we require the LHV model to be able to generate all the possible outputs according to QM without outputting a forbidden output, while in the second case the requirement is weaker; we only need to produce outputs that can be generated by QM. In other words, in pseudo-telepathy, we only need to avoid forbidden outputs, but we do not need to produce every possible output. Hardy’s proof relies on the fact that once in a while the experiment will generate $`++`$ on a $`\sigma _x\sigma _x`$ measurement, but for a pseudo-telepathy game it is of no help.
## 4 Conclusion
We have seen the formal definition of the three forms of LHV no-go theorems and how they differ one from another. We have thus shown that it is important to look closely at our models when we want to describe nature, for we could have seen that there is a qualitative difference between Hardy’s no-go theorems and most of the other Bell theorems without inequalities. In this hierarchy of no-go theorems, pseudo-telepathy is the stronger refutation of the local realistic viewpoint as it lies at the top of the LHV no-go theorem hierarchy. Not every Bell theorem is a Bell theorem without inequalities and not every Bell theorem without inequalities is a pseudo-telepathy game, while the converse is true. There exists states that can generate correlations strong enough for a Bell theorem without inequalities while the same state cannot yield a pseudo-telepathy game .
## Acknowledgements
The author would like to thank Gilles Brassard for many invaluable discussions on the subject and Anne Broadbent and Richard MacKenzie for very helpful comments. |
warning/0507/hep-th0507034.html | ar5iv | text | # AdS/CFT Casimir Energy for Rotating Black Holes
\]Research supported in part by DOE grant DE-FG03-95ER40917 and NSF grant INTO3-24081.
(July 4, 2005)
## Abstract
We show that if one chooses the Einstein Static Universe as the metric on the conformal boundary of Kerr-AdS spacetime, then the Casimir energy of the boundary conformal field theory can easily be determined. The result is independent of the rotation parameters, and the total boundary energy then straightforwardly obeys the first law of thermodynamics. Other choices for the metric on the conformal boundary will give different, more complicated, results. As an application, we calculate the Casimir energy for free self-dual tensor multiplets in six dimensions, and compare it with that of the seven-dimensional supergravity dual. They differ by a factor of $`5/4`$.
preprint: DAMTP-2005-60 MIFP-05-15 hep-th/0507034
An interesting application of the AdS/CFT correspondence mald ; guklpo ; wit is to consider for the bulk solution a rotating Kerr-anti-de Sitter spacetime. As was discussed in hawhuntay , the boundary theory in this case should describe a conformal field theory in a rotating Einstein universe, allowing one in principle to study the effects of rapid rotation on the thermodynamics of the system. Of particular interest is the behaviour of the conformal field theory (CFT) on the boundary as the rotation parameters achieve their maximal values consistent with the non-existence of closed timelike curves. The Chronology Protection Conjecture of Hawking hawk suggests that going beyond this limit should be physically impossible. This may be associated with a divergence of the free energy of the CFT as one approaches the limit, and with a possible non-unitarity of the CFT if one passes beyond it.
The study of the thermodynamics of such rotating systems is quite involved and subtle, both in the bulk, and in the passage to the boundary theory. In the bulk, there are subtleties concerning the definition of the energy, or mass, of the rotating AdS black hole. As we showed in gibperpop , it is important that one evaluate the energy and the angular velocities with respect to a frame that is asymptotically non-rotating at infinity, in order to obtain quantities that satisfy the first law of thermodynamics,
$$dE=TdS+\mathrm{\Omega }^idJ_i.$$
(1)
In particular, a commonly considered frame that rotates relative to the asymptotically static frame, and which arose when the Kerr-AdS metrics were first constructed in Boyer-Lindquist type coordinates, is particularly inappropriate for defining the energy and angular velocities, since its asymptotic rotation rate is dependent on the angular-momentum parameters of the metric gibperpop2 .
In gibperpop2 we addressed the question of how one should map between the bulk thermodynamic variables and the corresponding variables on the boundary. We showed that using the standard UV/IR connection, the bulk variables $`(E,T,S,\mathrm{\Omega }^i,J_i)`$ that satisfy the first law as in (1) map straightforwardly into boundary variables $`(E,T,S,\omega ^i,j_i)`$ that also satisfy the first law, now with the addition of a pressure term,
$$de=tds+\omega ^idj_ipdv.$$
(2)
This mapping is implemented, for an $`n`$-dimensional bulk spacetime, by imposing the relations
$`e={\displaystyle \frac{l}{y}}E,\omega ^i={\displaystyle \frac{l}{y}}\mathrm{\Omega }^i,t={\displaystyle \frac{l}{y}}T,s=S,`$
$`j_i=J_i,v=𝒜_{n2}y^{n2},p={\displaystyle \frac{e}{(n2)v}}.`$ (3)
Here $`l`$ is the radius of the asymptotically-AdS spacetime, in the sense that $`R_{\mu \nu }=(n1)l^2g_{\mu \nu }`$, $`𝒜_{n2}`$ is the volume of the unit $`(n2)`$-sphere, and $`y`$ is the radius of a large sphere, with round spherical metric, near the boundary in AdS. As we emphasised in gibperpop2 , the natural choice of boundary metric is that of the Einstein Static Universe, ESU<sub>n-1</sub> i.e. the standard product metric on $`\times S^{n2}`$. In other words, one introduces a set of coordinates in which the Kerr-AdS metric approaches the canonical AdS<sub>n</sub> metric at infinity, in the form
$`d\overline{s}_n^2`$ $`=`$ $`(1+y^2l^2)dt^2+{\displaystyle \frac{dy^2}{1+y^2l^2}}+y^2d\mathrm{\Omega }_{n2}^2,`$ (4)
where $`d\mathrm{\Omega }_{n2}^2`$ is the metric on the unit $`(n2)`$-sphere. The conformal boundary is then located at $`y\mathrm{}`$, and the induced metric has the standard ESU<sub>n-1</sub> form
$$d\overline{s}_{n1}^2=dt^2+l^2d\mathrm{\Omega }_{n2}^2.$$
(5)
More precisely, in an AdS type conformal compactification, the physical metric is given by $`g=\mathrm{\Omega }^2\stackrel{~}{g}`$, where $`\mathrm{\Omega }=0`$ and $`d\mathrm{\Omega }0`$ on the timelike conformal boundary $``$, with $`\stackrel{~}{g}`$ being non-singular in a neighbourhood of $``$. The metric induced on the boundary $``$ is $`h=\stackrel{~}{g}|_{}`$, and depends on the choice of the conformal factor $`\mathrm{\Omega }`$. If $`\mathrm{\Omega }f\mathrm{\Omega }`$, for some function $`f`$ which is non-vanishing in a neighbourhood of $``$, then the boundary metric undergoes a conformal rescaling $`hf^2h`$. In order to obtain the standard metric on ESU<sub>n-1</sub>, one simply chooses $`\mathrm{\Omega }=\frac{l}{y}`$.
Our result in gibperpop2 refuted a recent surprising claim in cai , where it was asserted that in order to get thermodynamic quantities that satisfied the first law on the boundary, it was necessary to start from thermodynamic quantities in the bulk that did not satisfy the first law (in fact, the quantities defined relative to the rotating frame we mentioned earlier). As we showed, the puzzling results in cai were associated with the use of a somewhat unnatural conformal factor $`\mathrm{\Omega }`$, and hence a complicated boundary metric $`h`$ whose spatial sections were not round spheres, and which was not ultrastatic, i.e. $`g_{tt}`$ depended on spatial position.
In gibperpop2 , our principal interest was in the properties and thermodynamics of the bulk theories, and so it was not relevant to consider the contribution of Casimir energies on the boundary. Such terms do play an important rôle in the boundary CFT, and much work has been done on calculating them. For the case of Schwarzschild-AdS spacetime, the boundary Casimir contribution was evaluated in balakrau . Casimir calculations have also been performed for the case of rotating Kerr-AdS black holes, in awadjohn and more recently in skenpapa . The result obtained in these papers for the four-dimensional boundary theory is
$$E_{\mathrm{Casimir}}=\frac{3\pi ^2l^2}{4\kappa ^2}\left(1+\frac{(\mathrm{\Xi }_a\mathrm{\Xi }_b)^2}{9\mathrm{\Xi }_a\mathrm{\Xi }_b}\right),$$
(6)
where $`\kappa ^2/(8\pi )`$ is Newton’s constant, $`\mathrm{\Xi }_a=1a^2l^2`$, $`\mathrm{\Xi }_b=1b^2l^2`$, and $`a`$ and $`b`$ are the rotation parameters of the five-dimensional Kerr-AdS black hole given in hawhuntay .
The expression (6) for the Casimir energy of the boundary is a somewhat surprising result, since it depends on the rotation parameters $`a`$ and $`b`$ of the Kerr-AdS black hole. As we have argued above, the most natural conformal frame in which to formulate the boundary CFT is one in which the metric approaches the form (4), and the boundary metric is that of ESU<sub>n-1</sub>, and since this metric is manifestly independent of $`a`$ and $`b`$ the Casimir energy will necessarily also be independent of $`a`$ and $`b`$. Evidently, therefore, the conformal boundary metric chosen in awadjohn ; skenpapa is not the one we are advocating. In the remainder of this paper, we shall endeavour to convince the reader that our proposed choice of conformal boundary metric for the CFT dual to the Kerr-AdS metric is by far the most natural one, and that it has the very satisfactory feature that it leads to a genuinely constant Casimir energy in the boundary theory.
In fact, in order to demonstrate our point we need only collect a few results from previous papers. We shall mostly discuss the general $`n`$-dimensional case, since it is just as easy to discuss it generally as in any specific dimension. The general Kerr-AdS metrics were obtained in gilupapo1 ; gilupapo2 . The metrics have $`N[(n1)/2]`$ independent rotation parameters $`a_i`$ in $`N`$ orthogonal 2-planes. We have $`n=2N+1`$ when $`n`$ is odd, and $`n=2N+2`$ when $`n`$ is even. The metrics can be described by introducing $`N`$ azimuthal angles $`\varphi _i`$, and $`(N+ϵ)`$ “direction cosines” $`\mu _i`$ obeying the constraint
$$\underset{i=1}{\overset{N+ϵ}{}}\mu _i^2=1,$$
(7)
where $`ϵ=(n1)`$ mod 2. In Boyer-Lindquist type coordinates that are asymptotically non-rotating, the metrics are given by gilupapo1 ; gilupapo2
$`ds_n^2`$ $`=`$ $`W(1+r^2l^2)dt^2+{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{r^2+a_i^2}{\mathrm{\Xi }_i}}\mu _i^2d\phi _i^2`$ (8)
$`+{\displaystyle \underset{i=1}{\overset{N+ϵ}{}}}{\displaystyle \frac{r^2+a_i^2}{\mathrm{\Xi }_i}}d\mu _i^2+{\displaystyle \frac{Udr^2}{V2m}}`$
$`+{\displaystyle \frac{2m}{U}}\left(Wdt{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{a_i\mu _i^2d\phi _i}{\mathrm{\Xi }_i}}\right)^2`$
$`{\displaystyle \frac{l^2}{W(1+r^2l^2)}}\left({\displaystyle \underset{i=1}{\overset{N+ϵ}{}}}{\displaystyle \frac{r^2+a_i^2}{\mathrm{\Xi }_i}}\mu _id\mu _i\right)^2,`$
where
$`W`$ $``$ $`{\displaystyle \underset{i=1}{\overset{N+ϵ}{}}}{\displaystyle \frac{\mu _i^2}{\mathrm{\Xi }_i}},\mathrm{\Xi }_i1a_i^2l^2,`$
$`U`$ $``$ $`r^ϵ{\displaystyle \underset{i=1}{\overset{N+ϵ}{}}}{\displaystyle \frac{\mu _i^2}{r^2+a_i^2}}{\displaystyle \underset{j=1}{\overset{N}{}}}(r^2+a_j^2),`$ (9)
$`V`$ $``$ $`r^{ϵ2}(1+r^2l^2){\displaystyle \underset{i=1}{\overset{N}{}}}(r^2+a_i^2),`$ (10)
and it is understood, in the even-dimensional case, that $`a_{N+1}=0`$. The metrics satisfy $`R_{\mu \nu }=(n1)l^2g_{\mu \nu }`$. The constant-$`r`$ spatial surfaces at large distance are inhomogeneously distorted $`(n2)`$-spheres. Making the coordinate transformations
$$\mathrm{\Xi }_iy^2\widehat{\mu }_i^2=(r^2+a_i^2)\mu _i^2,$$
(11)
where $`_i\widehat{\mu }_i^2=1`$, the metrics at large $`y`$ approach the standard AdS form given in (4), where
$$d\mathrm{\Omega }_{n2}^2=\underset{k=1}{\overset{N+ϵ}{}}d\widehat{\mu }_k^2+\underset{k=1}{\overset{N}{}}\widehat{\mu }_k^2d\phi _k^2,$$
(12)
with round $`(n2)`$-spheres of volume $`𝒜_{n2}y^{n2}`$ at radius $`y`$, where $`𝒜_{n2}`$ is the volume of the unit $`(n2)`$-sphere. Note, in particular, that the boundary metric does not depend on any of the black-hole parameters.
The boundary CFT will be defined on the surface $`y=`$constant at very large $`y`$. The Casimir energy in the boundary theory is clearly independent of the mass $`m`$ of the Kerr-AdS black hole, since the boundary metric does not depend upon $`m`$. Therefore, for convenience, we can evaluate it by first setting $`m=0`$ in (8), which implies that the metric becomes purely AdS itself, in an unusual coordinate system. As was shown in gilupapo1 ; gilupapo2 , if one now performs the coordinate transformation (11), the AdS metric becomes exactly the canonical metric given in (4). The calculation of the Casimir energy for the boundary of Kerr-AdS is therefore reduced to the standard calculation for the Einstein Static Universe, ESU<sub>n-1</sub>. The answer is obviously independent of the rotation parameters, since they do not appear in the boundary metric.
One may calculate the Casimir energy in a number of ways. In balakrau , it was shown that the use of the holographic stress tensor for the conformal anomaly for four-dimensional $`𝒩=4`$ super-Yang-Mills gives
$$E_{\mathrm{Casimir}}=\frac{3\pi ^2l^2}{4\kappa ^2}.$$
(13)
This value agrees with a direct calculation of the Casimir energy using zeta-function regularisation of the sums over energy eigenvalues for $`6N^2`$ scalar fields, $`4N^2`$ Weyl spinor fields and $`N^2`$ vector fields on $`S^3`$. For our choice of conformal boundary metric, (13) is therefore the Casimir energy in the four-dimensional boundary theory dual to the five-dimensional Kerr-AdS metric. Combined with our calculation of the bulk energy term obtained in gibperpop , the complete CFT energy is given by
$$E_{\mathrm{tot}}=\frac{2\pi ^2m(2\mathrm{\Xi }_a+2\mathrm{\Xi }_b\mathrm{\Xi }_a\mathrm{\Xi }_b)}{\kappa ^2\mathrm{\Xi }_a^2\mathrm{\Xi }_b^2}+\frac{3\pi ^2l^2}{4\kappa ^2}.$$
(14)
The calculations in any other dimension proceed similarly, again yielding Casimir energies that are necessarily independent of the black hole rotation parameters. For example, in the case of a seven-dimensional Kerr-AdS bulk spacetime, we obtain a total boundary energy given by
$$E_{\mathrm{tot}}=\frac{2m\pi ^3}{\kappa ^2(_i\mathrm{\Xi }_i)}\left(\frac{1}{\mathrm{\Xi }_1}+\frac{1}{\mathrm{\Xi }_2}+\frac{1}{\mathrm{\Xi }_3}\frac{1}{2}\right)\frac{5\pi ^3l^4}{16\kappa ^2},$$
(15)
where the first term is the bulk energy that we calculated in gibperpop , and we have read off the Casimir term by setting $`a=0`$ in equation (51) of awadjohn . This value, which came from the use of the holographic stress tensor, should be compared with the value obtained directly by zeta-function regularisation of the sums over energy eigenvalues of $`20N^3`$ scalars, $`8N^3`$ Weyl fermions and $`4N^3`$ self-dual 3-forms, making up $`4N^3`$ copies basfrotse of the $`(2,0)`$ tensor multiplet. Using the energies and degeneracies tabulated in kutlar , we find that the Casimir energy for $`N_0`$ scalars, $`N_{1/2}`$ Weyl fermions and $`N_T`$ self-dual 3-forms is
$$E_{\mathrm{Casimir}}=\frac{(124N_0+1835N_{1/2}+11460N_T)}{241920l}.$$
(16)
For $`4N^3`$ multiplet, with the seven-dimensional AdS/CFT relation $`N^3=3\pi ^3l^5/(2\kappa ^2)`$, this gives
$$E_{\mathrm{Casimir}}=\frac{25\pi ^3l^4}{64\kappa ^2}.$$
(17)
Clearly this does not agree with the Casimir term in (15); the ratio is in fact $`5/4`$, which is not the same as the $`4/7`$ ratio conjectured in awadjohn . Presumably these differences reflect the absence of the non-renormalisation theorems that appear to account for the agreement of the various methods in the four-dimensional case.
We have demonstrated that by choosing the conformal boundary metric defined by taking $`y=`$constant, we obtain a very simple framework for describing the thermodynamics of the boundary field theory, with a very simple expression for the Casimir contribution to the energy, which is independent of the parameters in the Kerr-AdS metric. In particular, this means that the first law of thermodynamics continues to hold in a straightforward manner, when the Casimir energy is included. This contrasts with the more complicated situation in the case of the conformal boundary metric chosen in awadjohn ; skenpapa , where, as was shown in skenpapa , an additional diffeomorphism term must be included in the first law in order to compensate for the dependence of the boundary metric on the rotation parameters of the black hole.
The Boyer-Lindquist form of the boundary metric (in non-rotating coordinates) is obtained by taking the limit $`r\mathrm{}`$ of $`(l^2/r^2)ds^2`$, where $`ds^2`$ is given by (8). This yields
$`ds_{n1}^2`$ $`=`$ $`Wdt^2+l^2{\displaystyle \underset{i=1}{\overset{N+ϵ}{}}}{\displaystyle \frac{1}{\mathrm{\Xi }_i}}d\mu _i^2+l^2{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{1}{\mathrm{\Xi }_i}}\mu _i^2d\phi _i^2`$ (18)
$`{\displaystyle \frac{l^2}{W}}\left({\displaystyle \underset{i=1}{\overset{N+ϵ}{}}}{\displaystyle \frac{\mu _id\mu _i}{\mathrm{\Xi }_i}}\right)^2.`$
If we introduce new coordinates $`\widehat{\mu }_i`$, satisfying $`_i\widehat{\mu }_i^2=1`$, by the transformations
$$\mu _i=\sqrt{W\mathrm{\Xi }_i}\widehat{\mu }_i,$$
(19)
then the metric (18) becomes
$$ds_{n1}^2=Wd\overline{s}_{n1}^2,$$
(20)
where $`d\overline{s}_{n1}^2`$ is the standard metric on ESU<sub>n-1</sub>, given by (5) and (12), with $`W`$ expressed in terms of the $`\widehat{\mu }_i`$ as $`W=_i\mathrm{\Xi }_i\widehat{\mu }_i^2`$. This is the generalisation to arbitrary dimension of the five-dimensional demonstration in gibperpop that the $`r=`$constant “Boyer-Lindquist” boundary metric (18) and the $`y=`$constant Einstein static universe boundary metric (5) are related by a Weyl rescaling, and thus they lie in the same conformal class. When considering the first law, one varies the rotation parameters $`a_i`$. If one uses the Boyer-Lindquist boundary metric (18), this variation induces an infinitesimal change of the conformal factor, and hence a change of the Casimir contribution to the energy, as described in detail in five dimensions in skenpapa (see their equations (6.51)-(6.54)).
We are not, of course, saying that the more complicated results for the Casimir energies obtained in awadjohn ; skenpapa are wrong, but rather, that they are the consequence of having made a less felicitous choice for the conformal boundary metric, or equivalently, the conformal factor. As one can see from the calculations presented in skenpapa , although a coordinate transformation was performed in order to simplify the conformal boundary metric, it did not lead to as great a simplification as can be achieved by using the transformation (11).
We shall conclude with a remark on the action of the asymptotic symmetry group. Since we are working on the universal covering space $`\stackrel{~}{\text{AdS}}_n`$ of AdS$`{}_{n}{}^{}SO(n1,2)/SO(n2,1)`$, this is an infinite covering $`\stackrel{~}{SO}(n1,2)`$ of $`SO(n1,2)`$. The bulk energy $`E`$ and the angular momenta $`J_i`$ transform properly under the action of the asymptotic symmetry group as the components of a $`5\times 5`$ antisymmetric tensor $`J_{AB}`$, which may be thought of as an element of the Lie algebra $`\stackrel{~}{𝔰𝔬}(n1,2)`$. The transformation rule is just the adjoint action.
The transformation properties of the Casimir energy are more subtle. In general, the asymptotic group acts on the boundary by consometries, i.e. preserving the boundary metric $`h`$ only up to a conformal factor. Depending upon one’s choice of the boundary metric, i.e, of the representative in its conformal equivalence class, there may be a subgroup which acts by isometries. If, as we have done, we choose the Einstein static universe ESU<sub>n-1</sub> as boundary metric, then this subgroup is $`\times SO(n1)`$ and is maximal. Indeed, it is an infinite covering of the maximal compact subgroup of $`SO(n1,2)`$. The Casimir contribution to the energy is clearly invariant under the subgroup of isometries of the boundary metric, but it transforms in a well-defined but more complicated and non-trivial fashion under those elements of the asymptotic symmetry group which do not induce isometries. Choosing the Einstein static universe ESU<sub>n-1</sub> as boundary metric minimises these complications.
In summary, we have seen from previous work that the description of the bulk thermodynamics of rotating black holes in an AdS background is greatly simplified if one refers the energy and the angular velocities of the black hole to a coordinate frame that is non-rotating at infinity. Especially, one should not choose a rotating coordinate frame whose asymptotic angular velocity depends upon the parameters of the black hole. In this paper, we have furthermore shown that the description of the boundary CFT is greatly simplified if one likewise defines a conformal boundary metric that does not depend upon the parameters of the black hole. This is straightforwardly achieved by applying the coordinate transformation (11) to the asymptotically-static form of the Kerr-AdS metrics given in (8), and defining the boundary metric as the section $`y=`$constant as $`y`$ tends to infinity. By this means, all the complications associated with the unnecessary introduction of black hole parameter dependence of the thermodynamical quantities in the boundary theory are avoided.
###### Acknowledgements.
We thank Adel Awad, David Berman, Clifford Johnson and Ioannis Papadimitriou for discussions. C.N.P. thanks the Relativity and Cosmology group, Cambridge, for hospitality during the course of this work. |
warning/0507/hep-th0507144.html | ar5iv | text | # 1 Introduction and Summary
## 1 Introduction and Summary
Higher-spin gauge fields<sup>1</sup><sup>1</sup>1The web site http://www.ulb.ac.be/sciences/ptm/pmif/Solvay1proc.pdf contains the Proceedings of the First Solvay Workshop on Higher-Spin Gauge Fields, with a number of contributions, including that are more closely related to this work, and many references to the original literature. are an old and fascinating corner of Field Theory, with many open questions, some of which appear of direct relevance for a deeper understanding of String Theory. Only a few Lorentz representations play a role in the present models of the Fundamental Interactions (two-tensors, vectors, scalars, spinors, and often their super-partners), while the term “higher spin” qualifies in principle fields belonging to all types of representations that are more complicated than these. In practice, however, one often restricts the attention to rank-$`s`$ totally symmetric tensors $`\phi _{\mu _1\mathrm{}\mu _s}`$, that generalize the metric tensor fluctuation $`h_{\mu \nu }`$ and possess the gauge transformations
$$\delta \phi _{\mu _1\mathrm{}\mu _s}=_{\mu _1}\mathrm{\Lambda }_{\mu _2\mathrm{}\mu _s}+\mathrm{},$$
(1.1)
with parameters $`\mathrm{\Lambda }_{\mu _1\mathrm{}\mu _{s1}}`$ that are themselves totally symmetric tensors, of rank-$`(s1)`$, or to rank-$`n`$ totally symmetric spinor-tensors $`\psi _{\mu _1\mathrm{}\mu _n}`$, that generalize the gravitino field $`\psi _\mu `$ of supergravity and possess the gauge transformations
$$\delta \psi _{\mu _1\mathrm{}\mu _n}=_{\mu _1}ϵ_{\mu _2\mathrm{}\mu _n}+\mathrm{},$$
(1.2)
with parameters $`ϵ_{\mu _1\mathrm{}\mu _{n1}}`$ that are themselves totally symmetric spinor-tensors, of rank-$`(n1)`$. For this class of higher-spin fields, explicit statements can often be made efficiently and concisely for arbitrary values of $`s`$ or $`n`$. It should be borne in mind, however, that these types of fields do not exhaust all available possibilities in more than four dimensions, and indeed tensors of mixed symmetry are an important part of the massive spectrum of String Theory. Previous results concerning the free theory have been consistently generalized to cases of mixed symmetry, albeit necessarily in a less explicit fashion, and hence, for the sake of clarity, here we shall restrict ourselves to the totally symmetric case, leaving for future work a detailed analysis of the more involved cases of mixed symmetry. Abiding to common practice, from now on we shall simply use the term “spin” for the rank of bosonic tensors, or for the rank of Fermi fields augmented by $`1/2`$.
The conventional formulation for free totally symmetric tensor gauge fields $`\phi _{\mu _1\mathrm{}\mu _s}`$ was originally deduced by Fronsdal , in the late seventies, from the massive Singh-Hagen Lagrangians . The key feature of this formulation is the need for a pair of constraints, one on the gauge parameter, $`\mathrm{\Lambda }_{\mu _1\mathrm{}\mu _{s1}}`$, whose *trace* $`\mathrm{\Lambda }_{\mu _3\mathrm{}\mu _{s1}}^{}\eta ^{\mu _1\mu _2}\mathrm{\Lambda }_{\mu _1\mathrm{}\mu _{s1}}`$ is required to vanish, and one on the gauge field itself, whose *double trace* $`\phi _{\mu _5\mathrm{}\mu _s}^{\prime \prime }\eta ^{\mu _1\mu _2}\eta ^{\mu _3\mu _4}\phi _{\mu _1\mathrm{}\mu _s}`$ is also required to vanish. In a similar fashion, the conventional formulation for free totally symmetric spinor-tensors, due to Fang and Fronsdal , requires that both the *$`\gamma `$ \- trace* $`\overline{)}ϵ_{\mu _2\mathrm{}\mu _{n1}}`$ of the gauge parameter and the *triple $`\gamma `$ \- trace*<sup>2</sup><sup>2</sup>2For symmetric spinor-tensors two $`\gamma `$ \- traces are equivalent to a trace. of the spinor gauge field $`\overline{)}\psi _{\mu _4\mathrm{}\mu _n}^{}`$ vanish. While these constraints result in a consistent free dynamics, it is difficult to regard them as natural ingredients of a complete formulation of higher-spin gauge fields.
There is a body of evidence that consistent higher-spin interactions generally require that an infinite number of such fields be mutually coupled. For instance, the gravitational coupling for a single higher-spin field suffers from inconsistencies, the so-called Aragone-Deser problem, that can only be avoided in special cases, and most notably in (A)dS backgrounds. Only in such special circumstances can these fields be considered in isolation as in flat space. On the other hand, interacting systems of higher-spin fields are bound to be very complicated, and hence are not fully under control to date, but possess the intriguing virtue of being of intermediate complexity between ordinary low-spin Field Theory and String Theory. Still, many important results are now available. Most notably, in the presence of a cosmological term, the perturbative definition of higher-spin interactions around (A)dS spaces can avoid the difficulties long recognized for their naive flat-space couplings . This crucial observation led Vasiliev to formulate a consistent set of coupled non-linear equations for totally symmetric tensors, first in four dimensions and, more recently, in arbitrary dimensions as well. The Vasiliev equations (see also for relevant contributions along these lines) represent the most encouraging result on higher-spin gauge fields available to date, although they are clearly non-Lagrangian and more work is required to arrive at an off-shell formulation. They generalize both the frame formulation of Einstein gravity and the Cartan integrable systems that have long emerged from supergravity , extending them to allow for non-polynomial scalar couplings.
The Vasiliev equations are based on an extension of the frame formalism for gravity, and as a result their fields, forms valued in representations of the tangent Lorentz group, can simply accommodate Fronsdal’s constraints. Still, other possibilities should be explored at this stage, and in metric-like formulations the trace constraints of should naturally be absent. With this motivation in mind, in we showed that it is possible to extend the Fronsdal construction to allow for unconstrained gauge fields and parameters. A nice outcome was the prompt emergence of the geometry underlying the field equations, that become
$$\frac{1}{\mathrm{}^p}^{[p]}{}_{;\alpha _1\mathrm{}\alpha _{2p+1}}{}^{}=0,$$
(1.3)
for odd spins $`s=2p+1`$, and
$$\frac{1}{\mathrm{}^{p1}}^{[p]}{}_{;\alpha _1\mathrm{}\alpha _{2p}}{}^{}=0,$$
(1.4)
for even spins $`s=2p`$, where $``$ denotes the linearized higher-spin curvature introduced by de Wit and Freedman in and the bracketed suffix denotes that $`p`$ \- fold traces are taken. Being *non local* for spin $`s>2`$, both these equations and the corresponding Lagrangians are not easy to use, but their geometrical nature is nonetheless quite suggestive. Similar constructions for higher-spin fermions were also presented in , and the bosonic construction of was generalized to mixed symmetry tensors in .
Local field equations for unconstrained fields can actually be obtained without much effort. Confining our attention for simplicity to bosonic fields, it possible to show that the non-local geometric equations (1.3) and (1.4) are equivalent to simple non-Lagrangian systems involving a new field, a spin - $`(s3)`$ compensator $`\alpha _{\mu _1\mathrm{}\mu _{s3}}`$, that under gauge transformations transforms as
$$\delta \alpha _{\mu _1\mathrm{}\mu _{s3}}=\mathrm{\Lambda }_{\mu _1\mathrm{}\mu _{s3}}^{}.$$
(1.5)
As we shall review in the next Section, the resulting local compensator form of the higher-spin equations is equivalent to eqs. (1.3) and (1.4), and is actually suggested by String Field Theory . It can also be extended in a relatively simple fashion to (A)dS backgrounds, and similar results apply to its fermionic analog of .
The role of the unconstrained gauge symmetry in the interactions of higher-spin gauge fields is less clear at the moment, but there are clues that off-shell they will eventually make use of it. To wit, while the four-dimensional Vasiliev construction of was based on a generalization of the two-component formalism for gravity, and as a result is strictly tied to the Fronsdal form, this is not necessarily the case for the more recent construction of . As stressed in , in some respects the new formulation of can be regarded as a step forward in the direction of an off-shell formulation, and indeed it allows to drop the trace conditions. When this is done , at the free level one recovers very nicely the local compensator equations of . At the interacting level, however, this choice entails some subtleties that are spelled out in detail in , and further work is required to settle the issue of its consistency, although a forthcoming microscopic analysis of the Vasiliev equations lends further, independent support to the role of the unconstrained symmetry . Recently, additional evidence for the potential role of the unconstrained symmetry for the off-shell description of higher-spin gauge fields was also provided in .
Complete off-shell formulations for free symmetric tensors were already introduced some time ago by Pashnev and Tsulaia . Like their more recent fermionic counterparts of , these constructions rest on the BRST formalism and describe spin - $`s`$ symmetric tensors via additional fields whose number grows proportionally to $`s`$. The resulting free systems are rather complicated, but nonetheless in it was shown that a judicious elimination of most of the additional fields reproduces the geometric equations (1.3) and (1.4) in the compensator form of . The present letter is devoted to displaying far simpler constructions, minimal Lagrangians where the trace constraints of the conventional Fronsdal formulation are eliminated adding *only two* fields, the compensator $`\alpha `$ and a single spin - $`(s4)`$ symmetric tensor $`\beta `$ playing the role of a Lagrange multiplier. The next Section reviews briefly the results of , while the two remaining Sections are devoted, respectively, to the minimal bosonic Lagrangians and to their fermionic counterparts.
## 2 Geometric and Local Compensator Equations
As anticipated in the Introduction, in this paper we restrict our attention to an important class of higher-spin fields, totally symmetric (spinor-)tensors, a choice that has the virtue of allowing a relatively handy discussion. The Fronsdal equations are
$$\mathrm{}\phi \phi +^2\phi ^{}=\mathrm{\hspace{0.17em}0},$$
(2.1)
where $``$ will be often referred to as the Fronsdal operator. For spin one and two, eqs. (2.1) reduce to the Maxwell equation for the vector potential $`A_\mu `$ and to the linearized Einstein equation for the metric fluctuation $`h_{\mu \nu }`$, while novelties begin to emerge for spin 3.
Let us pause briefly to explain our notation. In this letter, as in , primes (or bracketed suffixes) denote traces, while all indices carried by the symmetric tensors $`\phi _{\mu _1\mathrm{}\mu _s}`$ and $`\mathrm{\Lambda }_{\mu _1\mathrm{}\mu _{s1}}`$, by the metric tensor $`\eta _{\mu \nu }`$ or by derivatives are left implicit. In order to fully profit from this shorthand notation, where all terms are meant to be totally symmetrized so that, for instance, $`\phi `$ stands for $`_{\mu _1}\phi _{\mu _2\mathrm{}\mu _{s+1}}+\mathrm{}`$, one need only get accustomed to a few rules, that is convenient to display again, correcting also a misprint in :
$$\begin{array}{cc}\hfill \left(^p\phi \right)^{}& =\mathrm{}^{p2}\phi +\mathrm{\hspace{0.17em}2}^{p1}\phi +^p\phi ^{},\hfill \\ \hfill ^p^q& =\left(\genfrac{}{}{0pt}{}{p+q}{p}\right)^{p+q},\hfill \\ \hfill \left(^p\phi \right)& =\mathrm{}^{p1}\phi +^p\phi ,\hfill \\ \hfill \eta ^k& =\eta ^{k1},\hfill \\ \hfill \left(\eta ^k\phi \right)^{}& =\left[D+\mathrm{\hspace{0.17em}2}(s+k1)\right]\eta ^{k1}\phi +\eta ^k\phi ^{}.\hfill \end{array}$$
(2.2)
We shall work throughout in $`D`$ dimensions, with a mostly positive space-time signature, and in this notation the Fronsdal Lagrangian is simply
$$_0=\frac{1}{2}\phi \left(\frac{1}{2}\eta ^{}\right),$$
(2.3)
where, as will be always the case in the following, evident contractions between different fields in Lorentz invariant monomials are left implicit. For instance, here $`\phi `$ is implicitly contracted with the Einstein-like tensor $`\frac{1}{2}\eta ^{}`$.
This formulation rests crucially on two constraints, that first emerge for spin 3 and 4. The first concerns the gauge parameter, $`\mathrm{\Lambda }`$, that enters the gauge transformation of $`\phi `$ in the conventional fashion, as $`\delta \phi =\mathrm{\Lambda }`$, but is to be traceless in order to guarantee the gauge invariance of (2.1), since
$$\delta =\mathrm{\hspace{0.17em}3}^{\mathrm{\hspace{0.17em}3}}\mathrm{\Lambda }^{}.$$
(2.4)
The second concerns the gauge field $`\phi `$ itself, that is to be doubly traceless. This peculiar restriction originates from the “anomalous” Bianchi identity for the Fronsdal operator,
$$\frac{1}{2}^{}=\frac{3}{2}^{\mathrm{\hspace{0.17em}3}}\phi ^{\prime \prime },$$
(2.5)
since indeed, even with a traceless $`\mathrm{\Lambda }`$, and up to partial integrations that will be always left implicit in the following, the Lagrangian (2.3) *is not* gauge invariant, but varies into
$$\delta _0=\frac{s}{2}\mathrm{\Lambda }\left(\frac{1}{2}^{}\right).$$
(2.6)
In we showed that, making use of the gauge field $`\phi `$ only, one can build a sequence of pseudo-differential analogs of $``$,
$$^{(k+1)}=^{(k)}+\frac{1}{(k+1)(2k+1)}\frac{^{\mathrm{\hspace{0.33em}2}}}{\mathrm{}}_{}^{(k)}{}_{}{}^{}\frac{1}{k+1}\frac{}{\mathrm{}}^{(k)},$$
(2.7)
such that
$$\delta ^{(k)}=\left(2k+1\right)\frac{^{\mathrm{\hspace{0.33em}2}k+1}}{\mathrm{}^{k1}}\mathrm{\Lambda }^{[k]}.$$
(2.8)
For spin $`s=2k1`$ or $`s=2k`$, $`^{(k)}=\mathrm{\hspace{0.17em}0}`$ is thus the minimal fully gauge invariant modification of the Fronsdal equation. Interestingly, this result is directly linked to the higher-spin curvatures $``$ introduced by de Wit and Freedman , and indeed the fully gauge invariant equations can be written in the more suggestive *geometric* fashion
$$\frac{1}{\mathrm{}^p}^{[p]}{}_{;\alpha _1\mathrm{}\alpha _{2p+1}}{}^{}=0,$$
(2.9)
for odd spins $`s=2p+1`$, and
$$\frac{1}{\mathrm{}^{p1}}^{[p]}{}_{;\alpha _1\mathrm{}\alpha _{2p}}{}^{}=0,$$
(2.10)
for even spins $`s=2p`$.
In this formulation gauge invariance is manifest, and does not require any constraints on the gauge parameter, but the equivalence to the Fronsdal formulation entails a few subtleties, that are spelled out in . In addition, the Bianchi identity is also modified, so that
$$^{(k)}\frac{1}{2k}_{}^{(k)}{}_{}{}^{}=\left(1+\frac{1}{2k}\right)\frac{^{\mathrm{\hspace{0.33em}2}k+1}}{\mathrm{}^{k1}}\phi ^{[k+1]},$$
(2.11)
and this suffices to show that, for every given spin, there exists a lowest value of $`k`$ such that the generalized Einstein tensors
$$𝒢^{(k)}=\underset{pk}{}\frac{(1)^p}{2^pp!\left(\genfrac{}{}{0pt}{}{k}{p}\right)}\eta ^p^{(k)[p]}$$
(2.12)
are divergence-free, and hence the Lagrangians
$$=\frac{1}{2}\phi 𝒢^{(k)}$$
(2.13)
are fully gauge invariant.
In we also showed that the geometric equations (2.9) and (2.10) can be always turned into the form
$$=\mathrm{\hspace{0.17em}3}^3,$$
(2.14)
with $``$ a *non-local construct* of $``$, bound to transform as $`\delta =\mathrm{\Lambda }^{}`$ under a gauge transformation. In other words, $``$ behaves as a *compensator* for the trace of the gauge parameter.
In , drawing also from String Field Theory , we explored the implications of allowing in the theory an *independent field* $`\alpha `$, denoted as the “compensator” and such that $`\delta \alpha =\mathrm{\Lambda }^{}`$. The key result of this analysis was that with $`\phi `$ and $`\alpha `$ one can arrive at the two local field equations
$``$ $`=`$ $`3^{\mathrm{\hspace{0.17em}3}}\alpha ,`$ (2.15)
$`\phi ^{\prime \prime }`$ $`=`$ $`4\alpha +\alpha ^{},`$ (2.16)
that can also be obtained truncating the bosonic triplet of , are nicely consistent with the Bianchi identity (2.5), but *are not* Lagrangian. As shown in (see also ), from these one can readily recover the non-local geometric equations (2.9) and (2.10) building a sequence of equations for the non-local extensions $`^{(k)}`$ of (2.7), since eq. (2.15) implies that
$$^{(k)}=(2k+1)\frac{^{\mathrm{\hspace{0.17em}2}k+1}}{\mathrm{}^k}\alpha ^{[k]},$$
(2.17)
and finally, after the minimal number of iterations needed to produce a trace of $`\alpha `$ not allowed for a given spin $`s`$, eqs. (2.9) and (2.10).
In a similar fashion starting from the fermionic Fang-Fronsdal operator
$$𝒮=i\left(\overline{)}\psi \overline{)}\psi \right),$$
(2.18)
one can define a sequence of non-local extensions of $`𝒮`$, directly linked to the bosonic ones according to
$$𝒮_{n+1/2}^{(k)}\frac{1}{2k}\frac{}{\mathrm{}}\overline{)}\overline{)}𝒮_{n+1/2}^{(k)}=i\frac{\overline{)}}{\mathrm{}}_n^{(k)}(\psi ),$$
(2.19)
with $`_n^{(k)}(\psi )`$ a non-local extension of the Fronsdal operator for the spinor-tensor $`\psi `$, and arrive eventually at non-local geometric equations. Even in this case, non-Lagrangian equations for spin - $`(n+1/2)`$ spinor-tensors $`\psi _{\mu _1\mathrm{}\mu _n}`$ involving a single spin - $`(n3/2)`$ compensator $`\xi _{\mu _1\mathrm{}\mu _{n2}}`$ were obtained in . In flat space they read
$`𝒮`$ $`=`$ $`2i^2\xi ,`$
$`\overline{)}\psi ^{}`$ $`=`$ $`2\xi +\xi ^{}+\overline{)}\overline{)}\xi ,`$ (2.20)
and can be obtained truncating the fermionic triplet introduced in .
## 3 Local Lagrangians for unconstrained bosons
In the previous Section we have reviewed the salient features of the non-local geometric equations of and their relation with the non-Lagrangian equations of . As shown in , the latter also follow, albeit in a somewhat indirect fashion, from the BRST construction of Pashnev and Tsulaia .
In this Section we would like to present a simple alternative: local Lagrangians for unconstrained spin - $`s`$ tensor fields $`\phi `$ that involve at most two additional fields, the spin - $`(s3)`$ compensator $`\alpha `$ of , and an additional spin - $`(s4)`$ field $`\beta `$ that acts as a Lagrange multiplier for the relation between the double trace $`\phi ^{\prime \prime }`$ and the compensator $`\alpha `$ in eq. (2.16). These “minimal” Lagrangians are closely related to the geometric equations (2.9) and (2.10) of via the compensator system (2.15) and (2.16), and for spin $`s=3`$ reduce to the result already presented in .
A minimal gauge invariant Lagrangian for spin - $`s`$ symmetric tensors $`\phi `$ can be nicely determined resorting to the familiar Noether procedure, that allows one to deal with this problem in a systematic fashion and has the additional virtue of clarifying the origin of the difficulty met in a naive approach beyond the spin - 3 case. Let us therefore begin by considering the Fronsdal expression (2.3), now written for a field $`\phi `$ *not subject* to any trace constraints, and let us vary it without enforcing any constraints on the gauge parameter $`\mathrm{\Lambda }`$. The resulting complete variation,
$$\begin{array}{cc}\hfill \delta _0=& 3\left(\genfrac{}{}{0pt}{}{s}{4}\right)\phi ^{\prime \prime }\mathrm{\Lambda }9\left(\genfrac{}{}{0pt}{}{s}{4}\right)\phi ^{}\mathrm{\Lambda }^{}+\frac{15}{2}\left(\genfrac{}{}{0pt}{}{s}{5}\right)\phi ^{}\mathrm{\Lambda }^{\prime \prime }\hfill \\ & +\mathrm{\Lambda }^{}\left(\genfrac{}{}{0pt}{}{s}{3}\right)\left\{\frac{3}{4}^{}\frac{3}{2}\phi +\frac{9}{4}\mathrm{}\phi ^{}\right\},\hfill \end{array}$$
(3.1)
comprises a number of terms depending on $`\mathrm{\Lambda }^{}`$, that can be canceled adding
$$\begin{array}{cc}\hfill _1=& \alpha \left(\genfrac{}{}{0pt}{}{s}{3}\right)\left\{\frac{3}{4}^{}\frac{3}{2}\phi +\frac{9}{4}\mathrm{}\phi ^{}\right\}\hfill \\ & +\mathrm{\hspace{0.17em}9}\left(\genfrac{}{}{0pt}{}{s}{4}\right)\alpha \phi ^{}\frac{15}{2}\left(\genfrac{}{}{0pt}{}{s}{5}\right)\alpha ^{}\phi ^{},\hfill \end{array}$$
(3.2)
that depends linearly on the compensator $`\alpha `$.
Additional terms depending on $`\mathrm{\Lambda }^{}`$ now present themselves in the resulting variation of $`_0+_1`$, but can be eliminated adding
$$\begin{array}{cc}\hfill _2=& \frac{9}{4}\left(\genfrac{}{}{0pt}{}{s}{3}\right)\alpha \mathrm{}^2\alpha \mathrm{\hspace{0.17em}27}\left(\genfrac{}{}{0pt}{}{s}{4}\right)\alpha \mathrm{}\alpha +\mathrm{\hspace{0.17em}45}\left(\genfrac{}{}{0pt}{}{s}{5}\right)(\alpha )^2\hfill \\ & +\frac{45}{2}\left(\genfrac{}{}{0pt}{}{s}{5}\right)\alpha \mathrm{}\alpha ^{}\mathrm{\hspace{0.17em}45}\left(\genfrac{}{}{0pt}{}{s}{6}\right)\alpha \alpha ^{},\hfill \end{array}$$
(3.3)
that is quadratic in $`\alpha `$, so that the final remainder is
$$\delta \left\{_0+_1+_2\right\}=3\left(\genfrac{}{}{0pt}{}{s}{4}\right)\left\{\phi ^{\prime \prime }\mathrm{\hspace{0.17em}4}\alpha \alpha ^{}\right\}\mathrm{\Lambda }.$$
(3.4)
These terms vanish for $`s<4`$, and are proportional to the gauge invariant expression given by eq. (2.16). A fully gauge invariant unconstrained Lagrangian is thus finally obtained introducing, from spin $`s=4`$ onwards, the single additional term
$$_3=\mathrm{\hspace{0.17em}3}\left(\genfrac{}{}{0pt}{}{s}{4}\right)\beta \left(\phi ^{\prime \prime }\mathrm{\hspace{0.17em}4}\alpha \alpha ^{}\right),$$
(3.5)
where the spin - $`(s4)`$ *Lagrange multiplier* $`\beta `$ transforms as $`\delta \beta =\mathrm{\Lambda }`$.
Summarizing, the complete Lagrangians for unconstrained spin - $`s`$ totally symmetric tensors are
$$\begin{array}{cc}\hfill =& \frac{1}{2}\phi \left(\frac{1}{2}\eta ^{}\right)\left(\genfrac{}{}{0pt}{}{s}{3}\right)\alpha \left\{\frac{3}{4}^{}\frac{3}{2}\phi +\frac{9}{4}\mathrm{}\phi ^{}\right\}\hfill \\ & +\mathrm{\hspace{0.17em}9}\left(\genfrac{}{}{0pt}{}{s}{4}\right)\alpha \phi ^{}\frac{15}{2}\left(\genfrac{}{}{0pt}{}{s}{5}\right)\alpha ^{}\phi ^{}+\frac{9}{4}\left(\genfrac{}{}{0pt}{}{s}{3}\right)\alpha \mathrm{}^2\alpha \hfill \\ & \mathrm{\hspace{0.17em}27}\left(\genfrac{}{}{0pt}{}{s}{4}\right)\alpha \mathrm{}\alpha +\mathrm{\hspace{0.17em}45}\left(\genfrac{}{}{0pt}{}{s}{5}\right)(\alpha )^2+\frac{45}{2}\left(\genfrac{}{}{0pt}{}{s}{5}\right)\alpha \mathrm{}\alpha ^{}\hfill \\ & \mathrm{\hspace{0.17em}45}\left(\genfrac{}{}{0pt}{}{s}{6}\right)\alpha \alpha ^{}+3\left(\genfrac{}{}{0pt}{}{s}{4}\right)\beta \left(\phi ^{\prime \prime }\mathrm{\hspace{0.17em}4}\alpha \alpha ^{}\right),\hfill \end{array}$$
(3.6)
and are invariant under the gauge transformations
$$\begin{array}{cc}& \delta \phi =\mathrm{\Lambda },\hfill \\ & \delta \alpha =\mathrm{\Lambda }^{},\hfill \\ & \delta \beta =\mathrm{\Lambda }.\hfill \end{array}$$
(3.7)
We can now move on to clarify the connection with the non-Lagrangian system of eqs. (2.15) and (2.16). The starting point are the field equations determined by (3.6),
$`\phi `$ $`:\mathrm{\hspace{0.17em}3}^{\mathrm{\hspace{0.17em}3}}\alpha {\displaystyle \frac{1}{2}}\eta (^{}{\displaystyle \frac{1}{2}}^2\phi ^{\prime \prime }\mathrm{\hspace{0.17em}3}\mathrm{}\alpha \mathrm{\hspace{0.17em}4}^{\mathrm{\hspace{0.17em}2}}\alpha {\displaystyle \frac{3}{2}}^{\mathrm{\hspace{0.17em}3}}\alpha ^{})`$
$`+\eta ^2(\beta +{\displaystyle \frac{1}{2}}\alpha +\mathrm{}\alpha {\displaystyle \frac{1}{2}}\phi ^{})=\mathrm{\hspace{0.17em}0},`$ (3.8)
$`\beta `$ $`:\phi ^{\prime \prime }\mathrm{\hspace{0.17em}4}\alpha \alpha ^{}=0,`$ (3.9)
$`\alpha `$ $`:6\mathrm{}^2\alpha +\mathrm{\hspace{0.17em}18}\mathrm{}\alpha +\mathrm{\hspace{0.17em}12}^2\alpha +\mathrm{\hspace{0.17em}3}\mathrm{}^2\alpha ^{}+\mathrm{\hspace{0.17em}3}^3\alpha ^{}`$
$`\mathrm{\hspace{0.17em}3}\phi ^{}^{}+\mathrm{\hspace{0.17em}2}\phi \mathrm{\hspace{0.17em}3}\mathrm{}\phi ^{}+\mathrm{\hspace{0.17em}4}\beta `$
$`+\eta (3\mathrm{}\alpha +\alpha \phi ^{}+\mathrm{\hspace{0.17em}2}\beta )=\mathrm{\hspace{0.17em}0},`$ (3.10)
and the issue is to show their equivalence to eq. (2.15). A general argument to this effect can be built as follows.
Let us begin by noticing that, when $`\beta `$ is on-shell, i.e. when eq. (3.9) is enforced, the $`\phi `$ equation becomes of the form
$$𝒜\frac{1}{2}\eta 𝒜^{}+\eta ^2𝒞=\mathrm{\hspace{0.17em}0},$$
(3.11)
where
$$\begin{array}{cc}\hfill 𝒜& =\mathrm{\hspace{0.17em}3}^{\mathrm{\hspace{0.17em}3}}\alpha ,\hfill \\ \hfill 𝒞& =\beta +\frac{1}{2}\alpha +\mathrm{}\alpha \frac{1}{2}\phi ^{},\hfill \end{array}$$
(3.12)
and that, under the same condition, the double trace of $`𝒜`$ vanishes *identically*. One can then take successive traces of (3.11): whereas the first relates $`𝒜^{}`$ to $`𝒞`$ and $`𝒞^{}`$, the higher ones yield relations of the form
$$\left(\eta ^2𝒞\right)^{[k]}=\eta ^2𝒞^{[k]}+\underset{i=1}{\overset{k}{}}\rho _{2i+1}\eta 𝒞^{[k1]}+\underset{ij=2}{\overset{k}{}}\rho _{2i1}\rho _{2j}𝒞^{[k2]}=\mathrm{\hspace{0.17em}0}.$$
(3.13)
It should be appreciated that the $`𝒞^{[i]}`$ are all independent and do not vanish identically. As a result, these never reduce to trivial identities, since the coefficients $`\rho _kD+\mathrm{\hspace{0.17em}2}(sk)`$ are positive. Therefore, denoting with $`p`$ the integer part of $`\frac{s4}{2}`$ and taking $`p+2`$ traces of (3.11) gives
$$\underset{ij=2}{\overset{p+2}{}}\rho _{2i1}\rho _{2j}𝒞^{[p]}=\mathrm{\hspace{0.17em}0},$$
(3.14)
and hence, finally, $`𝒞^{[p]}=\mathrm{\hspace{0.17em}0}`$. Making use of this relation in the $`[p+1]`$-th trace, and working backwards, one can convince oneself that *on shell all traces of $`𝒞`$, including $`𝒞`$ itself, vanish*. In the first trace of (3.11), this result gives $`𝒜^{}=\mathrm{\hspace{0.17em}0}`$, and then finally eq. (3.11) turns into the desired form (2.15)<sup>3</sup><sup>3</sup>3$`s=2`$ and $`D=2`$ is a well-known exception, since in that case the trace of $`𝒜\frac{1}{2}\eta 𝒜^{}`$ vanishes identically, giving no indications on $`𝒜^{}`$..
As we have seen, the field equation for the Lagrange multiplier $`\beta `$ is the condition that the double trace of the dynamical field be pure gauge, and plays a crucial role in linking these Lagrangian equations to the geometrical ones. On the other hand, (3.10) has not played any role so far, in particular in the relation between the local Lagrangian (3.6) and the higher-spin geometry. There is a reason for this: (3.10) is a *consequence* of the field equations for $`\phi `$ and $`\beta `$. Indeed, taking the divergence of eq. (3.8) and using eq. (3.9) in the result, one arrives at an expression proportional to (3.10). More precisely, indicating with $`𝒢_{\phi ,\beta }(\alpha )`$ the field operator for the compensator $`\alpha `$, one can see that
$$\{𝒜\frac{1}{2}\eta 𝒜^{}+\eta ^2𝒞\}=\frac{\eta }{4}𝒢_{\phi ,\beta }(\alpha ).$$
(3.15)
It is then clear that if $`\phi `$ and $`\beta `$ satisfy their field equations, $`𝒢_{\phi ,\beta }(\alpha )`$ is forced to vanish. That is to say, $`\alpha `$ is forced to satisfy its field equation as well.
Actually, the role of the field equation for the compensator $`\alpha `$ can be better appreciated if the dynamical field $`\phi `$ is coupled to an external source $`𝒥`$. In the Fronsdal case, the Lagrangian equation is
$$\frac{1}{2}\eta ^{}=𝒥,$$
(3.16)
and its divergence gives
$$\frac{1}{2}\eta ^{}=𝒥.$$
(3.17)
Hence, while in the Maxwell and Einstein cases $`^{}`$ vanishes identically, so that the sources must be *divergence free*, in the conventional formulation for spin 3 or higher only the *traceless part of the divergence* is forced to vanish. In it was shown that even this weaker condition suffices to ensure that only physical polarizations contribute to the exchange of quanta between sources. On the other hand, if an external source is introduced in our Lagrangian (3.6) via the standard coupling $`\phi 𝒥`$, the field equations become
$`𝒜{\displaystyle \frac{1}{2}}\eta +\eta ^2𝒞=𝒥,`$ (3.18)
$`𝒢_{\phi ,\beta }(\alpha )=\mathrm{\hspace{0.17em}0},`$ (3.19)
$`\phi ^{\prime \prime }\mathrm{\hspace{0.17em}4}\alpha \alpha ^{}=\mathrm{\hspace{0.17em}0},`$ (3.20)
where $``$ is defined by comparing with eq. (3.8). Combining the divergence of (3.18) with (3.20) then yields
$$\frac{\eta }{4}𝒢_{\phi ,\beta }(\alpha )=𝒥,$$
(3.21)
a result apparently similar to Fronsdal’s . However, now the full Lagrangian system implies that, when the compensator $`\alpha `$ satisfies its field equation, the coupling can only be consistent if the current $`𝒥`$ is *divergence free*. Incidentally, this is just the expected Noether constraint for a source related to the gauge symmetry of the theory described by (3.6).
## 4 Local Lagrangians for unconstrained fermions
One can repeat the previous steps almost verbatim for fermion fields. Here the starting point is provided by the Fang-Fronsdal equation for a symmetric spin - $`(n+1/2)`$ spinor-tensor $`\psi `$,
$$𝒮i\left(\overline{)}\psi \overline{)}\psi \right)=0,$$
(4.1)
that is invariant under the gauge transformation $`\delta \psi =ϵ`$ only if the gauge parameter is *$`\gamma `$ \- traceless*, since
$$\delta 𝒮=\mathrm{\hspace{0.17em}2}i^{\mathrm{\hspace{0.17em}2}}\overline{)}ϵ.$$
(4.2)
In a similar fashion, the corresponding Lagrangian
$$=\frac{1}{2}\overline{\psi }\left(𝒮\frac{1}{2}\gamma \overline{)}𝒮\frac{1}{2}\eta 𝒮^{}\right)+h.c.$$
(4.3)
is gauge invariant only if the gauge field $`\psi `$ is *triply $`\gamma `$ \- traceless*, on account of the “anomalous” Bianchi identity
$$𝒮\frac{1}{2}𝒮^{}\frac{1}{2}\overline{)}\overline{)}𝒮=i^{\mathrm{\hspace{0.33em}2}}\overline{)}\psi ^{}.$$
(4.4)
In this case one begins by considering the Lagrangian (4.3), written however for an *unconstrained* Fermi field $`\psi `$ of spin $`n+1/2`$, and varied with an *unconstrained* gauge parameter $`ϵ`$. The resulting variation is then
$$\begin{array}{cc}\hfill \delta _0=& \frac{3i}{2}\left(\genfrac{}{}{0pt}{}{n}{3}\right)\overline{ϵ}\overline{)}\psi ^{}+\frac{3i}{4}\left(\genfrac{}{}{0pt}{}{n}{3}\right)\overline{ϵ}^{}\overline{)}\psi ^{}3i\left(\genfrac{}{}{0pt}{}{n}{3}\right)\overline{ϵ}^{}\overline{)}\psi \hfill \\ & \frac{3i}{4}\left(\genfrac{}{}{0pt}{}{n}{3}\right)\overline{ϵ}^{}\mathrm{}\overline{)}\psi ^{}+3i\left(\genfrac{}{}{0pt}{}{n}{4}\right)\overline{ϵ}^{}\overline{)}\psi ^{}+2i\left(\genfrac{}{}{0pt}{}{n}{2}\right)\overline{\overline{)}ϵ}\psi \hfill \\ & 2i\left(\genfrac{}{}{0pt}{}{n}{2}\right)\overline{\overline{)}ϵ}\overline{)}\overline{)}\psi i\left(\genfrac{}{}{0pt}{}{n}{2}\right)\overline{\overline{)}ϵ}\mathrm{}\psi ^{}+\frac{9i}{2}\left(\genfrac{}{}{0pt}{}{n}{3}\right)\overline{\overline{)}ϵ}\psi ^{}\hfill \\ & 3i\left(\genfrac{}{}{0pt}{}{n}{4}\right)\overline{\overline{)}ϵ}^{}\psi ^{}+h.c..\hfill \end{array}$$
(4.5)
In complete analogy with the bosonic case, all terms involving the $`\gamma `$ \- trace $`\overline{)}ϵ`$ of the gauge parameter can be canceled by additional terms linear in the compensator field $`\xi `$, that are collected in
$$\begin{array}{cc}\hfill _1=& \frac{3i}{4}\left(\genfrac{}{}{0pt}{}{n}{3}\right)\overline{\overline{)}\xi }\overline{)}\psi ^{}+\mathrm{\hspace{0.17em}3}i\left(\genfrac{}{}{0pt}{}{n}{3}\right)\overline{\overline{)}\xi }\overline{)}\psi +\frac{3i}{4}\left(\genfrac{}{}{0pt}{}{n}{3}\right)\overline{\overline{)}\xi }\mathrm{}\overline{)}\psi ^{}\hfill \\ & \mathrm{\hspace{0.17em}3}i\left(\genfrac{}{}{0pt}{}{n}{4}\right)\overline{\overline{)}\xi }\overline{)}\psi ^{}\mathrm{\hspace{0.17em}2}i\left(\genfrac{}{}{0pt}{}{n}{2}\right)\overline{\xi }\psi +\mathrm{\hspace{0.17em}2}i\left(\genfrac{}{}{0pt}{}{n}{2}\right)\overline{\xi }\overline{)}\overline{)}\psi \hfill \\ & +i\left(\genfrac{}{}{0pt}{}{n}{2}\right)\overline{\xi }\mathrm{}\psi ^{}\frac{9i}{2}\left(\genfrac{}{}{0pt}{}{n}{3}\right)\overline{\xi }\psi ^{}+\mathrm{\hspace{0.17em}3}i\left(\genfrac{}{}{0pt}{}{n}{4}\right)\overline{\xi }^{}\psi ^{}+h.c..\hfill \end{array}$$
(4.6)
Additional terms depending on $`\overline{)}ϵ`$ generated by the variation of $`_1`$ can then be eliminated adding
$$\begin{array}{cc}\hfill _2=& \frac{15i}{2}\left(\genfrac{}{}{0pt}{}{n}{3}\right)\overline{\overline{)}\xi }\mathrm{}\xi i\left(\genfrac{}{}{0pt}{}{n}{2}\right)\overline{\xi }\mathrm{}\overline{)}\xi +\mathrm{\hspace{0.17em}3}i\left(\genfrac{}{}{0pt}{}{n}{3}\right)\overline{\xi }\overline{)}\xi \hfill \\ & +18i\left(\genfrac{}{}{0pt}{}{n}{4}\right)\overline{\overline{)}\xi }\xi +\mathrm{\hspace{0.17em}6}i\left(\genfrac{}{}{0pt}{}{n}{4}\right)\overline{\overline{)}\xi }\mathrm{}\xi ^{}\hfill \\ & \mathrm{\hspace{0.17em}15}i\left(\genfrac{}{}{0pt}{}{n}{5}\right)\overline{\overline{)}\xi }\xi ^{}+h.c.,\hfill \end{array}$$
(4.7)
and the total variation is finally
$$\delta \left\{_0+_1+_2\right\}=\frac{3i}{2}\left(\genfrac{}{}{0pt}{}{n}{3}\right)\overline{ϵ}(\overline{)}\psi ^{}\mathrm{\hspace{0.17em}2}\xi \overline{)}\overline{)}\xi \xi ^{})+h.c..$$
(4.8)
This residual contribution is proportional to the constraint relating $`\overline{)}\psi ^{}`$ to the compensator in eqs. (2.20). One can finally introduce a Lagrange multiplier field $`\lambda `$, a spinor-tensor of spin $`n\mathrm{\hspace{0.17em}5}/2`$ such that $`\delta \lambda =ϵ`$, to compensate (4.8) by the additional term
$$_3=\frac{3i}{2}\left(\genfrac{}{}{0pt}{}{n}{3}\right)\overline{\lambda }\left(\overline{)}\psi ^{}\mathrm{\hspace{0.17em}2}\xi \overline{)}\overline{)}\xi \xi ^{}\right)+h.c..$$
(4.9)
Summarizing, the complete Lagrangian for an *unconstrained* spin - $`(n+1/2)`$ spinor-tensor is
$$\begin{array}{cc}\hfill =& \frac{1}{2}\overline{\psi }\left(𝒮\frac{1}{2}\gamma \overline{)}𝒮\frac{1}{2}\eta 𝒮^{}\right)\hfill \\ & \frac{3i}{4}\left(\genfrac{}{}{0pt}{}{n}{3}\right)\overline{\overline{)}\xi }\overline{)}\psi ^{}+\mathrm{\hspace{0.17em}3}i\left(\genfrac{}{}{0pt}{}{n}{3}\right)\overline{\overline{)}\xi }\overline{)}\psi +\frac{3i}{4}\left(\genfrac{}{}{0pt}{}{n}{3}\right)\overline{\overline{)}\xi }\mathrm{}\overline{)}\psi ^{}\hfill \\ & \mathrm{\hspace{0.17em}3}i\left(\genfrac{}{}{0pt}{}{n}{4}\right)\overline{\overline{)}\xi }\overline{)}\psi ^{}\mathrm{\hspace{0.17em}2}i\left(\genfrac{}{}{0pt}{}{n}{2}\right)\overline{\xi }\psi +2i\left(\genfrac{}{}{0pt}{}{n}{2}\right)\overline{\xi }\overline{)}\overline{)}\psi \hfill \\ & +i\left(\genfrac{}{}{0pt}{}{n}{2}\right)\overline{\xi }\mathrm{}\psi ^{}\frac{9i}{2}\left(\genfrac{}{}{0pt}{}{n}{3}\right)\overline{\xi }\psi ^{}+\mathrm{\hspace{0.17em}3}i\left(\genfrac{}{}{0pt}{}{n}{4}\right)\overline{\xi }^{}\psi ^{}\hfill \\ & \frac{15i}{2}\left(\genfrac{}{}{0pt}{}{n}{3}\right)\overline{\overline{)}\xi }\mathrm{}\xi i\left(\genfrac{}{}{0pt}{}{n}{2}\right)\overline{\xi }\mathrm{}\overline{)}\xi +\mathrm{\hspace{0.17em}3}i\left(\genfrac{}{}{0pt}{}{n}{3}\right)\overline{\xi }\overline{)}\xi \hfill \\ & +\mathrm{\hspace{0.17em}18}i\left(\genfrac{}{}{0pt}{}{n}{4}\right)\overline{\overline{)}\xi }\xi +\mathrm{\hspace{0.17em}6}i\left(\genfrac{}{}{0pt}{}{n}{4}\right)\overline{\overline{)}\xi }\mathrm{}\xi ^{}\mathrm{\hspace{0.17em}15}i\left(\genfrac{}{}{0pt}{}{n}{5}\right)\overline{\overline{)}\xi }\xi ^{}\hfill \\ & +\frac{3i}{2}\left(\genfrac{}{}{0pt}{}{n}{3}\right)\overline{\lambda }\left(\overline{)}\psi ^{}\mathrm{\hspace{0.17em}2}\xi \overline{)}\overline{)}\xi \xi ^{}\right)+h.c.,\hfill \end{array}$$
(4.10)
and is invariant under the gauge transformations
$$\begin{array}{cc}& \delta \psi =ϵ,\hfill \\ & \delta \xi =\overline{)}ϵ,\hfill \\ & \delta \lambda =ϵ.\hfill \end{array}$$
(4.11)
The corresponding field equations are:
$`\overline{\psi }`$ $`:i𝒮+2^2\xi {\displaystyle \frac{1}{2}}\eta (i𝒮^{}+{\displaystyle \frac{1}{2}}\overline{)}\psi ^{}{\displaystyle \frac{1}{2}}\overline{)}\overline{)}\xi +\mathrm{\hspace{0.17em}2}\mathrm{}\xi +\mathrm{\hspace{0.17em}3}\xi +^2\xi ^{})`$
$`{\displaystyle \frac{1}{2}}\gamma (i\overline{)}𝒮+\mathrm{\hspace{0.17em}2}^2\overline{)}\xi +\mathrm{\hspace{0.17em}2}\overline{)}\xi )+{\displaystyle \frac{1}{4}}\gamma \eta (\psi ^{}\mathrm{}\overline{)}\xi \overline{)}\xi \mathrm{\hspace{0.17em}2}\lambda )=\mathrm{\hspace{0.17em}0},`$ (4.12)
$`\overline{\lambda }`$ $`:\overline{)}\psi ^{}\mathrm{\hspace{0.17em}2}\xi \overline{)}\overline{)}\xi \xi ^{}=\mathrm{\hspace{0.17em}0},`$ (4.13)
$`\overline{\xi }`$ $`:\mathrm{}\psi ^{}+\mathrm{\hspace{0.17em}2}\overline{)}\overline{)}\psi \mathrm{\hspace{0.17em}2}\psi +{\displaystyle \frac{3}{2}}\psi ^{}\mathrm{\hspace{0.17em}2}\mathrm{}\overline{)}\xi {\displaystyle \frac{5}{2}}\mathrm{}\overline{)}\xi \mathrm{\hspace{0.17em}2}\overline{)}\xi `$
$`\mathrm{\hspace{0.17em}3}^2\overline{)}\xi \lambda +\eta ({\displaystyle \frac{1}{2}}\psi ^{}\mathrm{}\overline{)}\xi {\displaystyle \frac{1}{2}}\overline{)}\xi \lambda )`$
$`+\gamma ({\displaystyle \frac{1}{4}}\overline{)}\psi ^{}+\overline{)}\psi +{\displaystyle \frac{1}{4}}\mathrm{}\overline{)}\psi ^{}+{\displaystyle \frac{1}{4}}\overline{)}\psi ^{}`$
$`{\displaystyle \frac{5}{2}}\mathrm{}\xi {\displaystyle \frac{3}{2}}\xi {\displaystyle \frac{1}{2}}\mathrm{}\xi ^{}{\displaystyle \frac{1}{2}}^2\xi ^{}{\displaystyle \frac{1}{2}}\overline{)}\lambda )=\mathrm{\hspace{0.17em}0}.`$ (4.14)
As in the bosonic case, we can now relate them to the simple non-Lagrangian system (2.20). The basic observation is to recognize that, when (4.13) is satisfied, the equation for $`\psi `$ takes the form
$$𝒲\frac{1}{2}\gamma \overline{)}𝒲\frac{1}{2}\eta 𝒲^{}+\frac{i}{4}\eta \gamma 𝒵=\mathrm{\hspace{0.17em}0},$$
(4.15)
where
$$\begin{array}{cc}\hfill 𝒲& =𝒮+\mathrm{\hspace{0.17em}2}i^2\xi ,\hfill \\ \hfill 𝒵& =\psi ^{}\mathrm{}\overline{)}\xi \overline{)}\xi \mathrm{\hspace{0.17em}2}\lambda .\hfill \end{array}$$
(4.16)
Moreover, under the same conditions the triple $`\gamma `$ \- trace of $`𝒲`$ vanishes. It is then possible to rephrase the iterative argument presented for bosonic fields: if $`p`$ is the integer part of $`\frac{s3}{2}`$, where $`s3`$, taking $`(p+1)`$ successive traces of (4.15) one arrives at the condition that the highest $`\gamma `$ \- trace vanish:
$$\gamma \left(\frac{1}{4}\eta \gamma 𝒵\right)^{[p+1]}=\mathrm{\hspace{0.17em}0}.$$
(4.17)
Inserting (4.17) in the lower $`\gamma `$ \- traces of (4.15), one can recursively show that *all* $`\gamma `$ \- traces of $`𝒵`$ vanish as well. Consequently, the first trace and the $`\gamma `$ \- trace of (4.15) imply that $`𝒲^{}=\mathrm{\hspace{0.17em}0}`$ and $`\overline{)}𝒲=\mathrm{\hspace{0.17em}0}`$, and in conclusion the Lagrangian (4.10) leads indeed to the local compensator equations (2.20). In complete analogy with the bosonic case, making use of (4.13) one can show that the field equation (4.14) for the compensator $`\xi `$ is proportional to the divergence of (4.12).
## Acknowledgments
We are grateful to the CPhT of the Ecole Polytechnique, to the LPT-Orsay and to the CERN Theory Division for the kind hospitality extended to us, and to C. Iazeolla, J. Mourad and M. Porrati for stimulating discussions. The visits of A.S. were partly supported by a “poste rouge” CNRS, by the CNRS PICS n. 3059 and by the CERN TH unit, while the visits of D.F. were partly supported by the CNRS PICS n. 3059 and by the APC group of U. Paris VII. The present work was also supported in part by INFN, by the MIUR-COFIN contract 2003-023852, by the EU contracts MRTN-CT-2004-503369 and MRTN-CT-2004-512194, by the INTAS contract 03-51-6346, and by the NATO grant PST.CLG.978785. |
warning/0507/astro-ph0507647.html | ar5iv | text | # The DEEP2 Galaxy Redshift Survey: Clustering of Groups and Group Galaxies at 𝑧∼1
## 1. Introduction
Groups of galaxies populate an intermediate range in density-contrast between galaxies and clusters and occupy a regime that is critical to understanding hierarchical galaxy formation in $`\mathrm{\Lambda }`$CDM models. Merger events between galaxies likely occur within groups rather than clusters due to their lower velocity dispersions (e.g., Ostriker, 1980; Barnes, 1985). Galaxy groups should also be more easily related to dark matter halos than galaxies themselves, which can have a complicated halo-occupation function that depends significantly on galaxy properties. In order to understand and test galaxy formation and evolution models it is useful to relate galaxies to observable groups as a proxy for their parent dark matter halos. The current halo model paradigm (e.g., Seljak, 2000; Ma & Fry, 2000; Peacock & Smith, 2000; Cooray & Sheth, 2002; Kravtsov et al., 2004) provides a statistical analytic measure for relating galaxies to their dark matter halos. A key statistic in the halo model is the halo occupation distribution (HOD) which measures the probability of a halo of a given mass hosting $`N`$ galaxies. The halo model also naturally explains the small deviations seen in the clustering of galaxies at $`z0`$ from a power-law model, where there is a transition from galaxies within a single halo and between different halos. The clustering of groups and group galaxies depends not only on halo model parameters but also on the nature of bias and the details of hierarchical structure formation, as well as cosmological parameters such as $`\mathrm{\Omega }_m`$ and $`\sigma _8`$.
Observationally, clusters of galaxies have been shown to be very strongly clustered (Bahcall, 1988), with the clustering strength depending on the richness of the cluster. Kaiser (1984) show that the large clustering scale-length of massive and rare Abell clusters can be explained by a simple model in which these clusters formed in regions where the primordial density enhancement was unusually high. Objects forming in the densest peaks would naturally be biased tracers of the underlying dark matter field such that massive clusters would have a higher correlation amplitude than that of galaxies. Since their masses are intermediate, galaxy groups are therefore expected to have clustering properties between those of galaxies and clusters.
The first papers analyzing the clustering of groups in local redshift surveys at $`z0`$ present conflicting results and were hampered by small samples and cosmic variance (e.g., Jing & Zhang, 1988; Maia & da Costa, 1990; Ramella, Geller, & Huchra, 1990). Trasarti-Battistoni, Invernizzi, & Bonometto (1997) investigated the effect of changing the linking-length parameters in the Friends-of-Friends (FoF) algorithm and its effect on the clustering signal and found that these early papers used too large a linking length, which led to a diminished clustering strength due to the presence of interlopers. Trasarti-Battistoni, Invernizzi, & Bonometto (1997) found in their data from the Perseus-Pisces redshift survey, using two fields and $`50`$ and 200 groups in each field, that groups are approximately twice as clustered as galaxies, but with significant error bars. All of these papers showed that determining the clustering properties of groups is a tricky endeavor, which can depend quite sensitively on the volume and magnitude depth of the survey, the handling of the selection function and varying completeness, and the method used to identify groups. In addition, these analyses all suffered from significant uncertainties, from both Poisson statistics, due to the small number of groups in the surveys, and cosmic variance, due to the small volumes surveyed, which was not quantified in any of these papers.
Significant advances have recently been made with much larger datasets. Girardi, Boschin, & da Costa (2000), with the combined CfA2 and SSRS2 surveys, have a sample of 885 groups in a volume of $`3\times 10^5h^3\mathrm{Mpc}^3`$, much larger than earlier surveys. With this large sample size, they are able to construct volume-limited subsamples and investigate the dependence of clustering on group properties, finding that groups with more members and/or larger internal velocity dispersion are more strongly clustered. Merchán, Maia, & Lambas (2000) use a sample of 517 groups from the Updated Zwicky Catalog and 104 groups from the SSRS2 to show that groups are at least twice as clustered as galaxies, and that more massive groups have larger clustering strength.
Most recently, the 2dF Galaxy Redshift Survey has provided a vast dataset with which to study large-scale structure at $`z0.2`$. Using data from the 100k release, Zandivarez, Merchán, & Padilla (2003) measure the clustering of groups as a function of virial mass and find that more massive groups are more clustered and that their measurements match the clustering of dark matter halos in a $`\mathrm{\Lambda }`$CDM N-body simulation. Using the completed 2dF survey, Padilla et al. (2004) analyze group clustering as a function of luminosity and show that while the least luminous groups actually cluster less than the galaxies in the survey, there is a strong relation between group luminosity and correlation length, and the most luminous groups (with $`L4\times 10^{11}h^2L_{\mathrm{}}`$) are $`10`$ times more clustered as the least luminous groups (with $`L2\times 10^{10}h^2L_{\mathrm{}}`$). The relation between clustering scale length and mean group separation that they find continues the trend seen on larger scales for clusters. They find very good agreement between their data and mock catalogs constructed from $`\mathrm{\Lambda }`$CDM simulations and semi-analytic galaxy evolution recipes. These same conclusions are reached by Yang et al. (2005c), who also use the 2dF data to measure the clustering of groups as a function of luminosity. It appears that at $`z0`$ there is now convergence among group clustering analyses.
The extensive 2dF group catalogs have now also allowed studies of groups beyond simple measures of the correlation function of groups. Several authors measure the radial profile of galaxies in groups using 2dF data (Collister & Lahav, 2005; Diaz et al., 2005; Yang et al., 2005b), and analyze the cross-correlation between group centers and galaxies (Yang et al., 2005b). These papers find that galaxies in groups are less concentrated than dark matter particles in simulations, and also find that, locally, the centers of groups are preferentially populated by red galaxies compared to blue galaxies. The HOD has now also been measured directly at $`z0`$ using counts of galaxies in groups of different masses (Collister & Lahav, 2005; Yang et al., 2005a), constraining halo model parameters locally. Clustering measures of the correlation function of all galaxies (not just those in groups) at intermediate redshifts indicate that the HOD does not change significantly between $`z1`$ and $`z0`$ (Yan, Madgwick, & White, 2003; Phleps et al., 2005).
These measurements have only been performed with local samples; group catalogs have not been available at intermediate- or high-redshift. In this paper we present the first analysis of group clustering at $`z1`$, using group catalogs from the DEEP2 Galaxy Redshift Survey (Gerke et al., 2005). The high resolution of the DEEP2 data allows us to identify groups in three dimensions regardless of their galaxy properties, using only their overdensity in space. We focus here on the clustering of groups and galaxies within groups at $`z1`$, as these measures, when combined with similar measures at $`z0`$, will provide constraints on galaxy evolution and structure formation models. We also show how these measures can constrain the HOD at $`z1`$.
Four measures of two-point clustering in the DEEP2 dataset between $`0.7z1.0`$ are analyzed: 1) the group correlation function, 2) the galaxy correlation function for DEEP2 galaxies, 3) the galaxy correlation function for galaxies in groups, and 4) the group-galaxy cross-correlation function. These clustering measures can be used to constrain the halo model parameters by comparing the data to mock catalogs with different HODs. The first clustering measure, when combined with the observed number density of groups in our sample, is used to estimate the typical dark matter masses of the halos the groups studied reside in. The second measure provides constraints on the HOD, though not in an entirely unique way; further constraints on the HOD are provided by the last two measures. For the clustering of galaxies in groups, we empirically distinguish between ‘one-halo’ and ‘two-halo’ terms using pairs of galaxies within the same group and pairs in different groups, respectively. The clustering of galaxies in groups and the group-galaxy cross-correlation function can also constrain the radial distribution of galaxies within groups when compared with mock catalogs.
Lastly, groups may be used to constrain cosmological parameters like the dark energy equation of state, $`w`$, if their bivariate distribution in redshift $`z`$ and velocity dispersion $`\sigma `$ can be accurately measured (Newman et al., 2002). This test, however, requires that the relation between group velocity dispersion and dark matter halo mass be known and accurately calibrated. It has recently been suggested (Majumdar & Mohr, 2004; Lima & Hu, 2004) that the clustering properties of galaxy clusters may be used for ”self-calibration”, since the clustering properties of halos can be predicted as a function of mass and compared to the measured clustering. The group correlation function results presented here should be useful as such a self-calibration procedure for future DEEP2 studies.
An outline of the paper is as follows: §2 briefly describes the DEEP2 Galaxy Redshift Survey and the sample of galaxy groups used here, as well as the mock galaxy catalogs constructed for the survey. §3 discusses the methods used to calculate the two-point correlation functions. We present clustering results for groups in §4, where we compare with simulations and estimate the minimum dark matter halo mass for our groups. In §5 we analyze the clustering of the full galaxy sample and galaxies in groups in the DEEP2 data and in mock catalogs and show the contribution to the correlation function for galaxies in groups from the ‘one-halo’ and ‘two-halo’ terms. §6 presents the cross-correlation between the full galaxy sample and group centers, which depends upon the radial profile of galaxies in groups. The relative biases between galaxies in groups and all galaxies and between groups and galaxies are presented in §7. Mock catalogs with different halo models are used to illustrate how these various clustering measures can be used to constrain parameters of the halo model in §8. We conclude in §9.
## 2. Data Sample and Mock Catalogs
In this paper we use data from the DEEP2 Galaxy Redshift Survey, which is an ongoing project using the DEIMOS spectrograph (Faber et al., 2002) on the 10-m Keck II telescope to survey optical galaxies at $`z1`$ in a comoving volume of approximately 5$`\times `$10<sup>6</sup> $`h^3`$ Mpc<sup>3</sup>. Using $`1`$ hr exposure times, the survey will measure redshifts for $`40,000`$ galaxies in the redshift range $`0.7z1.5`$ to a limiting magnitude of $`R_{\mathrm{AB}}=24.1`$ (Coil et al., 2004a; Faber et al., 2005). Spectroscopic targets are pre-selected using a color cut in $`BR`$ \- $`RI`$ space to ensure that most galaxies lie beyond $`z0.75`$. This color-cut results in a sample with $``$90% of the targeted objects at $`z>0.75`$, missing only $``$3% of the $`z>0.75`$ galaxies which meet our magnitude limit (Davis et al., 2002). Due to the high dispersion ($`R5,000`$) of our spectra, our redshift errors, determined from repeated observations, are $`30`$km s<sup>-1</sup>. Restframe $`(UB)_0`$ colors have been derived as described in (Willmer et al., 2005). Details of the observations, catalog construction and data reduction can be found in Davis et al. (2002); Coil et al. (2004b); Davis, Gerke, & Newman (2004); Faber et al. (2005).
The completed survey will cover 3 square degrees of the sky over four widely separated fields to limit the impact of cosmic variance. Each field is comprised of two to four contiguous photometric ’pointings’ of size $`0.5`$ by $`0.67`$ degrees. Here we use data from six of our most complete pointings to date, in three separate fields. We use data from pointings 1 and 2 in the DEEP2 fields 2, 3 and 4. Each DEEP2 pointing corresponds to a volume of comoving dimensions $`20\times 27\times 550`$ $`h^1`$ Mpc in a $`\mathrm{\Lambda }`$CDM model for $`0.7z1.0`$. The total volume of the sample used in this paper is thus $`1.8\times 10^6h^3\mathrm{Mpc}^3`$. To convert measured redshifts to comoving distances along the line of sight we assume a flat $`\mathrm{\Lambda }`$CDM cosmology with $`\mathrm{\Omega }_\mathrm{m}=0.3`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$. We define $`hH_0/(100\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1)`$ and quote correlation lengths, $`r_0`$, in comoving $`h^1`$ Mpc. Throughout the paper, we quote empirical errors calculated from the variance across the six pointings. In so doing we have treated the six DEEP2 pointings as being entirely independent, and they are not: there are two adjacent pointings in each of three independent fields. We estimate from Monte Carlo simulations using the methods outlined in Newman et al. (2002) that, for the distribution of pointings used here, the measured standard error should be increased by $`17\pm 5`$%, for errors which are dominated by cosmic variance; this is a conservative assumption, as (especially for the group-group correlations) other contributions such as shot noise are significant.
A description of the methods employed to detect groups is presented in Gerke et al. (2005), along with details of the group catalog. A Voronoi-Delaunay Method (VDM) group-finder (Marinoni et al., 2002) is used to identify galaxy groups. This method searches for galaxy overdensities in redshift space, using an asymmetric search window to account for redshift-space distortions. The advantage of the method over traditional FoF group-finding methods lies in the fact that it has no fixed length scale, but instead uses an adaptive search radius based on estimated group richness. The VDM group-finder thus avoids a common problem of FoF methods, in which clustered groups are merged together along filaments.
We use three galaxy samples in this paper; all are drawn from the same volume as the group sample. The full galaxy sample includes a total of 9787 galaxies in the redshift range $`0.7z1.0`$ in the same six pointings used for the group analysis. We also split the full galaxy sample into subsamples of field and group galaxies, with 5947 and 3840 objects in each. Throughout this paper we use the terms “all galaxies” and “the full galaxy sample” interchangeably; they include both the field and group galaxy samples.
Our main group sample consists of all DEEP2 groups with an estimated velocity dispersion $`\sigma 200`$ km s<sup>-1</sup>, with a total of 460 groups and $`50100`$ groups per pointing. We do not make a distinction between groups and clusters - clusters are simply the larger groups; they are included in all analyses, but have minimal impact on this work due to their rarity. The $`\sigma 200`$ km s<sup>-1</sup> group sample has similar completeness ($`74\pm 5`$%) and purity ($`57\pm 3`$%) as samples with higher $`\sigma `$ cutoffs, as estimated using mock catalogs described below. Here completeness is defined as the fraction of real groups (defined as a group of galaxies belonging to the same parent halo, where halos are identified using FoF in real space on the dark matter particles in the simulations) that are successfully identified in the recovered group catalog created by the group-finder, and purity is the fraction of recovered groups that correspond to real groups (see Gerke et al. (2005) for more details). As shown in Fig. 8 of Gerke et al. (2005), these statistics are nearly independent of group velocity dispersion. There are more recovered groups than real groups in these mock catalogs by a factor of 1.4. The $`\sigma 200`$ km s<sup>-1</sup> group sample has a relatively high “galaxy success rate” of $`70\pm 1`$%, defined as the fraction of galaxies in real groups that are identified as group galaxies in our recovered group catalog, and an interloper fraction of $`43\pm 1`$%. We do not use the $`\sigma 350`$ km s<sup>-1</sup> cut that was used in Gerke et al. (2005), as that cut was used to define a sample that recovered different physical properties than the ones relevant here; instead of the distribution of groups in redshift and velocity dispersion, the relevant parameters here are the positions of groups and the interloper fraction (the fraction of identified group galaxies that are not actually in groups).
The observed richness distribution of our groups in the DEEP2 data is roughly a power-law, with most groups having two observed galaxies; the largest groups have $`10`$ galaxies, though there are only a few groups this large. Fig. 1 shows the observed richness distribution for groups in the DEEP2 data as well as for real and recovered groups in our mock catalogs, for groups with $`0.7z1.0`$ and $`\sigma 200`$ km s<sup>-1</sup>.
A histogram of the redshift distribution of our group sample is shown in Fig. 2. We restrict the analyses here to the redshift range $`0.7z1.0`$ to minimize systematic effects. While groups are found in the survey to higher redshifts, those groups are likely to be more massive and hence not as representative of groups in our sample as a whole. Additionally, the $`R`$-band target selection of the survey corresponds to a bluer restframe color-selection at higher redshift; this results in fewer red galaxies being targeted at higher redshift compared to blue galaxies (Willmer et al., 2005), which could systematically affect the group richness and velocity dispersion estimates at $`z1.0`$. The spatial distribution of DEEP2 groups in three of the six pointings used is shown in Fig. 3. Shown are pointings 2 in the DEEP2 fields 2, 3 and 4.
The mock galaxy catalogs used throughout the paper are described in Yan, White, & Coil (2004); relevant details are repeated here. The mock catalogs are constructed from N-body simulations of $`512^3`$ dark matter particles with a particle mass of $`m=1\times 10^{10}h^1`$ $`M_{\mathrm{}}`$ in a box with dimensions 256 $`h^1`$ Mpc on a side, for a $`\mathrm{\Lambda }`$CDM cosmology with $`\sigma _8=0.9`$. Dark matter halos were identified using FoF in real space and galaxies were placed within halos using a halo model prescription. To populate dark matter halos with galaxies, two functions need to be specified. The first is the halo occupation distribution function (HOD), which is the probability that a halo of mass M hosts N galaxies, $`P(N|M)`$. The second is the spatial distribution of galaxies within halos. The first moment of the HOD function, the average number of galaxies as a function of halo mass M, is shown in Fig. 1 in Yan, White, & Coil (2004) and in Fig. 12 of this paper (labeled as “B256”) for the HOD used in the mock galaxy catalogs here. This HOD is constrained at $`z0`$ with the 2dF luminosity function and luminosity dependent two-point correlation function of galaxies and at $`z1`$ with the DEEP2 $`\xi (r)`$ and the COMBO-17 luminosity functions. In §8 of this paper we compare clustering results for mock catalogs with two different halo models; for the bulk of the paper, however, we use lightcones with the HOD as published in Yan, White, & Coil (2004). Once the number of galaxies in each halo and their corresponding luminosities are known, the most luminous galaxy is assigned to the center of mass of the halo, and positions and velocities for the other galaxies are drawn randomly from those of the dark matter particles. No other radial or velocity bias is included in assigning galaxies to particles, so that galaxies trace the mass and velocity distributions of the dark matter particles in the halo. The spatial distribution of galaxies follow the dark matter density profile, which on average is an NFW profile. Although no velocity bias is included, the velocity dispersions of the galaxies are systematically smaller than those of the mass (see Fig.8 in Yan, White & Coil 2004), due to the fact that the most luminous galaxy is always assigned to the particle at the center of mass. Here we use a set of twelve independent catalogs which each have the same spatial extent as a single DEEP2 pointing ($`0.5`$ by $`0.67`$ degrees).
The center of each DEEP2 group (needed for the clustering measures performed here) is measured as the median of the positions (in comoving x, y, and z) of the galaxies identified in that group. Errors on the positions of the recovered groups are estimated using the differences in the mock catalogs between the centers of real groups and recovered groups. As stated above, there is a factor of 1.4 more recovered groups in the mock catalogs than real groups for $`\sigma 200`$ km s<sup>-1</sup>. In the mock catalogs 74% of recovered groups have a real group within $`r_p=1`$ $`h^1`$ Mpc, while 48% have a real group within $`r_p=0.2`$ $`h^1`$ Mpc and 41% have a real group within $`r_p=0.1`$ $`h^1`$ Mpc. Here $`r_p`$ is the projected distance on the plane of the sky, which is the relevant distance for the projected clustering used in §5 and later in the paper; in §4 we use the redshift-space correlation function, but only measure it on scales $`s>2`$ $`h^1`$ Mpc. We therefore believe that on scales larger than 1 $`h^1`$ Mpc our results are robust to errors in the group positions, while on smaller scales, where we will measure the cross-correlation between group centers and galaxies, there is likely to be some degradation in the signal due to position errors; this is discussed more in §9. However, the mock catalogs have been treated in an identical manner as the data so that comparisons between the data and mock catalogs are unaffected.
## 3. Methods
The two-point correlation function $`\xi (r)`$ is defined as a measure of the excess probability above Poisson of finding an object in a volume element $`dV`$ at a separation $`r`$ from another randomly chosen object,
$$dP=n[1+\xi (r)]dV,$$
(1)
where $`n`$ is the mean number density of the object in question (Peebles, 1980).
Measuring $`\xi (r)`$ requires constructing a random catalog with the same selection criteria and observational effects as the data, to serve as an unclustered distribution with which to compare. For each data sample we create a random catalog with the same overall sky coverage and redshift distribution as the data. This is achieved by first applying the two-dimensional window function of our data in the plane of the sky to the random catalog. Our overall redshift success rate is $`gtrsim`$70% and is not entirely uniform across the survey; some slitmasks are observed under better conditions than others and therefore yield a slightly higher completeness. This spatially-varying redshift success rate is taken into account in the spatial window function which is applied to both the random catalog and the mock catalogs, such that regions of the sky with a higher completeness have a correspondingly higher number of random points or more objects in the mock galaxy catalogs. This ensures that there is no bias introduced when computing correlation statistics or comparing the data to the mock catalogs. We also mask the regions of the random and mock catalogs where the photometric data have saturated stars and CCD defects.
We then apply a selection function, $`\varphi (z)`$, defined as the probability of observing a group as a function of redshift, so that the random catalog has the same overall redshift distribution as the data. The selection function for groups in the DEEP2 survey is determined by smoothing the observed redshift distribution of groups in the catalog used here, as shown in Fig. 2. The apparent overdensity at $`z0.85`$ seen in this figure is due to structures in several pointings and does not significantly affect our results, which are averaged over all pointings and a wide redshift range. As the group catalog includes three separate fields, the data are combined when calculating $`\varphi (z)`$, which reduces the effects of cosmic variance. Using a smoothed redshift distribution to estimate the selection function can cause a systematic bias, but we estimate that this bias is less than the errors on $`\xi (r)`$. In what follows, for both the data and the mock catalogs the redshift range over which we compute $`\xi (r)`$ is limited to $`0.7z1.0`$, the redshift range over which the selection function varies by less than a factor of two. Each of the six pointings has between $`50100`$ groups, and each of the twelve independent mock catalogs contains $`100`$ groups. We use $`3,000`$ random points in each pointing to calculate the correlation function for groups and $`20,000`$ random points to calculate the correlation function for galaxies.
The two-point correlation function is measured for both groups and galaxies using the Landy & Szalay (1993) estimator,
$$\xi =\frac{1}{RR}\left[DD\left(\frac{n_R}{n_D}\right)^22DR\left(\frac{n_R}{n_D}\right)+RR\right],$$
(2)
where $`DD,DR`$, and $`RR`$ are pair counts in a given separation range in the data-data, data-random, and random-random catalogs, and $`n_D`$ and $`n_R`$ are the mean number densities in the data and random catalogs. This estimator has been shown to perform as well as the Hamilton estimator (Hamilton, 1993) but is preferred as it is relatively insensitive to the size of the random catalog and handles edge corrections well (Kerscher, Szapudi, & Szalay, 2000). As the magnitude limit of our survey results in a non-uniform selection function, a standard $`J_3`$-weighting scheme is applied (Davis & Huchra, 1982), which attempts to weight each volume element equally while minimizing the variance on large scales. As the redshift range is limited to $`0.7z1.0`$ for all analyses in this paper, this effect is not large ($`20`$%).
Measurements of the cross-correlation between two samples (presented in §6) use a symmetrized version of Eqn. 2. Each data sample, with pair counts labeled $`D_1`$ and $`D_2`$, has an associated random catalog, with pair counts $`R_1`$ and $`R_2`$, with the same selection function as the data. After normalizing each data and random catalog by its number density, the cross-correlation is estimated using
$$\xi =\frac{1}{R_1R_2}\left[D_1D_2D_1R_2D_2R_1+R_1R_2\right].$$
(3)
Redshift-space distortions due to peculiar velocities along the line-of-sight will introduce systematic effects to the estimate of $`\xi (r)`$. At small separations, random motions within a virialized overdensity cause an elongation along the line-of-sight (“fingers of God”), while on large scales, coherent infall of galaxies into forming structures causes an apparent contraction of structure along the line-of-sight (the “Kaiser effect”). What is actually measured then is $`\xi (s)`$, where $`s`$ is the redshift-space separation between a pair of galaxies. As we will show in the next section using mock group catalogs, for galaxy groups over the scales relevant here, the systematic effects of redshift-space distortions are of order 20%. Mock catalogs are used to correct for redshift-space distortions and infer $`\xi (r)`$ from $`\xi (s)`$by multiplying the observed $`\xi (s)`$for DEEP2 groups by the ratio of $`\xi (r)`$to $`\xi (s)`$for groups in the mock catalogs.
However, redshift-space distortions are more significant when measuring $`\xi `$ for galaxies, as smaller scales are probed. We are able to measure $`\xi `$ on small scales for the galaxy sample both because of the larger sample size, which allows us to stably measure the clustering on small scales, and also because there is no exclusion radius, unlike with groups, which typically have a radius of $`R1`$ $`h^1`$ Mpc. To uncover the real-space clustering properties of galaxies we measure $`\xi `$ in two dimensions, both perpendicular to and along the line of sight. Following Fisher et al. (1994), $`𝐯_\mathrm{𝟏}`$ and $`𝐯_\mathrm{𝟐}`$ are the redshift positions of a pair of galaxies, $`𝐬`$ is the redshift-space separation ($`𝐯_\mathrm{𝟏}𝐯_\mathrm{𝟐}`$), and $`𝐥=\frac{1}{2}`$($`𝐯_\mathrm{𝟏}+𝐯_\mathrm{𝟐})`$ is the mean distance to the pair. The separation between the two galaxies across ($`r_p`$) and along ($`\pi `$) the line of sight are defined as
$$\pi =\frac{𝐬𝐥}{|𝐥|},$$
(4)
$$r_p=\sqrt{𝐬𝐬\pi ^2}.$$
(5)
In applying the Landy & Szalay (1993) estimator to galaxies, pair counts are computed over a two-dimensional grid of separations to estimate $`\xi (r_p,\pi )`$.
To recover $`\xi (r)`$ $`\xi (r_p,\pi )`$ is projected along the $`r_p`$ axis. As redshift-space distortions affect only the line-of-sight component of $`\xi (r_p,\pi )`$, integrating over the $`\pi `$ direction leads to a statistic $`w_p(r_p)`$, which is independent of redshift-space distortions. Following Davis & Peebles (1983),
$$w_p(r_p)=2_0^{\mathrm{}}𝑑\pi \xi (r_p,\pi )=2_0^{\mathrm{}}𝑑y\xi (r_p^2+y^2)^{1/2},$$
(6)
where $`y`$ is the real-space separation along the line of sight. If $`\xi (r)`$ is modeled as a power-law, $`\xi (r)=(r/r_0)^\gamma `$, then $`r_0`$ and $`\gamma `$ can be readily extracted from the projected correlation function, $`w_p(r_p)`$, using an analytic solution to Equation 6:
$$w_p(r_p)=r_p\left(\frac{r_0}{r_p}\right)^\gamma \frac{\mathrm{\Gamma }(\frac{1}{2})\mathrm{\Gamma }(\frac{\gamma 1}{2})}{\mathrm{\Gamma }(\frac{\gamma }{2})},$$
(7)
where $`\mathrm{\Gamma }`$ is the gamma function. A power-law fit to $`w_p(r_p)`$ will then recover $`r_0`$ and $`\gamma `$ for the real-space correlation function, $`\xi (r)`$.
## 4. Group Clustering Results
This section presents results on the clustering of groups at $`z1`$. The two-point correlation function for groups in the DEEP2 data is measured. The clustering properties of groups in our mock catalogs are analyzed, which allows us to quantify our systematic errors and to correct the observed clustering in the DEEP2 data for redshift-space distortions and the effects of the group-finder. The number density of groups in the DEEP2 data is then used to infer the minimum dark matter halo mass which hosts our groups; we then compare the clustering of halos with this mass to the clustering of our group sample.
### 4.1. Clustering of Groups in DEEP2 Data and Mock Catalogs
The left panel of Fig. 4 shows the observed two-point correlation function in redshift space, $`\xi (s)`$, for our $`\sigma 200`$ km s<sup>-1</sup>group sample, for both the full sample (solid line) with two or more galaxies in each group ($`N2`$) and for a subsample with four or more galaxies in each group ($`N4`$). Also shown is $`\xi (s)`$ for the recovered group sample in the mock catalogs with $`\sigma 200`$ and $`N2`$ (dotted line). We show the standard error across the six pointings; this empirical error therefore includes both cosmic variance and Poisson error. The $`N4`$ data sample shows a stronger clustering amplitude than the $`N2`$ data sample. This is likely due to the $`N4`$ sample containing more massive groups, on average, than the $`N2`$ sample. The $`N4`$ sample is also less likely to have interlopers, which may also increase the observed clustering.
Power-law fits to $`\xi (s)`$ over the range $`s=320`$ $`h^1`$ Mpc are given in Table 1. Below $`s3`$ $`h^1`$ Mpc $`\xi (s)`$ does not continue to rise as a single power-law in our data; this lack of pairs on scales $`s<3`$ $`h^1`$ Mpc is likely due to the finite physical extent of groups. To take into account the covariance between $`s_0`$ and $`\gamma `$, we perform a $`\chi ^2`$ minimization and marginalize over each parameter separately. This procedure leads to errors which are roughly a factor of two larger than if we neglected this covariance.
The mock galaxy catalogs described in §2 are used to quantify systematic errors in our group clustering analysis. Effects that may be introduced by the group-finder are tested by comparing $`\xi (s)`$ in redshift space for both real and recovered groups. The results are shown in Table 1. The recovered groups have a slightly higher clustering amplitude (5%) than the real groups. Redshift-space distortions are quantified by measuring $`\xi (s)`$ in redshift space and $`\xi (r)`$ in real space for real groups. Redshift-space distortions appear to enhance the clustering properties of groups by $`20`$% on scales $`s215`$ $`h^1`$ Mpc and decrease the clustering amplitude on smaller scales. The clustering scale length increases by $`20`$% when measured in redshift space, while there is no change to the slope over the scales used here. The effects of our slitmask target selection algorithm on the clustering of groups are also tested using the mock catalogs. Our target selection code determines which galaxies would be observed on slitmasks; because spectra from neighboring galaxies can not overlap on the CCD, not all galaxies can be selected to be observed. In particular, in overdense regions on the plane of the sky the number of galaxies which can be observed decreases in a non-trivial way. By running our slitmask target selection code on the mock catalogs we can quantify the effect on the measured $`\xi (r)`$.
Under the assumption that corrections for the above effects should be proportional to the clustering strength, these results from the mock catalogs can now be used to correct the observed $`\xi (s)`$ for groups in the DEEP2 data for target selection effects, redshift space distortions, and the group-finder, in order to estimate $`\xi (r)`$ in real space for real groups. We can either apply corrections to the power-law fits themselves or to the data points as a function of scale. If we apply the corrections to $`s_0`$ and $`\gamma `$ as measured for the DEEP2 groups, the corrections would infer that $`r_0=6.8\pm 0.6`$ $`h^1`$ Mpc and $`\gamma =1.4\pm 0.2`$ for the real-space correlation function of real groups in the DEEP2 data with $`N2`$. If we explicitly correct the observed $`\xi (s)`$ for the DEEP2 groups as a function of scale, using the ratio of $`\xi (r)`$/$`\xi (s)`$ in the mock catalogs, the resulting $`\xi (r)`$ has a power-law fit of $`r_0=6.2\pm 0.4`$ $`h^1`$ Mpc and $`\gamma =1.5\pm 0.2`$, within the 1$`\sigma `$ error of the inferred values. The corrected $`\xi (r)`$ is shown as the solid line in the right panel of Fig. 4, for both the $`N2`$ and $`N4`$ samples. The clustering of groups in the DEEP2 data is therefore 1-3$`\sigma `$ lower than the clustering of real groups before target selection in the mock catalogs, where $`r_0=7.4\pm 0.4`$ and $`\gamma =1.6\pm 0.2`$, for $`N2`$.
### 4.2. Minimum Group Mass Derived from Clustering Results
The minimum dark matter halo mass that our groups reside in can be inferred from the observed number density and clustering strength of our groups using either analytic formulations of the dark matter halo mass function or by comparing directly to simulations. Here we estimate a minimum dark matter halo mass from the observed number density using analytic theory and then compare the corresponding predicted clustering strength of those halos in both theory and simulations with the observed clustering strength of our groups.
The $`\sigma 200`$ km s<sup>-1</sup> group sample has an observed density of $`n=4.5\times 10^4h^3\mathrm{Mpc}^3`$, calculated from the observed number of groups between $`z=0.750.85`$, where the group selection function is the highest, divided by the comoving volume occupied by the groups. The mock catalogs are used to estimate the effects on the observed number density due to our slitmask target selection (which decreases the number of observed groups, as most groups have only two or three observed galaxies) and the false interloper rate due to the group-finder. Correcting for these effects, the actual comoving number density is estimated to be $`n=6\times 10^4h^3\mathrm{Mpc}^3`$. This corresponds to a mean inter-group spacing of $`d=11.8`$ $`h^1`$ Mpc. For comparison to the sample used by Gerke et al. (2005), the number density of groups with $`\sigma 350`$ km s<sup>-1</sup> is $`n=2.4\times 10^4h^3\mathrm{Mpc}^3`$, after applying the above corrections, which corresponds to an intergroup spacing of $`d=16`$ $`h^1`$ Mpc.
For a $`\mathrm{\Lambda }`$CDM cosmology with $`\mathrm{\Omega }_b=0.05`$, $`\mathrm{\Lambda }=0.7`$, $`\mathrm{\Omega }_m=0.3`$, $`h=0.7`$, $`\sigma _8=0.9`$, a comoving number density of $`n=6\times 10^4h^3\mathrm{Mpc}^3`$ results in $`M_{min}=5.9\times 10^{12}h^1M_{\mathrm{}}`$ at $`z=0.8`$ and $`M_{min}=5.5\times 10^{12}h^1M_{\mathrm{}}`$ at $`z=1`$ for a Sheth & Tormen (1999) mass function. A comoving density of $`n=2.4\times 10^4h^3\mathrm{Mpc}^3`$, estimated for the $`\sigma 350`$ km s<sup>-1</sup> sample, corresponds to $`M_{min}=1.2\times 10^{13}h^1M_{\mathrm{}}`$ at $`z=0.8`$ and $`M_{min}=1.1\times 10^{13}h^1M_{\mathrm{}}`$ at $`z=1`$. These masses are only approximate as the group number and volume are both just estimates. The errors on the minimum dark matter halo masses inferred from the observed number densities are likely $`50`$%, given cosmic variance errors and the uncertainties in the corrections made to the number densities from the mock catalogs.
We check for consistency between the observed and predicted clustering of these halos. Mo & White (2002) use the Sheth & Tormen (1999) model to predict the evolution in the clustering of dark matter halos and find that halos of mass $`M_{min}=5.5\times 10^{12}h^1M_{\mathrm{}}`$ will have a clustering amplitude of $`\sigma _8=1.0`$ at $`z=1`$. Here $`\sigma _8`$ is defined as the standard deviation of halo count fluctuations in a sphere of radius 8 $`h^1`$ Mpc; it can be preferable to quoting a scale-length, $`r_0`$, as it removes the significant covariance with $`\gamma `$. $`\sigma _8`$ can be calculated from a power-law fit to $`\xi (r)`$ using the formula,
$$(\sigma _8)^2J_2(\gamma )\left(\frac{r_0}{8h^1\mathrm{Mpc}}\right)^\gamma ,$$
(8)
where
$$J_2(\gamma )=\frac{72}{(3\gamma )(4\gamma )(6\gamma )2^\gamma }$$
(9)
(Peebles, 1980). Using the power-law fits to $`\xi (r)`$ for groups in the DEEP2 data, we find $`\sigma _8=1.0`$, in agreement with the predicted value of Mo & White (2002).
Kravtsov et al. (2004) find using dark matter simulations with the same concordance cosmology that a number density of $`n=6\times 10^4h^3\mathrm{Mpc}^3`$ at $`z=1`$ corresponds to a minimum mass of $`M_{min}=5.9\times 10^{12}h^1M_{\mathrm{}}`$. This value is comparable to the Sheth & Tormen (1999) value quoted above. They predict a clustering amplitude of $`r_0=5.2`$ $`h^1`$ Mpc and $`\gamma =2.16`$ for these halos, which corresponds to $`\sigma _8=1.05`$, in good agreement with our observed value of $`\sigma _8`$ for the DEEP2 groups.
We verify from the mock catalogs that the actual minimum dark matter halo masses for these group samples are similar to those estimated above. For the group sample defined as having an estimated $`\sigma 200`$ km s<sup>-1</sup>, the mass distribution has a rough minimum dark matter halo mass of $`M_{min}=23\times 10^{12}h^1M_{\mathrm{}}`$ and the $`\sigma 350`$ km s<sup>-1</sup> sample has a rough minimum dark matter halo mass of $`M_{min}=45\times 10^{12}h^1M_{\mathrm{}}`$. The mass distributions for the group samples in the mock catalogs do not have a very clearly defined lower-mass cutoff, however, and these values are roughly half of the values quoted above. Even though the galaxies in the mock catalogs are randomly drawn from the velocity distribution of individual dark matter particles, the fact that only a small number of galaxies are observed in a single group leads to significant scatter between the estimated $`\sigma `$ from the observed galaxies and the actual $`\sigma `$ and therefore the mass of the underlying dark matter particles; this scatter likely accounts for the factor of two discrepancy between the halos in the mock catalogs and the predictions of Sheth & Tormen (1999).
## 5. Clustering of Galaxies in Different Environments
In this section the clustering properties for galaxies in groups relative to the full galaxy population and to field galaxies are investigated. Unlike for the group sample, which has larger Poisson errors and can only be measured on scales $`r3`$ $`h^1`$ Mpc, where redshift-space distortions are small, here we are interested in measuring clustering properties of galaxies on small scales where redshift-space distortions are much more significant. Instead of inferring $`\xi (r)`$ in real space from measurements of $`\xi (s)`$ in redshift space, we measure the projected two-point correlation function, $`w_p(r_p)`$, from which $`\xi (r)`$ can be more directly inferred. Color information additionally allows us to divide the sample of group galaxies into red and blue populations and measure $`w_p(r_p)`$ for each. We test the effects of our slitmask target selection algorithm and group-finder on these results using mock catalogs. Finally, we are able to empirically separate the observed $`w_p(r_p)`$ for galaxies in groups into a ‘one-halo’ and a ‘two-halo’ term, by keeping pair counts where both galaxies are in the same or in different groups.
### 5.1. Clustering of Full Galaxy Sample and Galaxies in Groups in DEEP2 data
The full galaxy sample here refers to all DEEP2 galaxies between $`0.7z1.0`$ in the same six pointings for which we have group catalogs. The group and field galaxy samples are identified using the $`\sigma 200`$ km s<sup>-1</sup> group catalog, and these samples combined make up the full galaxy sample. Fig. 5 shows the spatial distribution of galaxies in three DEEP2 pointings, with different symbols for group (open triangles) and field (crosses) galaxies. In the DEEP2 data 39 +/-4% of all galaxies are identified as belonging in recovered groups which have $`\sigma 200`$ km s<sup>-1</sup> in the redshift range $`0.7z1.0`$, where the error quoted is the standard error across the six pointings. This rate is artificially increased by false interlopers, but it is also decreased by our slitmask target selection by roughly the same amount, as estimated using mock galaxy catalogs. In the mock catalogs, 27% of all galaxies are identified as belonging in recovered groups after target selection (24% are in real groups before target selection), significantly less than in the DEEP2 data. This difference is most likely due to the free parameters in the group-finding algorithm having been tuned to reproduce the observed $`n(\sigma ,z)`$ for $`\sigma 350`$ km s<sup>-1</sup>, not the $`\sigma 200`$ km s<sup>-1</sup> cutoff in the present sample (see Fig. 6 in Gerke et al. (2005) for details). The clustering results shown here do not depend on the absolute number of galaxies in each of the samples (all, group, and field galaxies).
Fig. 6 presents $`w_p(r_p)`$ for the full galaxy sample (solid lines), galaxies in groups (dashed lines) and for field galaxies (dotted lines), in both the DEEP2 data (top) and in the mock catalogs (bottom). The top left panel compares $`w_p(r_p)`$ as measured in the observed DEEP2 data (thin lines with no error bars) with $`w_p(r_p)`$ after correcting for effects due to target selection and the group-finder (thick lines with error bars). To make this correction, we have used the ratio of $`w_p(r_p)`$ as a function of scale in the mock catalogs for group galaxies, field galaxies, and all galaxies separately, identified using real groups before target selection (thick lines with error bars in the bottom right panel), to $`w_p(r_p)`$ for galaxies identified using recovered groups after target selection (thick lines with error bars in the bottom left panel). Power-law fits to the corrected $`w_p(r_p)`$ points are shown in the top right panel of Fig. 6 and are listed in table 2, along with fits to the observed points. Field galaxies are well fit by a power-law on scales $`r_p=120`$ $`h^1`$ Mpc. On smaller scales the correlation function is negative, as field galaxies are not found within $`0.6`$ $`h^1`$ Mpc of each other. Pairs of galaxies within that distance are likely to be part of a group.
We will present updated results for $`w_p(r_p)`$ for galaxies in the DEEP2 dataset as a function of galaxy color, luminosity, redshift, etc. in a future paper (Coil et al. 2005, in preparation). Here we focus on the difference between the clustering of galaxies in groups relative to the full galaxy sample (throughout this paper we use the terms “all galaxies” and “the full galaxy sample” interchangeably; they include both field and group galaxies). The full galaxy sample used here includes a total of 9787 galaxies in the redshift range $`0.7z1.0`$ in the same six pointings used for the group analysis. This represents a great advance over the sample used in Coil et al. (2004a), which contained 2219 galaxies in one pointing over the redshift range $`0.7z1.35`$. The fits for $`r_0`$ and $`\gamma `$ here agree with our earlier results, but the errors, both Poisson and cosmic variance, are much smaller in the current sample. For example, we can now address whether the significant rise in slope of $`w_p(r_p)`$ on small scales predicted by Kravtsov et al. (2004) for galaxies at $`z=1`$ based upon their simulations is present in the DEEP2 data. They find that for a galaxy sample with a comoving number density of $`n=1.5\times 10^2h^3\mathrm{Mpc}^3`$ (similar to the DEEP2 sample at $`z0.8`$; see Lin et al. (2004)) that the slope of $`\xi (r)`$ changes from $`\gamma 1.65`$ when measured on scales of $`r=0.310`$ $`h^1`$ Mpc, where $`r_0`$ is measured to be $`4`$ $`h^1`$ Mpc, to $`\gamma 1.9`$ over scales $`r0.110`$ $`h^1`$ Mpc, where $`r_03.5`$ $`h^1`$ Mpc. Fitting our corrected $`w_p(r_p)`$ for the full galaxy sample over the same range in $`r_p`$ results in $`r_p=0.310`$ $`h^1`$ Mpc, $`r_0=3.64\pm 0.07`$ $`h^1`$ Mpc and $`\gamma =1.73\pm 0.04`$, while for scales $`r_p=0.0110`$ $`h^1`$ Mpc, $`r_0=3.64\pm 0.07`$ $`h^1`$ Mpc and $`\gamma =1.73\pm 0.03`$, with no change at all in either amplitude or slope. Therefore no evidence is found for a rise in the slope on small scales as predicted by Kravtsov et al. (2004). We also do not find a significant difference in the slope on scales $`r_p<1`$ $`h^1`$ Mpc compared to scales $`r_p>1`$ $`h^1`$ Mpc, as is seen for the galaxies in groups. Fitting for a power-law on scales $`r_p=0.051`$ $`h^1`$ Mpc, $`r_0=3.52\pm 0.16`$ $`h^1`$ Mpc and $`\gamma =1.78\pm 0.06`$, while on scales $`r_p=120`$ $`h^1`$ Mpc, $`r_0=3.70\pm 0.10`$ and $`\gamma =1.77\pm 0.06`$. The full galaxy sample therefore appears to be consistent with a single power-law slope over the range $`r_p=0.120`$ $`h^1`$ Mpc. We note that in the mock catalogs, the slope measured for $`w_p(r_p)`$ is consistent with the data, though the amplitude of $`r_0`$ is $`7`$% higher. As stated before, the mock catalogs were designed to match the previously published measurements of $`w_p(r_p)`$ for the full DEEP2 galaxy sample.
The correlation function for galaxies in groups is relatively well-fit by a power-law over all scales; a broken power-law fit with a break at $`r_p=1`$ $`h^1`$ Mpc results in a steeper slope on small scales, with low significance. The slope for a single power-law is $`\gamma =2.12\pm 0.06`$ for scales $`r_p=0.0520`$ $`h^1`$ Mpc, while it increases to $`\gamma =2.16\pm 0.11`$ for scales $`r_p=0.051`$ $`h^1`$ Mpc and decreases to $`\gamma =2.02\pm 0.15`$ for scales $`r_p=120`$ $`h^1`$ Mpc. Note the significantly different shape of $`w_p(r_p)`$ for group galaxies in the mock catalogs, which exhibit a strong break at $`r_p1`$ $`h^1`$ Mpc. This sharp rise on small scales is not seen in the DEEP2 data. As we will show in the next subsection, this difference between the mock catalogs and the data is not due to any systematic effect from our slitmask algorithm or in our group identification, as the general shape of the correlation function in the mock catalogs is not changed by these. This difference in the clustering of galaxies in groups between the mock catalogs and the data is the first indication that the mock catalogs, which were constructed to match the $`z1`$ luminosity function and clustering of all galaxies in the DEEP2 data, do not reproduce additional properties of the data. It appears that the spatial distribution of galaxies in groups is less concentrated in the data than in the mock catalogs. We discuss the implications of this in §9.
We also investigate $`w_p(r_p)`$ for group galaxies as a function of color. Fig. 7 shows $`w_p(r_p)`$ for galaxies in the DEEP2 data which are identified as belonging to groups and have red or blue colors, defined by the minimum in the color bi-modality in restframe $`(UB)_0`$, at $`(UB)_0=1.05`$. In the group galaxy sample, 20% of group galaxies are red, while for the full galaxy sample 15% of galaxies are red. Corrections for target selection and our group-finder are made using the mock catalogs as above. Red galaxies in groups have a steeper slope in $`w_p(r_p)`$ and a higher correlation length than blue galaxies; power-law fits result in $`r_0=4.77\pm 0.20`$ $`h^1`$ Mpc and $`\gamma =2.15\pm 0.05`$ for blue group galaxies and $`r_0=5.81\pm 0.45`$ $`h^1`$ Mpc and $`\gamma =2.27\pm 0.11`$ for red group galaxies. The steeper slope for the red galaxy sample implies that red galaxies are more centrally concentrated in their parent dark matter halos; we investigate this more directly in §6. Colors are not currently included in the mock catalogs and so we can not present this measurement for the mock catalogs.
### 5.2. Effects of Slitmask Target Selection and Group-finder
As for the clustering of groups in the DEEP2 data, mock catalogs are used to quantify systematic errors in our measurements of $`w_p(r_p)`$ for each galaxy sample, where real groups are used to define galaxies in groups or in the field. The bottom left panel of Fig. 6 shows $`w_p(r_p)`$ measured for all galaxies, group galaxies, and field galaxies in the mock catalogs, before (thin lines without error bars) and after (thick lines with error bars) target selection. Separate power-law least-squares fits to $`w_p(r_p)`$ in the mock catalogs are performed for the full galaxy sample and for galaxies in groups and in the field; the results before and after target selection are shown in Table 2. Field galaxies are only affected by the target selection on the very smallest scales, $`r_p0.2`$ $`h^1`$ Mpc, and there is no significant change to the correlation function of field galaxies on the scales over which we measure a power-law, $`r_p=120`$ $`h^1`$ Mpc. For the full galaxy sample, our target selection algorithm causes $`w_p(r_p)`$ to be slightly underestimated on small scales ($`r_p2`$ $`h^1`$ Mpc) due to our inability to target all close neighbors. Before target selection is applied, a power-law fits well for the full galaxy sample on scales $`r_p=120`$ $`h^1`$ Mpc, with a steeper slope on small scales, $`r_p=0.11`$ $`h^1`$ Mpc. This change in slope on small scales in the mock catalogs is more significant (2.5$`\sigma `$ vs 1.3$`\sigma `$) for the sample before target selection than after, due to the smaller error bars for the larger sample. We note that the scale at which the slope changes is larger than predicted by Kravtsov et al. (2004) in their simulations, and that this difference in slope is not seen in the DEEP2 data, as discussed above.
For galaxies in groups, however, the effects of target selection are more complicated; it increases their observed clustering on all scales. In further tests, we find that target selection does not affect the measured clustering of galaxies in mock catalogs known to belong in real groups, so long as the group membership is identified before target selection. However, galaxies can be identified as belonging to a group after target selection only if two or more observed galaxies are in that group. We can identify only half of the groups after target selection that were detectable before target selection, and the groups which are preferentially lost are those with a low richness. This causes $`w_p(r_p)`$ for galaxies in groups to be greater when groups are identified after target selection, as only galaxies in the richest, and presumably highest mass and most clustered, groups will be included in the observed sample.
The mock catalogs are also used to investigate the effect of our group-finder on the measured clustering of galaxies in groups, as we want to be sure that our group-finding algorithm is not imposing (even indirectly) a preferred scale for groups, which could potentially artificially cause a change in slope at small scales in $`w_p(r_p)`$ for group galaxies. This is tested by comparing the clustering of galaxies in real and recovered groups in the mock catalogs, where both samples have had the DEEP2 target selection algorithm applied. The results are shown in the bottom right panel of Fig. 6, where both real and recovered groups are seen to have an inflection in $`w_p(r_p)`$ such that the slope rises on small scales ($`r_p<1`$ $`h^1`$ Mpc). the group-finder, then, is not imposing this scale on the clustering results. We also find that the group-finder effectively cancels much of the effect of our slitmask target selection, in that $`w_p(r_p)`$ for galaxies in recovered groups after target selection is similar to $`w_p(r_p)`$ for galaxies in real groups before target selection. The overall correction applied to the observed $`w_p(r_p)`$ for group galaxies in the DEEP2 data is therefore small.
Throughout the paper we compare $`w_p(r_p)`$ for the DEEP2 data with mock catalogs by applying corrections to the observed correlation function to account for effects of our slitmask target selection and group-finder and compare to results in the mock catalogs before target selection for real groups. None of our results change if instead we compare results for the observed $`w_p(r_p)`$ in the data to results in the mock catalogs before target selection for real groups.
### 5.3. One- and Two-Halo Terms of the Group Galaxy Correlation Function
As it is known which group each of the group galaxies belongs to, we can empirically measure the contribution to $`w_p(r_p)`$ from pairs of galaxies in the same or in different groups. This is akin to determining the ‘one-halo’ and ‘two-halo’ terms of $`w_p(r_p)`$ in the halo model language, where each group is identified with a single dark matter halo. The total correlation function is then the sum of these two terms:
$$\xi (r)=[1+\xi _{1h}(r)]+\xi _{2h}(r).$$
(10)
We calculate the ‘one-halo’ and ‘two-halo’ correlation functions using the following estimators:
$$\xi _1=\frac{1}{RR}\left[DD_1\left(\frac{n_R}{n_D}\right)^22DR\left(\frac{n_R}{n_D}\right)+RR\right]$$
(11)
$$\xi _2=\frac{1}{RR}\left[DD_2\left(\frac{n_R}{n_D}\right)^22DR\left(\frac{n_R}{n_D}\right)+RR\right],$$
(12)
where $`DD_1`$ and $`DD_2`$ are pair counts of galaxies within the same group and in different groups, respectively. We then sum these along the line of sight (in the $`\pi `$ direction, to $`\pi _{max}`$) to obtain $`w_{p,1h}`$ and $`w_{p,2h}`$. The projected correlation functions sum such that
$$w_p(r_p)=[\pi _{max}+w_{p,1h}(r_p)]+w_{p,2h}(r_p).$$
(13)
Fig. 8 shows $`w_{p,1h}(r_p)`$ and $`w_{p,2h}(r_p)`$ for galaxies in groups in both the DEEP2 data (left) and for real groups in the mock catalogs (right). The data have been corrected for effects due to our target slitmask and group-finder algorithms.
In the DEEP2 data the scale at which the ‘one-halo’ and ‘two-halo’ terms intersect is $`r_p=1.0`$ $`h^1`$ Mpc; the scale in the mock catalogs is $`r_p=0.5`$ $`h^1`$ Mpc. Exactly where the break occurs between the ‘one-halo’ and ‘two-halo’ terms will presumably depend on the type of groups we are probing; larger groups may have this break at a larger radius. The change in slope seen in $`w_p(r_p)`$ for group galaxies in the mock catalogs is easily understood as the scale at which the ‘one-halo’ and ‘two-halo’ terms equally contribute to $`w_p(r_p)`$. However, the ‘one-halo’ term has a very different shape in the DEEP2 data than in the mock catalog, such that galaxies in the data which belong to the same group are not as clustered on small scales as in the mock catalogs. This is likely due the mock catalogs having an incorrect spatial distribution of galaxies within dark matter halos on small scales; we discuss this further in §9.
Yang et al. (2005c) measure $`w_p(r_p)`$ for group galaxies in 2dF Galaxy Redshift Survey data at $`z0`$ and find that it is not well fit by a single power-law; the ‘one-halo’ term is enhanced relative to the ‘two-halo’ term and there is a rise in $`w_p(r_p)`$ on scales $`r_p12`$ $`h^1`$ Mpc. The strength of the rise depends on the abundance and luminosities of the groups; galaxies in more luminous (and presumably more massive) groups have a larger ‘one-halo’ term and a stronger rise in the slope of $`w_p(r_p)`$ on small scales. It is only for the full galaxy population (including field galaxies) that they find a single power-law fit to $`w_p(r_p)`$. For our group sample at $`z1`$, the shape for $`w_p(r_p)`$ for galaxies in groups is similar to what is seen by Yang et al. (2005c) for groups with a comparable number density. Yang et al. (2005c) find a small rise in the slope of $`w_p(r_p)`$ on scales below $`r_p=1`$ $`h^1`$ Mpc but do not quantify this. Our results at $`z1`$ appear to be similar to their findings at $`z0`$, though with larger errors due to our smaller sample size.
## 6. Group-Galaxy Cross-Correlation Function
In this section we present the cross-correlation function between group centers and the full galaxy sample, which is sensitive to the radial profile of galaxies in and around groups. As with the group and galaxy correlation functions, to avoid redshift-space distortions we measure the projected cross-correlation, $`w_p(r_p)`$. As discussed in §2, errors in the positions of group centers will have some effect on scales $`r_p<1`$ $`h^1`$ Mpc; for this reason we do not plot results for $`r_p<0.3`$ $`h^1`$ Mpc. However, comparisons between data and mock catalogs are unaffected, as the mock catalogs have been treated in an identical manner as the data.
Fig. 9 shows the projected cross-correlation between group centers and the full galaxy sample in the DEEP2 data (left) and the mock catalogs (right). The left panel shows the observed $`w_p(r_p)`$ as a dashed line and the corrected $`w_p(r_p)`$ as a solid line, where corrections for target selection and the group-finding algorithm as a function of scale are made using the ratio of $`w_p(r_p)`$ in the mock catalogs between real group centers and all galaxies before (solid line, right panel) target selection and between recovered group centers and all galaxies after (dashed line, right panel) target selection. The dotted line in the right panel shows the cross-correlation between real groups and all galaxies after target selection. The target selection algorithm has the effect of increasing the cross-correlation on scales $`r_p>0.4`$ $`h^1`$ Mpc, while decreasing the amplitude on smaller scales. The small-scale decrease is due to our inability to target galaxies which are in close projection on the plane of the sky; this causes us to undersample close neighbors. On large scales, the effect is due to the slitmask target algorithm affecting which groups we identify; after target selection we lose many of the pairs of galaxies which were identified as groups before such that we preferentially identify the groups with more observed members, which are presumably more massive and therefore more clustered. The effect of the group-finding algorithm is to increase the cross-correlation on scales $`r_p>0.4`$ $`h^1`$ Mpc, where the group-finder has by definition targetted overdensities in the galaxy distribution.
Comparing the solid lines in the two panels of Fig. 9, which shows $`w_p(r_p)`$ for real groups before target selection to the corrected data, the overall shape of the cross-correlation agrees reasonably well, though the amplitude is somewhat higher in the mock catalogs on both small and large scales. This is consistent with what is seen for the correlation function of galaxies in groups shown in Fig. 6, which are more strongly clustered in the mock catalog than in the DEEP2 data.
We also investigate the dependence of the radial distribution of galaxies in groups on galaxy color. Fig. 10 shows the projected cross-correlation function between either red or blue galaxies and group centers, where again the galaxy sample has been split at the bi-modality in the restframe $`(UB)_0`$ color distribution at $`(UB)_0=1.05`$. Within groups, on small scales, $`r_p0.5`$ $`h^1`$ Mpc, red galaxies are much more strongly clustered than blue galaxies, i.e., red galaxies are preferentially found near the centers of groups. In the DEEP2 data 20% of group galaxies in our sample are red, while 13% of field galaxies and 15% of the full galaxy sample (used in this cross-correlation) are red.
Similar trends are seen at $`z0`$ by Yang et al. (2005b), who measure the cross-correlation between group centers and all galaxies in 2dF and SDSS data. They also find a difference between the radial distribution of red and blue galaxies, though it is only apparent for groups with masses $`M10^{13}h^1M_{\mathrm{}}`$, and the differences are smaller than those found here at $`z1`$.
## 7. Relative Bias Between Groups and Galaxies
Measuring the clustering properties of groups, all galaxies, and galaxies in groups in the DEEP2 data allows us to measure the relative bias between galaxies in groups and all galaxies and between groups and galaxies. Fig. 11 plots the relative bias of group galaxies to the full galaxy sample, which we define as the square root of $`w_p(r_p)`$ for group galaxies (dashed lines in Fig. 6) divided by $`w_p(r_p)`$ for the full galaxy sample (solid lines in Fig. 6), as a function of scale for $`r_p=0.120`$ $`h^1`$ Mpc in both the DEEP2 data (top left panel) and the the mock catalogs (bottom left panel), after correcting for target selection and the group-finder. The bias seen between group galaxies and all galaxies is not surprising, as group galaxies reside in more massive dark matter halos. There is a clear scale-dependence to the relative bias between group galaxies and all galaxies in the DEEP2 data, which falls from $`b_{rel}2.5\pm 0.3`$ at $`r_p=0.1`$ $`h^1`$ Mpc to $`b_{rel}1\pm 0.5`$ at $`r_p=10`$ $`h^1`$ Mpc. The mock catalogs have a much higher relative bias on small scales ($`r_p1`$ $`h^1`$ Mpc) which does not match the bias seen in the data. This reflects the strong rise in slope of the correlation function of galaxies in groups seen on small scales in the mock catalogs.
The ratio of the group center-full galaxy sample cross-correlation function (Fig. 9) to the galaxy correlation function (solid lines in Fig. 6) provides a measure of the relative bias of groups to galaxies, which is shown on the right side of Fig. 11. We have corrected for slitmask target selection and the group-finder. There is some scale-dependence to the relative bias between groups and galaxies in the DEEP2 data (top right panel) and the weighted mean relative bias is $`b_{rel}=1.17\pm 0.04`$ over scales $`r_p=0.515`$ $`h^1`$ Mpc. The mock catalogs have a mean value of $`b_{rel}=1.23\pm 0.02`$ for $`r_p=0.515`$ $`h^1`$ Mpc, in reasonable agreement with the data, though the mock catalogs again show a higher bias on small scales, below $`r_p1`$ $`h^1`$ Mpc, and the agreement with the data is better on scales $`r_p>1`$ $`h^1`$ Mpc. These measures of the relative biases of groups to galaxies and galaxies in groups to all galaxies at $`z1`$ are further constraints which simulations and galaxy evolution models must match, in addition to measures of $`w_p(r_p)`$ for all galaxies and for groups. We discuss the implications of these differences between the data and mock catalogs in the next two sections.
## 8. Effect of Varying the Halo Model Parameters
The differences seen between the clustering of group galaxies on small scales in the DEEP2 data and the mock catalogs could be due to the mock catalogs having the wrong spatial distribution for galaxies within their parent dark matter halos and/or the wrong HOD, which specifies the probability that a dark matter halo of mass M hosts N galaxies, $`P(N|M)`$. To illustrate how much these differences may be due to the HOD used to create the mock catalogs, we investigate the clustering of group galaxies in two mock catalogs with similar number densities and different HODs. In addition to the mock catalog used throughout this paper (labeled as “B256”), we also analyze a mock catalog in which a different HOD was applied to the same dark matter simulation; this model is labeled as “C256” and was chosen as one of the most discrepant HOD models that has an observed $`w_p(r_p)`$ for the full galaxy sample that, by design, matches the results for the DEEP2 data published in Coil et al. (2004a). The HODs for galaxies with $`L>L`$ for these two models are shown in the upper left of Fig. 12, where model B256 is seen to have a lower minimum halo mass hosting a single galaxy, and a steeper slope on larger mass scales ($`0.5`$ compared to $`0.26`$ for the C256 model), which results in having a greater fraction of galaxies residing in massive halos. The curves for galaxies with lower luminosity thresholds have a similar shape and higher amplitude than what is shown here (see Fig. 1 in Yan, White, & Coil (2004)). Mock catalogs made with the C256 model have 35% of galaxies in recovered groups after target selection, similar to what is found for the DEEP2 data (39%), and higher than the value found in the B256 model (27%), even though the C256 model has relatively fewer galaxies in more massive halos. This is due to the parameters of the group-finding algorithm having been tuned to match the observed $`n(\sigma ,z)`$ of the DEEP2 data for $`\sigma 350`$ km s<sup>-1</sup>; we have not re-tuned the group-finder to the C256 mock catalogs or our $`\sigma 200`$ km s<sup>-1</sup> cutoff. Both the B256 and C256 mock catalogs have the same number density for the full galaxy sample. The clustering measures shown here have all been corrected for slitmask target effects and the group-finder.
Fig. 12 shows the correlation function for all galaxies (top right panel) and for group galaxies (bottom left panel) in each of the two halo model mock catalogs. The $`w_p(r_p)`$ for the full galaxy sample is very similar in the two catalogs; the only differences are on scales less than $`r_p0.5`$ $`h^1`$ Mpc, where the C256 model has a slightly higher clustering amplitude. For $`w_p(r_p)`$ for group galaxies, the overall shape is similar for the two models but the amplitude in model B256 is higher at all scales, as this model has more galaxies in massive halos, such that the group galaxies will be more clustered. These figures show that the amplitude of the group galaxy correlation function adds an additional constraint on the HOD, which is not gained from the correlation function of all galaxies alone. It also shows that the general shape of the group galaxy correlation function, and in particular, the rise in slope on small scales, is not sensitive to the parameters of the HOD used.
The bottom right panel of Fig. 12 presents the cross-correlation function of group centers and all galaxies in both mock catalogs. Here there is a difference in the shape of the cross-correlation on small scales for the different HODs. The C256 model has a lower amplitude than the B256 model over almost all scales but shows a distinct rise on the smallest scales, $`r_p0.5`$ $`h^1`$ Mpc, which is not seen in the B256 model. Indeed, both the correlation function for all galaxies and the group-galaxy cross-correlation function in the C256 model show a rise on small scales that is not seen in the B256 model; this results from the C256 model having preferentially more galaxies in smaller mass halos which dominate the pair counts at small separations.
Results from the DEEP2 data are also compared to the different halo model mock catalogs in Fig. 12. By design, $`w_p(r_p)`$ for all galaxies matches both mock catalogs well, though the data do not show the rise on the smallest scales that is seen for the C256 model. The shape of the correlation function for group galaxies in the DEEP2 data does not match either of the halo model mock catalogs; the data show a significantly shallower slope on small scales. The amplitude of the cross-correlation function agrees better with the B256 model than the C256 model, and the shape of the cross-correlation disagrees with the C256 model on small scales. The significant difference in $`w_p(r_p)`$ for group galaxies on small scales is presumably due to a difference in the spatial distribution of galaxies in groups in the data and the mock catalogs, as it does not appear to be reconcilable by altering the HOD. This implies that the spatial distribution of galaxies in groups in the mock catalogs is incorrect. This will be discussed further in the next section.
We note that if the C256 mock catalogs had been used to correct the observed $`w_p(r_p)`$ for the full DEEP2 galaxy sample and galaxies in groups as presented in Fig. 6 for slitmask target effects, none of our conclusions in the paper would change, as the relative differences before and after target selection are similar in the two mock catalogs. The differences for all DEEP2 galaxies if using the C256 mock catalogs to correct for target selection effects are negligable, well within the $`1\sigma `$ errors quoted in Table 2. The differences for group galaxies are within the $`2\sigma `$ errors, with both the corrected $`r_0`$ and $`\gamma `$ being lower ($`r_0=4.72\pm 0.23`$ $`h^1`$ Mpc and $`\gamma =2.05\pm 0.06`$), and there is still no significant difference in the slope of $`w_p(r_p)`$ on small and large scales in the DEEP2 data.
## 9. Discussion and Conclusions
Groups bridge the gap between galaxies and clusters in both mass and scale, and are also the likely locations of galaxy mergers. Measurements of the clustering of groups can constrain cosmological parameters, and measurements of the clustering of galaxies in groups can constrain both halo model parameters and the spatial profile of galaxies in their parent dark matter halos. Here we present the first results on the clustering of groups and galaxies in groups at $`z1`$. We measure four types of correlation function statistics in the DEEP2 dataset: 1) the group correlation function, 2) the galaxy correlation for the full galaxy sample, 3) the galaxy correlation function for galaxies in groups, and 4) the group-galaxy cross-correlation function. The first clustering measure probes the dark matter halo-halo correlation function on mass scales of galaxy groups, which is well-understood from dark matter simulations alone. The clustering of groups in the DEEP2 data at $`z1`$ matches predictions for a $`\mathrm{\Lambda }`$CDM cosmology and is used to estimate the typical dark matter masses of the halos the groups studied reside in.
The second clustering measure, the galaxy correlation function for the full DEEP2 galaxy sample, is an update on results using early DEEP2 data presented in Coil et al. (2004a). Here we present this measurement for the full galaxy sample using a much larger dataset, with over four times as many galaxies covering three fields in the sky; the statistical error on $`r_0`$ is now 2%. We show that $`w_p(r_p)`$ for the full galaxy sample provides constraints on the halo occupation distribution (HOD), the number of galaxies that reside in a dark matter halo of a given mass), though not in an unique way. There is some leeway in how halos can be populated which results in a correlation function that matches our measurements, as shown in the upper right panel of Fig. 12.
The third clustering measure, the correlation function of galaxies in groups, is similar to the second but is restricted to galaxies in more massive halos, as measuring the clustering of galaxies in groups is sensitive to a higher halo mass range than for the full galaxy sample. The contribution from the ‘one-halo’ term is necessarily higher for galaxies in groups as these galaxies are identified as belonging in halos with several other galaxies. This provides further constraints on the halo model parameters than those from the measurement of $`\xi (r)`$ for all galaxies alone. We also find that red galaxies in groups have a steeper slope and higher clustering amplitude than blue galaxies in groups.
The fourth clustering measure, group-galaxy cross-correlation function, reflects the spatial distribution of galaxies within dark matter halos above a given mass, and depends as well upon the parameters of the halo model. We find using the group center-galaxy cross-correlation function that red galaxies are found preferentially in the centers of groups compared to blue galaxies, which has also been seen locally (Collister & Lahav, 2005; Yang et al., 2005b). We find that this trend is in place at $`z1`$.
All four of these measurements are compared to mock catalogs constructed from N-body simulations and depend differently on, and can therefore simultaneously constrain, both parameters of the halo model and the spatial profile of galaxies within halos. Comparing these clustering measurements in the DEEP2 data with the mock catalogs of Yan, White, & Coil (2004), three of the four measures agree fairly well with the simulations, with the exception of the correlation function for galaxies in groups. The clustering amplitude for the full galaxy sample roughly matches the mock catalogs; this is by design: the catalogs were constructed with an HOD that is consistent with earlier DEEP2 clustering results for all galaxies. The HOD used is not uniquely determined however; the observed $`w_p(r_p)`$ for the full DEEP2 galaxy sample can be matched with substantially different HODs (two samples are shown in the upper left panel of Fig. 12). We leave an improved HOD reconstruction from $`w_p(r_p)`$ for the full galaxy sample for a future paper, where we will study the clustering as a function of galaxy properties such as luminosity, color, redshift, etc., in volume-limited samples; here we focus on comparing the clustering of galaxies in groups to all galaxies and the different constraints they provide on the HOD. We do note that $`w_p(r_p)`$ for all galaxies is fit by a single power-law on scales $`r_p=0.0520`$ $`h^1`$ Mpc, with $`r_0=3.63\pm 0.07`$ and $`\gamma =1.74\pm 0.03`$.
While the mock catalogs have similar projected clustering for the full galaxy population to the DEEP2 data by design, there is a strong discrepancy in the clustering of galaxies in groups. The DEEP2 data do not show a significant rise in the slope of $`w_p(r_p)`$ on small scales, for either group galaxies or the full galaxy sample (upper panels of Fig. 6). In contrast, our mock catalogs have a very strong rise on small scales for $`w_p(r_p)`$ for group galaxies, though not for the full galaxy sample (bottom panels of Fig. 6). To test whether this discrepancy can be accounted for by the halo model parameters used, we analyze mock catalogs constructed with a different HOD but similar clustering for the full galaxy sample (Fig. 12). We find that there is still a rise in slope for the clustering of galaxies in groups in the second mock catalogs (model C256) which is not seen in the data. This result is unaffected by our definition of the group center.
The slope of $`w_p(r_p)`$ for group galaxies on small scales should depend quite sensitively on the spatial distribution of galaxies within dark matter halos. We therefore conclude that there is a difference in the spatial distribution of galaxies within their parent dark matter halos in the DEEP2 data and our mock catalogs. The mock catalogs assume no spatial bias, except for the assumption of a central galaxy; the most luminous galaxy in a halo is placed at the center of the halo, while all subsequent galaxies are assigned to random dark matter particles, following an NFW profile. Assuming that the brightest galaxy occupies the very center of the halo is likely not to be correct, as groups at $`z=1`$ are not expected to have a large, dominant, bright galaxy in their centers. This assumption will result in a higher correlation function on small scales.
In our mock catalogs the satellite galaxies are assumed to follow the same NFW profile as the dark matter particles; this appears to not be the correct spatial profile for the galaxy population at $`z=1`$. There is evidence in both simulations and data at $`z0`$ that satellite galaxies do not follow the same spatial profile as the dark matter particles. Simulations have found that subhalos have a shallower spatial profile than the dark matter particles at $`z0`$ (e.g., Gao et al., 2004; Diemand, Moore, & Stadel, 2004; Nagai & Kravtsov, 2005), though exactly how galaxies are related to subhalos is still not entirely known. Observationally, several authors have measured the spatial profiles of galaxies in groups in data at $`z0`$. Using the Two Micron All Sky Survey, Lin, Mohr, & Stanford (2004) stack groups and clusters to measure the radial mass-to-light profile and find that galaxies are less concentrated in the centers of groups and clusters than the dark matter. Collister & Lahav (2005) use the 2dF 2PIGG group catalog to directly measure the radial profile of galaxies within groups and find that galaxies are less centrally concentrated than what is seen for dark matter particles in simulations. Similar results are found by Hansen et al. (2005) for clusters in SDSS, and by Diaz et al. (2005) and Yang et al. (2005b) for groups in SDSS and 2dF.
The cross-clustering between groups and galaxies matches the mock catalogs well on scales $`r_p>0.5`$ $`h^1`$ Mpc; on smaller scales the mock catalogs have a slightly steeper slope. The cross-correlation between groups and galaxies is linearly proportional to the radial distribution of galaxies in groups, while the correlation of group galaxies is proportional to the second power of the radial distribution; this may be why the shape agreement between the data and the mock catalogs is better for the cross-correlation than for the correlation of group galaxies. The group-galaxy cross-correlation function can also be affected by uncertainties in the location of the center of each group which can dilute the signal on small scales, unlike for the correlation function of galaxies in groups. Yang et al. (2005b) find that the cross-correlation between group centers and galaxies at $`z=0.1`$ in 2dF and SDSS data is lower on scales $`r_p<0.1`$ $`h^1`$ Mpc than in their mock catalogs, which do not have a spatial bias with respect to the dark matter distribution. They create a series of mock catalogs with NFW profiles for the galaxies with lower concentration parameters, $`c`$, than in the dark matter and find that catalogs with concentration values of about one-third the value for the dark matter halos match their data well. These mock catalogs show the same trend which is required here, namely a lower cross-correlation on small scales of $`r_p<0.1`$ $`h^1`$ Mpc.
We show here that the clustering properties of galaxies in groups can be used to break degeneracies among different HODs that can not be distinguished by the clustering of all galaxies alone. We find that galaxies in the DEEP2 data do not have the same spatial profile as in our mock catalogs, which assumes a central galaxy surrounded by satellite galaxies following an NFW profile. Using the clustering statistics presented in this paper will allow us to now construct more realistic mock catalogs for the DEEP2 survey that have a better constrained HOD and radial profile for galaxies within dark matter halos.
We would like to thank the anonymous referee for helpful comments and Zheng Zheng for useful discussions. This project was supported by the NSF grant AST-0071048. J.A.N. acknowledges support by NASA through Hubble Fellowship grant HST-HF-01165.01-A awarded by the Space Telescope Science Institute, which is operated by AURA Inc. under NASA contract NAS 5-26555. C.-P. Ma is supported in part by NASA grant NAG5-12173 and NSF grant AST-0407351. S.M.F. would like to acknowledge the support of a Visiting Miller Professorship at UC Berkeley. The DEIMOS spectrograph was funded by a grant from CARA (Keck Observatory), an NSF Facilities and Infrastructure grant (AST92-2540), the Center for Particle Astrophysics and by gifts from Sun Microsystems and the Quantum Corporation. The DEEP2 Redshift Survey has been made possible through the dedicated efforts of the DEIMOS staff at UC Santa Cruz who built the instrument and the Keck Observatory staff who have supported it on the telescope. The data presented herein were obtained at the W.M. Keck Observatory, which is operated as a scientific partnership among the California Institute of Technology, the University of California and the National Aeronautics and Space Administration. The Observatory was made possible by the generous financial support of the W.M. Keck Foundation. The DEEP2 team and Keck Observatory acknowledge the very significant cultural role and reverence that the summit of Mauna Kea has always had within the indigenous Hawaiian community and appreciate the opportunity to conduct observations from this mountain. |
warning/0507/math0507406.html | ar5iv | text | # Singularités des courants d’Ahlfors
## 0. Introduction
<sup>1</sup><sup>1</sup>1”Mots-clés : Kobayashi hyperbolicity, entire curves, currents. <sup>2</sup><sup>2</sup>2”Class.AMS : 32Q45, 32U40.
Soit $`X`$ une variété compacte complexe hermitienne non hyperbolique au sens de Kobayashi. Elle possède alors une courbe entière, i.e. une application holomorphe non constante de $``$ dans $`X`$.
Depuis Ahlfors (voir par exemple ), on sait associer à cette courbe entière un courant positif fermé. Celui-ci s’obtient comme limite de courants d’intégration normalisés $`\frac{[\mathrm{\Delta }_n]}{\text{aire}(\mathrm{\Delta }_n)}`$ sur des disques concentriques de la courbe entière, choisis de sorte que la longueur de leur bord devienne négligeable devant leur aire : $`\text{long}(\mathrm{\Delta }_n)=\text{o(aire}(\mathrm{\Delta }_n))`$.
Ces courants apparaissent notamment dans la preuve par M. McQuillan de la conjecture de Green-Griffiths pour certaines surfaces (, voir aussi ). Celle-ci stipule qu’une courbe entière dans une variété algébrique de type général dégénère algébriquement.
L’objet de cet article est de comprendre la partie singulière de ces courants. Précisément, depuis Siu (voir par exemple ), on sait décomposer tout courant positif fermé en une partie singulière et une partie diffuse. La partie diffuse est peu singulière : elle a des nombres de Lelong nuls hors d’un ensemble dénombrable de points. La partie singulière rend compte des courbes analytiques chargées par le courant : c’est une combinaison des courants d’intégration sur ces courbes.
A la suite de M. Paun (voir qui contient aussi des réponses partielles), il est naturel de conjecturer qu’une courbe analytique chargée par un courant issu d’une courbe entière est rationnelle ou elliptique. En effet, la perte d’hyperbolicité au sens de Kobayashi due à la présence de la courbe entière doit se refléter à l’endroit où le courant se concentre le plus. Nous vérifions cette conjecture dans le cas où la variété $`X`$ est algébrique. La méthode repose sur une propriété de relèvement des arcs de la courbe algébrique à la courbe entière. Celle-ci est liée au caractère laminaire du courant issu de la courbe entière (comparer avec ).
Techniquement, il est commode d’élargir la classe des courants considérés en oubliant les courbes entières et en ne retenant que les suites de disques, voire d’unions finies de disques, dont ils sont limites :
Définition. Un courant $`T`$ est dit d’Ahlfors s’il existe une suite $`(\mathrm{\Delta }_n)`$ de réunions finies de disques holomorphes dans $`X`$ telle que long($`\mathrm{\Delta }_n)=o(\text{aire}(\mathrm{\Delta }_n))`$ et $`T=lim\frac{[\mathrm{\Delta }_n]}{\text{aire}(\mathrm{\Delta }_n)}`$.
Cette définition couvre aussi les courants issus de courbes entières après régulari-sation à la Nevanlinna (voir et ). Voici notre résultat :
###### Th\\'eor\\\`eme
Soit $`T`$ un courant d’Ahlfors dans une variété algébrique $`X`$. On suppose que $`T`$ charge une courbe irréductible $`C`$. Alors $`C`$ est rationnelle ou elliptique.
La démonstration procède par l’absurde en supposant $`C`$ de genre au moins deux. On trouve alors un lacet $`\gamma `$ avec un unique point double dans $`C`$ -une fi-gure huit- qui engendre dans le groupe fondamental de $`C`$ un groupe libre à deux générateurs. Comme $`T`$ se concentre sur $`C`$, on peut relever $`\gamma `$ en un graphe $`\gamma _n`$ dans $`\mathrm{\Delta }_n`$ convergeant vers $`\gamma `$ et possédant beaucoup de cycles (de l’ordre de l’aire de $`\mathrm{\Delta }_n`$). Comme $`\mathrm{\Delta }_n`$ est simplement connexe, le graphe $`\gamma _n`$ va y découper beaucoup de disques (de l’ordre de l’aire de $`\mathrm{\Delta }_n`$). On exhibe ainsi une suite de disques holomorphes d’aire bornée a priori dont le bord converge vers un lacet non trivial de $`\gamma `$. Par un résultat classique de compacité, cette suite produit un disque limite dans $`C`$ bordant ce lacet, c’est la contradiction.
La propriété de relèvement des arcs de $`C`$ à $`\mathrm{\Delta }_n`$, cruciale, s’obtient en analysant $`\mathrm{\Delta }_n`$ sous une projection centrale via la théorie d’Ahlfors de recouvrement des surfaces (voir par exemple ). C’est là qu’intervient l’hypothèse d’algébricité de $`X`$.
Avant de préciser ce schéma, je tiens à remercier chaleureusement M. Paun pour m’avoir soumis ce problème.
## 1. Préliminaires
Enonçons pour commencer le peu de théorie d’Ahlfors et le résultat de compacité des disques dont nous aurons besoin. Dans la suite, les espaces projectifs seront munis de leur métrique standard.
a) Un peu de théorie d’Ahlfors.
Nous renvoyons à la monographie pour les détails de cette théorie. Elle a été créée par Ahlfors pour traduire en termes géométriques la théorie de Nevanlinna de distribution des valeurs.
Précisément, soit $`f:\mathrm{\Sigma }\mathrm{\Sigma }_0`$ une application holomorphe entre deux surfaces de Riemann compactes à bord. On munit $`\mathrm{\Sigma }_0`$ d’une métrique qui induit via $`f`$ une pseudométrique sur $`\mathrm{\Sigma }_0`$. Le résultat central de la théorie est une inégalité de Riemann-Hurwitz approchée pour $`f`$
$$\chi ^{}(\mathrm{\Sigma })s\chi (\mathrm{\Sigma }_0)+c_0l.$$
Ici, la caractéristique d’Euler négative néglige les composantes de $`\mathrm{\Sigma }`$ qui sont des disques ou des sphères, $`s`$ est le nombre moyen de feuillets de $`f`$ : $`s=\frac{\text{aire}(\mathrm{\Sigma })}{\text{aire}(\mathrm{\Sigma }_0)}`$, $`l`$ est la longueur du bord relatif de $`f`$ : $`l=\text{long}(\mathrm{\Sigma }f^1(\mathrm{\Sigma }_0))`$, et la constante $`c_0`$ ne dépend que de $`\mathrm{\Sigma }_0`$ et de sa métrique.
Ahlfors en déduit son théorème des îles, l’analogue du second théorème fondamental de Nevanlinna :
###### Th\\'eor\\\`eme des \\^\iles
Soit $`f:\mathrm{\Delta }P^1()`$ une application holomorphe entre une réunion finie $`\mathrm{\Delta }`$ de disques et $`P^1()`$. On se fixe $`k`$ disques topologiques disjoints dans $`P^1()`$ et on appelle îles de $`f`$ les composantes connexes de la préimage des disques sur lesquelles $`f`$ est propre. Alors le nombre $`i`$ d’îles vérifie $`is(k2)cl`$.
Ici, comme plus haut, $`s`$ désigne le nombre de feuillets moyen de $`f`$, $`l`$ la longueur du bord de $`\mathrm{\Delta }`$, et $`c`$ ne dépend que de la configuration de disques topologiques.
Il suffit en effet d’appliquer l’inégalité de Riemann-Hurwitz à $`\mathrm{\Sigma }_0=P^1()`$ privé des k disques et $`\mathrm{\Sigma }=f^1(\mathrm{\Sigma }_0)`$, sachant que le nombre de feuillets moyen de $`f`$ sur $`\mathrm{\Sigma }_0`$ diffère de $`s`$ par un terme contrôlé par $`l`$.
Ceci entraîne qu’une suite de fonctions méromorphes sur des disques se comporte asymptotiquement comme un revêtement si la longueur du bord devient négligeable devant l’aire. Elle possède la propriété de relèvement suivante :
###### Rel\\\`evement
Soit $`f_n:\mathrm{\Delta }_nP^1()`$ une suite d’applications holomorphes entre des réunions finies $`\mathrm{\Delta }_n`$ de disques et $`P^1()`$. On suppose, avec les notations précédentes, que $`l_n=o(s_n)`$. On se fixe $`ϵ>0`$ et des disques topologiques disjoints dans $`P^1()`$ en nombre $`k>4/ϵ`$. Alors un de ces disques topologiques possède au moins $`s_n(1ϵ)`$ relèvements via $`f_n`$ pour $`n`$ assez grand.
Ici, un relèvement est une composante connexe de la préimage du disque sur laquelle $`f_n`$ est bijective.
En effet, par le théorème des îles, un des disques topologiques va posséder au moins $`s_n(1\frac{ϵ}{2})`$ îles dans sa préimage pour $`n`$ assez grand. Or la restriction de $`f_n`$ à chaque île est un revêtement ramifié au-dessus du disque. Ce revêtement sera donc de degré 1 pour au moins $`s_n(1ϵ)`$ îles, puisque le nombre moyen de feuillets de $`f_n`$ sur ce disque est asymptotiquement $`s_n`$.
b) Compacité de disques holomorphes.
Voici l’énoncé qui nous sera utile. Il remonte au moins à Gromov dans un contexte bien plus large (voir par exemple ).
###### Compacit\\'e
Soit $`f_n:DP^2()`$ une suite d’applications holomorphes du disque unité $`D`$ dans $`P^2()`$. On suppose l’aire de $`f_n(D)`$ uniformément bornée. Alors, quitte à extraire,
i) $`f_n`$ converge localement uniformément vers $`f_{\mathrm{}}`$ sur $`D`$ privé d’un nombre fini de points d’explosion;
ii) $`f_{\mathrm{}}`$ se prolonge en un disque holomorphe encore noté $`f_{\mathrm{}}:DP^2()`$;
iii) en un point d’explosion $`e`$, il existe une suite de disques $`d_n`$ dans $`D`$ tendant vers $`e`$ telle que $`f_n(d_n)`$ converge au sens de Hausdorff vers une courbe rationnelle (une bulle).
En particulier, si l’aire de $`f_n(D)`$ reste bornée par $`\frac{1}{2}`$ il ne se produira pas d’explosion puisqu’une courbe rationnelle est d’aire au moins 1. On aura dans ce cas convergence uniforme locale de $`f_n`$ vers $`f_{\mathrm{}}`$ sur $`D`$ après extraction.
## 2. Démonstration du théorème
Notons que, par plongement de $`X`$ dans un espace projectif, il suffit de regarder les courants d’Ahlfors de $`P^N()`$ chargeant une courbe projective. La démonstration procède en trois étapes : la première est une réduction à $`P^2()`$, la seconde traite de la propriété de relèvement des arcs, tandis que la troisième détaille la fin de l’argument.
Pour des précisions sur les courants et leurs nombres de Lelong, nous renvoyons le lecteur à . Dans la suite on ne distinguera plus les disques paramétrés de leur image géométrique. En particulier les aires des images seront comptées avec multiplicité. Que le lecteur nous pardonne cette liberté.
a) Réduction de dimension.
Soit donc un courant d’Ahlfors $`T`$ de $`P^N()`$ chargeant une courbe projective $`C`$. On projette $`T`$ dans un hyperplan projectif par une projection $`\pi `$ de centre $`p`$. Si $`p`$ est choisi assez générique, le nombre de Lelong de $`T`$ en $`p`$ sera nul et la courbe $`\pi (C)`$ aura même genre que $`C`$ tant que $`N3`$. Le lemme suivant permet donc de se ramener au cas $`N=2`$ par projections centrales successives.
###### Lemme
Le courant projeté $`\stackrel{~}{T}`$ est un courant d’Ahlfors chargeant la courbe projective $`\pi (C)`$.
###### Demonstration Démonstration
La difficulté réside dans la singularité de $`\pi `$ en $`p`$. Précisons d’abord la définition de $`\stackrel{~}{T}`$ : si $`B_r`$ désigne la boule de centre $`p`$ et de rayon $`r`$, on pose $`\stackrel{~}{T}=lim_{r0}\pi _{}(T|_{(B_r)^c})`$. Clairement $`\stackrel{~}{T}`$ charge $`\pi (C)`$. Reste à vérifier que c’est un courant d’Ahlfors.
Pour cela, partons de la suite de réunions finies $`\mathrm{\Delta }_n`$ de disques holomorphes donnant $`T`$. Elle satisfait $`lim\frac{[\mathrm{\Delta }_n]}{a_n}=T`$ et long($`\mathrm{\Delta }_n)=o(a_n)`$ avec $`a_n=\text{aire}(\mathrm{\Delta }_n)`$. La suite analogue $`\stackrel{~}{\mathrm{\Delta }}_n`$ pour $`\stackrel{~}{T}`$ s’obtient essentiellement en projetant par $`\pi `$ la partie de $`\mathrm{\Delta }_n`$ loin de $`p`$. Le contrôle de la longueur de son bord va provenir de l’annulation du nombre de Lelong de $`T`$ en $`p`$. Précisons ceci.
Soit $`r>0`$ fixé provisoirement. Dans ce qui suit, les inégalités seront à constante multiplicative indépendante de $`r`$ et $`n`$ près. Voici comment se traduit la nullité du nombre de Lelong de $`T`$ en $`p`$ pour $`n`$ assez grand :
$$\text{aire}(\mathrm{\Delta }_nB_r)ϵ(r)r^2a_n.$$
$`1`$
Ici $`ϵ(r)`$ tendra vers $`0`$ avec $`r`$.
L’inégalité de coaire (voir par exemple ) fournit alors $`r_n`$ avec $`\frac{r}{2}<r_n<r`$ et
$$\text{long}(\mathrm{\Delta }_nB_{r_n})ϵ(r)ra_n.$$
$`2`$
Remarquons que $`\mathrm{\Delta }_n(B_{r_n})^c`$ est une réunion de disques troués. On rebouche ces trous par les composantes connexes correspondantes de $`\mathrm{\Delta }_nB_{r_n}`$. On obtient ainsi une réunion $`\mathrm{\Delta }_n^{}`$ de disques dont le bord consiste en deux parties, l’une incluse dans $`\mathrm{\Delta }_nB_{r_n}`$, l’autre étant $`\mathrm{\Delta }_n(B_{r_n})^c`$. Notons $`\stackrel{~}{\mathrm{\Delta }}_n`$ la projection de $`\mathrm{\Delta }_n^{}`$ par $`\pi `$. C’est la suite recherchée.
Puisque $`\pi `$ dilate au plus les longueurs d’un facteur $`\frac{1}{r}`$ sur $`(B_{r_n})^c`$, on a, d’après (2),
$$\text{long}(\stackrel{~}{\mathrm{\Delta }}_n)ϵ(r)a_n$$
$`3`$
pour $`n`$ assez grand.
Comparons maintenant $`a_n`$ avec $`\stackrel{~}{a}_n`$ l’aire de $`\stackrel{~}{\mathrm{\Delta }}_n`$. Pour cela, notons $`\omega `$ et $`\stackrel{~}{\omega }`$ les formes de Fubini-Study respectives de $`P^N()`$ et de l’hyperplan sur lequel on projette. On a $`a_n=_{\mathrm{\Delta }_n}\omega `$ et $`\stackrel{~}{a}_n=_{\mathrm{\Delta }_n^{}}\pi ^{}\stackrel{~}{\omega }`$. Remarquons que $`\pi ^{}\stackrel{~}{\omega }\omega =dd^cu`$$`u`$ est une fonction lisse hors de $`p`$ avec une singularité logarithmique en $`p`$. En particulier $`d^cu`$ est de l’ordre de $`\frac{1}{r}`$ sur $`(B_{r_n})^c`$. Il s’ensuit, par le théorème de Stokes et les estimées (1) et (3), que
$$|\stackrel{~}{a}_na_n|_{\mathrm{\Delta }_n\mathrm{\Delta }_n^{}}\omega +|_{\mathrm{\Delta }_n^{}}d^cu|\text{aire}(\mathrm{\Delta }_nB_{r_n})+\text{long}(\stackrel{~}{\mathrm{\Delta }}_n)ϵ(r)a_n$$
pour $`n`$ assez grand. Donc long($`\stackrel{~}{\mathrm{\Delta }}_n)ϵ(r)\stackrel{~}{a}_n`$.
Faisons tendre maintenant $`r`$ vers 0. Après extraction diagonale sur les suites $`(\stackrel{~}{\mathrm{\Delta }}_n)`$ produites à $`r`$ fixé, on vérifie que $`\stackrel{~}{T}=lim\frac{[\stackrel{~}{\mathrm{\Delta }}_n]}{\stackrel{~}{a}_n}`$ et long($`\stackrel{~}{\mathrm{\Delta }}_n)=o(\stackrel{~}{a}_n)`$. C’est donc bien un courant d’Ahlfors. $`\mathrm{}`$
Remarque. Le même type d’estimée montre que l’aire de $`\pi (\mathrm{\Delta }_n^{}B_{r_n})`$ est contrôlée par $`ϵ(r)a_n`$ pour $`n`$ assez grand.
b) Relèvement des arcs.
On s’est ramené à un courant d’Ahlfors $`T`$ de $`P^2()`$ chargeant une courbe projective $`C`$. Précisément $`T|_C=\nu [C]`$ avec $`\nu >0`$.
Analysons encore $`T`$ sous une projection $`\pi `$ de centre $`p`$ sur une droite projective $`L`$. Si $`p`$ est choisi assez générique, on modifie comme plus haut les réunions $`\mathrm{\Delta }_n`$ de disques définissant $`T`$ pour que les longueurs des bords de $`\mathrm{\Delta }_n`$ et $`\pi (\mathrm{\Delta }_n)`$ deviennent négligeables devant l’aire de $`\mathrm{\Delta }_n`$, celle-ci soit équivalente à l’aire de $`\pi (\mathrm{\Delta }_n)`$, et
$$\underset{r}{lim}\underset{n}{lim\; sup}\frac{\text{aire}(\pi (\mathrm{\Delta }_nB_r))}{\text{aire}(\mathrm{\Delta }_n)}=0.$$
$`4`$
Notons $`s_n`$ l’aire de $`\pi (\mathrm{\Delta }_n)`$. C’est le nombre moyen de feuillets de $`\pi :\mathrm{\Delta }_nL`$. Dans la suite on appellera arc un chemin injectif dans $`C`$ qui se projette injectivement par $`\pi `$ dans $`L`$. On dira qu’un arc $`\alpha `$ est $`ϵ`$-relevable si on peut l’épaissir en une bande $`\beta `$ dans $`C`$ possédant au moins $`s_n(\nu ϵ)`$ relevés dans $`\mathrm{\Delta }_n`$ après extraction, qui convergent vers $`\beta `$. On parlera de relevés proches. Ici un relevé de $`\beta `$ est une composante $`\beta ^{}`$ de $`\pi ^1(\pi (\beta ))`$ dans $`\mathrm{\Delta }_n`$ telle que $`\pi :\beta ^{}\pi (\beta )`$ soit bijective.
###### Lemme
Tout arc $`\alpha `$ de $`C`$ peut se perturber en un arc $`ϵ`$-relevable.
Démonstration (comparer à ). On se fixe un voisinage $`V`$ de $`\alpha `$ dans $`C`$ qui se projette injectivement par $`\pi `$ sur $`U`$. S’il est suffisamment fin, la masse de $`T`$ dans $`\pi ^1(U)`$ est inférieure à $`\frac{ϵ}{4}`$. Soient maintenant $`k`$ perturbations disjointes de $`\alpha `$ dans $`V`$, avec $`k>\frac{8}{ϵ}`$. On les épaissit un peu en $`k`$ fines bandes disjointes.
Appliquons la propriété de relèvement 1.a) à $`\pi :\mathrm{\Delta }_nL`$ et aux $`k`$ disques disjoints qui sont les images par $`\pi `$ de ces bandes. On obtient qu’une de ces bandes $`\beta _0`$ a une projection $`\pi (\beta _0)=\stackrel{~}{\beta }`$ possédant $`s_n(1\frac{ϵ}{2})`$ relevés dans $`\mathrm{\Delta }_n`$ pour $`\pi `$, si $`n`$ est assez grand et après extraction. Parmi ceux-ci, au moins $`s_n(1ϵ)`$ sont d’aire inférieure à $`1/2`$ par le contrôle de la masse de $`T`$ dans $`\pi ^1(U)`$. On les note $`\beta _{n,i}`$. Reste à en localiser une partie près de $`\beta _0`$.
Introduisons pour cela un peu de terminologie. Soient $`Q`$ le cône $`\pi ^1(\stackrel{~}{\beta })`$ et $`B`$ l’ensemble des disques holomorphes $`\beta `$ de $`Q`$ d’aire inférieure à 1/2 et se projetant bijectivement par $`\pi `$ sur $`\stackrel{~}{\beta }`$. On peut voir les disques de $`B`$ comme des sections de $`\pi `$ au-dessus de $`\stackrel{~}{\beta }`$.
Par 1.b) une suite dans $`B`$ possèdera toujours une suite extraite convergente puisque la borne d’aire interdit les explosions. La limite est encore dans $`B`$ ou vaut $`p`$. Ainsi $`B`$ est relativement compact pour la convergence uniforme locale sur $`\stackrel{~}{\beta }`$, et l’espace $`M`$ des mesures positives sur $`B`$ de masse au plus 1 est compact pour la convergence faible.
A une mesure $`\lambda `$ de $`M`$ on associe un courant “géométrique” $`T_\lambda `$ dans $`Q`$ par la formule $`T_\lambda =_B[\beta ]𝑑\lambda (\beta )`$. Cette application de $`M`$ sur l’espace $`G`$ des courants géométriques étant continue, on en déduit que $`G`$ est compact pour la topologie faible.
Revenons maintenant à nos relevés $`\beta _{n,i}`$. Par construction, la suite $`(\frac{1}{s_n}_i[\beta _{n,i}])`$ est une suite de courants géométriques $`(T_{\lambda _n})`$ pour $`\lambda _n=\frac{1}{s_n}_i\delta _{\beta _{n,i}}`$. Elle converge après extraction vers $`T_\lambda `$ dans $`B`$. Comme $`T_{\lambda _n}\frac{1}{s_n}[\mathrm{\Delta }_nQ]`$ on a $`T_\lambda T|_Q`$.
Par ailleurs $`\lambda `$ a une masse minorée par $`1ϵ`$ car le nombre de disques $`\beta _{n,i}`$ s’échappant vers $`p`$ est négligeable d’après (4). Donc $`\pi _{}(T_\lambda )>(1ϵ)[\stackrel{~}{\beta }].`$ Comme $`\pi _{}(T|_Q)=[\stackrel{~}{\beta }]`$ on en déduit
$$\pi _{}(T|_QT_\lambda )<ϵ[\stackrel{~}{\beta }].$$
Par hypothèse $`T`$ charge $`C`$ avec un poids $`\nu `$ donc $`T|_{\beta _0}=\nu [\beta _0]`$. Par construction $`T_\lambda |_{\beta _0}=\mu [\beta _0]`$$`\mu `$ est la charge de $`\beta _0`$ pour $`\lambda `$. Donc
$$\pi _{}(T|_QT_\lambda )(\nu \mu )[\stackrel{~}{\beta }].$$
Ainsi la charge de $`\beta _0`$ pour $`\lambda `$ est supérieure à $`\nu ϵ`$. Il en est de même de la charge pour $`\lambda _n`$ de tout voisinage de $`\beta _0`$ dans $`B`$ pour $`n`$ assez grand. Autrement dit, au moins $`s_n(\nu ϵ)`$ disques $`\beta _{n,i}`$ convergent vers $`\beta _0`$. $`\mathrm{}`$
Un graphe $`\gamma `$ de $`C`$ est une réunion d’un nombre fini $`a`$ d’arcs $`\alpha _j`$ se coupant en un nombre fini $`i`$ de points. Deux arcs sont disjoints ou se coupent exactement une fois, soit en leurs extrémités (arcs consécutifs), soit en leurs intérieurs (arcs en croix). Le graphe $`\gamma `$ sera dit relevable si ses arcs sont $`ϵ`$-relevables pour $`ϵ=\frac{\nu }{2a+6i}`$. Grâce au lemme on peut toujours perturber un graphe en un graphe relevable en préservant sa combinatoire d’incidence. Notons $`\mathrm{\Gamma }`$ la réunion des bandes $`\beta _j`$ épaississant les arcs et $`\mathrm{\Gamma }_n`$ la réunion de leurs relevés proches dans $`\mathrm{\Delta }_n`$. C’est l’épaississement d’un graphe $`\gamma _n`$ dans $`\mathrm{\Delta }_n`$ qui converge vers $`\gamma `$.
###### Fait
On a $`\chi (\mathrm{\Gamma }_n)s_n(\chi (\mathrm{\Gamma })+\frac{1}{2})\nu `$ pour $`n`$ assez grand.
En effet les bandes $`\beta _j`$ sont disjointes ou se coupent sur des disques qui épais-sissent les intersections des $`\alpha _j`$. Donc $`\chi (\mathrm{\Gamma })=ai`$. Par ailleurs les relevés proches des bandes sont disjoints ou se recollent sur des relevés proches de ces disques. Donc $`\chi (\mathrm{\Gamma }_n)=a_ni_n`$$`a_n`$ désigne le nombre total de relevés proches de bandes et $`i_n`$ le nombre total de recollements.
Estimons ces quantités. Pour cela remarquons que les relevés proches d’une bande $`\beta _j`$ donnée (ou d’un disque) sont en nombre majoré par $`s_n(\nu +ϵ)`$ pour $`n`$ assez grand. On aurait sinon $`T|_{\beta _j}(\nu +ϵ)[\beta _j]`$ par passage à la limite.
Donc $`a_ns_n(\nu +ϵ)a`$. Par ailleurs, au-dessus d’un disque intersection de deux bandes, on a au moins $`s_n(\nu ϵ)`$ relevés proches provenant de chacune des bandes. Il s’ensuit nécessairement au moins $`s_n(\nu 3ϵ)`$ recollements au-dessus de ce disque. Ainsi $`i_ns_n(\nu 3ϵ)i`$. On obtient donc par notre choix de $`ϵ`$
$$\chi (\mathrm{\Gamma }_n)s_n(\chi (\mathrm{\Gamma })\nu +ϵ(a+3i))s_n(\chi (\mathrm{\Gamma })+\frac{1}{2})\nu .$$
c) Fin de l’argument.
Supposons par l’absurde $`C`$ de genre au moins 2. Soit $`\gamma `$ une figure huit (un lacet avec un unique point double) dans $`C`$ dont le groupe fondamental s’injecte dans celui de $`C`$. On reprend ici la terminologie du paragraphe précédent. On peut toujours supposer après perturbation que $`\gamma `$ est un graphe relevable.
Comme $`\chi (\gamma )=1`$ on a, par l’estimée précédente, $`\chi (\gamma _n)\frac{s_n\nu }{2}`$ pour $`n`$ assez grand. Par un argument de dualité d’Alexander, $`\gamma _n`$ découpe donc dans $`\mathrm{\Delta }_n`$ au moins $`\frac{s_n\nu }{2}`$ disques. Or $`s_n`$ équivaut à l’aire de $`\mathrm{\Delta }_n`$. En sélectionnant pour chaque $`n`$ le disque $`\delta _n`$ composante de $`\mathrm{\Delta }_n\gamma _n`$ d’aire minimale, on produit ainsi une suite de disques holomorphes d’aire bornée a priori.
Analysons la convergence de $`(\delta _n)`$. Commençons par leur bord $`\delta _n`$. Il est constitué d’un cycle de relevés proches d’arcs de $`\gamma `$. Par construction ceux-ci s’épaississent dans $`\mathrm{\Delta }_n`$ en relevés proches des demi-bandes correspondantes dans $`\mathrm{\Gamma }`$. Ces relevés forment un anneau $`A_n`$ dans $`\delta _n`$ bordant $`\delta _n`$ et d’aire comparable à la longueur de $`\delta _n`$. Il s’ensuit que le nombre de relevés proches d’arcs de $`\gamma `$ dans $`\delta _n`$ est borné a priori. Donc, quitte à extraire, $`\delta _n`$ converge vers un cycle non trivial d’arcs de $`\gamma `$. De même $`A_n`$ converge vers un anneau $`A`$ immergé dans $`C`$. En particulier le module de $`A_n`$ reste minoré. Cela permet de choisir un paramétrage holomorphe $`f_n`$ des disques $`\delta _n`$ par le disque unité $`D`$ tel que la préimage de $`A_n`$ contienne un anneau fixe $`(r<|z|<1)`$.
Appliquons l’énoncé de compacité 1.b). Comme $`C`$ n’est pas rationnelle, les éventuelles explosions ne peuvent se produire que dans $`(|z|r)`$. Ainsi le disque limite $`\delta _{\mathrm{}}=f_{\mathrm{}}(D)`$ contient $`A`$. Par prolongement analytique il est entièrement dans $`C`$. Or $`\delta _{\mathrm{}}`$ est un lacet non trivial de $`\gamma `$ donc non homotope à zéro dans $`C`$. C’est la contradiction. $`\mathrm{}`$ |
warning/0507/astro-ph0507246.html | ar5iv | text | # Higher Criticism Statistic: Detecting and Identifying Non-Gaussianity in the WMAP First Year Data
## 1 Introduction
The Standard Inflationary model solves the horizon, the flatness, and the monopole problems, and provides a framework for the formation of structure in the universe (Guth 1981, Guth & Pi 1982). Regarding the latter, the Standard Inflationary model predicts the existence of quantum density fluctuations that are amplified during the inflationary period and that grew, through gravitational instabilities, into the galaxies and clusters that populate our universe. These primordial density fluctuations are predicted to form a homogeneous and isotropic Gaussian field. The predicted statistical distribution of the Cosmic Microwave Background (CMB) temperature fluctuations reflects that of the primordial density fluctuations. Testing this prediction has been the aim of many works in the literature. In particular, since the release of the first year of data collected by the WMAP (Wilkinson Microwave Anisotropy Probe) satellite (Bennett et al. 2003), a considerable number of papers have presented different statistical analyses based on data in real space, spherical harmonics, and wavelet space (Park 2004, Eriksen et al. 2004a, Eriksen et al. 2004b, Eriksen et al. 2005, Larson & Wandelt 2004, Hansen et al. 2004, Prunet et al. 2005, Vielva et al. 2004, Cruz et al. 2005, Mukherjee & Wang 2004, Mc.Ewen et al. 2004). All these works have claimed the detection of deviations from the predictions of the Standard Inflationary model in the WMAP data set optimal for CMB studies (see Komatsu et al. 2003). Several other works have presented statistical methods shown to be very powerful in detecting deviations from the Standard Inflationary model in the so-called Internal Linear Combination (ILC) map (Chiang et al. 2003, Coles et al. 2004, Copi et al. 2004, Naselsky et al. 2003, Chiang & Naselsky 2004). In this case the most convincing source of deviations is foreground related.
Deviations from the predictions of the Standard Inflationary model can have a cosmological origin. Non-Gaussianity can be generated under different conditions. A review on the predictions from several alternative scenarios including multi-field inflation, inhomogeneous reheating, non-linearities in the gravitational potential, and the curvaton-based model can be found in Bartolo et al. 2004. However, these deviations can also be caused by systematic effects or noise associated with the experiment as well as foregrounds (Galactic or extra-galactic). The first year of data collected by the WMAP satellite has undergone careful characterization and examination by the WMAP team in an attempt to fully understand the data set (Page et al. 2003, Hinshaw et al. 2003, Bennett et al. 2003b, Barnes et al. 2003, Jarosik et al. 2003). The detection of deviations indicated in the previous paragraph show the presence of some asymmetry between the northern and southern hemispheres. Moreover, analyses in wavelet space in different regions (north, south, northeast, northwest, southeast, southwest of the Galactic plane) performed by Vielva et al. 2004 and Cruz et al. 2005 indicate that the source of deviations might be a cold spot located at $`(l209^{},b57^{})`$. The nature of the observed deviations is still not clear.
The implementation and development of new statistical methods is indispensable to improve our understanding of the deviations from the Standard Inflationary model predictions, in particular those observed in the WMAP data. There are many kinds of non-Gaussianity, and each type may be sensitive to some statistical tests but immune to others. We need more types of tests, such as Higher Criticism in order to better understand non-Gaussianity. The Higher Criticism ($`HC`$) statistic was first proposed in Donoho & Jin (2004) for testing a very subtle problem: given $`n`$ Gaussian observations with same standard deviations but different means, we want to test whether $`n`$ means are all zeros versus the alternative that a small fraction of them is nonzero. The fraction of nonzero means is too small to have any effect on the bulk of the observations, so all statistics based on moments (e.g. excess kurtosis ($`\kappa `$)) would generally have no power for detection. With careful calibrations, $`HC`$ was proved to be optimal in detecting such situations. $`HC`$ is also competitive for detection in other scenarios. An analysis of simulated CMB on flat patches of the sky using the $`HC`$ statistics was presented in Jin et al. 2004. The simulated images included CMB signals, as well as cosmic strings and Sunyaev-Zeldovich type emission. Analyses were performed in wavelet and curvelet spaces (Candés & Donoho 2000). The very particular non-Gaussianity introduced by the simulated effects, as well as the restricted size of the sample, resulted in the lower sensitivity of this method in comparison with the traditional $`\kappa `$ statistic. However, theoretical studies indicate that the $`HC`$ statistic will outperform the $`\kappa `$ when the sample size gets larger.
Besides its competitive detection power, $`HC`$ is also valuable in identifying the origin of detected deviations, offering a way to determine which portion of the data contributes most to the deviation; this property is not available for many other statistics such as $`\kappa `$. As mentioned above, analyses of the WMAP data in wavelet space suggest that the source of the detected deviations is due to the cold temperatures of certain pixels within a spot located in the southern hemisphere. The $`HC`$ test is especially sensitive to abnormally high amounts of moderately large observations and, as shown in this paper, its application to WMAP data shows deviations from the predictions of the Standard Inflationary model. Moreover, inspection of individual pixel $`HC`$ values in the WMAP data provides a direct means to identify the pixels at the source of the deviation. No additional region-by-region analysis is needed. The computational cost is therefore strongly reduced, making it a promising test for analyzing future higher resolution data, such as data from the Planck mission<sup>1</sup><sup>1</sup>1http://www.rssd.esa.int/index.php?project=PLANCK.
The paper is organized as follows. The $`HC`$ statistical test is described in detail in Section 2. A comparison between the performance of this statistic and that of the kurtosis ($`\kappa `$) and the Maximum ($`Max`$) is included in subsection 2.1. The WMAP data is analyzed in real and wavelet spaces. A description of the combination of WMAP data in different channels, the Monte Carlo simulations, and the process followed in wavelet space is included in Section 3. The analysis of the WMAP data is presented in Section 4. Section 5 is dedicated to conclusions and discussion.
## 2 Higher Criticism
The $`HC`$ statistic was first proposed in Donoho & Jin 2004, Jin 2004. Given $`n`$ independent observations of a distribution which is thought to be slightly deviated from the standard Gaussian, one can compare the fraction of observed significances at a given $`\alpha `$-level (i.e. the number of observed values exceeding the upper-$`\alpha `$ quantile of the standard Gausian) to the expected fraction under the standard Gaussian assumption:
$$\sqrt{n}|(FractionofSignificanceat\alpha )\alpha |/(\alpha (1\alpha ))^{0.5}.$$
The $`HC`$ statistic is then defined as the maximum of the above quantities over all significance levels $`0<\alpha <1`$. Given $`n`$ individual observations $`X_i`$ from a distribution which is thought to be symmetric and slightly deviated from the standard Gaussian, there is a simpler equivalent form of HC defined as follows. First, we convert the individual $`X_i`$’s into individual $`p`$-values: $`p_i=P\{|N(0,1)|>|X_i|\}`$, then we let $`p_{(1)}<p_{(2)}<\mathrm{}<p_{(n)}`$ denote the $`p`$-values sorted in ascending order, define:
$$HC_{n,i}=\sqrt{n}\left|\frac{i/np_{(i)}}{\sqrt{p_{(i)}(1p_{(i)})}}\right|,$$
the $`HC`$ statistic is then:
$$HC_n^{}=\underset{i}{\mathrm{max}}HC_{n,i}$$
or in a modified form:
$$HC_n^+=\underset{\{i:\mathrm{\hspace{0.33em}1}/np_{(i)}11/n\}}{\mathrm{max}}HC_{n,i};$$
we let $`HC_n`$ refer either to $`HC_n^{}`$ or $`HC_n^+`$ whenever there is no confusion. The above definition is slightly different from that in Donoho & Jin 2004, but the ideas are essentially the same.
$`HC`$ is useful in non-Gaussianity detection when $`X_i`$ are truly from $`N(0,1)`$, with the result that $`HC_{n,i}`$ is approximately distributed as $`N(0,1)`$ for almost every $`i`$. Thus an unusually large $`HC_n^{}`$ or $`HC_n^+`$ value strongly implies non-Gaussianity. Moreover, $`HC_{n,i}`$ also provides localized information on deviations from Gaussianity. We can track down the source of non-Gaussianity by studying which portion of the data gives unusually large $`HC_{n,i}`$.
Previous works have claimed the detection of deviations from the predictions of the Standard Inflationary model in the First Year WMAP data, using the $`\kappa `$ statistic in wavelet space (Vielva et al. 2004, Mukherjee & Wang 2004, McEwen et al. 2004). In order to establish a comparison between this statistic and the $`HC`$ we will include calculations of $`\kappa `$ in our analysis as well. For completeness, we will also use the so-called $`Max`$ statistic. The definitions for these two statistics are provided below. A discussion about the theoretical power of the three statistics to detect deviations from Gaussianity is included in the following subsection. $`Kurtosis`$, $`\kappa `$. For a (symmetric) random variable $`X`$, $`\kappa `$ is $`\kappa (X)=\frac{E[X^4]}{(E[X^2])^2}3`$, which uses the 4th moment to measure the departure from Gaussianity. $`\kappa `$ is useful in non-Gaussianity detection as $`\kappa (X_1,X_2,\mathrm{},X_n)\kappa (X)`$ for large $`n`$. $`Max`$. The largest (absolute) observation is a classical statistic:
$$M_n=max\{|X_1|,|X_2|,\mathrm{},|X_n|\}.$$
$`Max`$ is useful in non-Gaussianity detection because $`M_n\sqrt{2\mathrm{log}n}`$ when $`X_i`$ are truly from $`N(0,1)`$; thus a significant difference between $`M_n`$ and $`\sqrt{2\mathrm{log}n}`$ implies non-Gaussianity.
### 2.1 Comparison of Higher Criticism, Max and Excess Kurtosis
The aim of this section is to establish a simple theoretical comparison between the three statistics applied in this paper. We show the power of the different statistical tests in detecting the distortion generated by a faint non-Gaussian signal (modeled as a function of the decaying rate of the tail of the distribution) superposed on a Gaussian signal. Similar to the analysis presented in Jin et al. 2004, the superposed image can be thought of as $`Y=N+G`$, where $`Y`$ is the observed image, $`N`$ is the non-Gaussian component, and $`G`$ is the Gaussian component (assumed to have mean zero and dispersion one, $`N(0,1)`$). We study the power of the three statistics in testing whether $`N=0`$ or not.
One can do such a test either in real space (the space of the observations) or, as it is done in this paper, in wavelet space. For sufficiently fine resolution, the wavelet coefficients $`X_i`$ of $`Y`$ can be modeled as:
$$X_i=\sqrt{1\lambda }z_i+\sqrt{\lambda }w_i,1in,$$
where $`n`$ is the number of observations, $`1\lambda 0`$ is a parameter, $`z_i\stackrel{iid}{}N(0,1)`$ are the transform coefficients of the Gaussian component $`G`$, and $`w_i\stackrel{iid}{}W`$ are the transform coefficients of the non-Gaussian component $`N`$. $`W`$ is some unknown distribution with vanishing first and third order moment (this constraint is only imposed to simplify the possible range of non-Gaussian distributions, and it will be representative of a Gaussian distribution modified by a few large values contributing to a tail). Without loss of generality, we assume the standard deviation of both $`z_i`$ and $`w_i`$ are $`1`$. Phrased in statistical terms, the problem of detecting the existence of a non-Gaussian component is equivalent to discriminating between the null hypothesis and the alternative hypothesis:
$$H_0:X_i=z_i,$$
$$H_1:X_i=\sqrt{1\lambda }z_i+\sqrt{\lambda }w_i,0\lambda 1.$$
$`N0`$ being equivalent to $`\lambda 0`$.
In order to obtain a quantitative estimation of the ability of the three statistics to detect the non-Gaussian component, we parametrize the tail probability of $`W`$ as follows.
$$\underset{x\mathrm{}}{lim}x^\alpha P\{|W|>x\}=C_\alpha ,C_\alpha \text{ is a constant. }$$
To model increasingly challenging situations as the number of observations increases, we calibrate $`\lambda `$ to decay with $`n`$ as:
$$\lambda =\lambda _n=n^r.$$
To constrain the detectability of such a non-Gaussian distribution superposed to a Gaussian one, we search the $`r`$ and $`\alpha `$ parameter space. We define the regions of this space that will be detectable under different statistical tests. It was shown in Jin et al. 2004 that there is a curve in the $`r`$-$`\alpha `$ plane that separates the detectable regions of parameters from the undetectable regions. That curve is given by:
$$r=\{\begin{array}{cc}2/\alpha ,\hfill & \alpha 8,\hfill \\ 1/4,\hfill & \alpha >8.\hfill \end{array}$$
In Figure 1, we compare the results of $`HC`$, $`Max`$ and $`\kappa `$. When it is possible to detect, $`HC`$ or $`Max`$ are better than $`\kappa `$ when $`\alpha 8`$. $`\kappa `$ is better than $`HC`$ or $`Max`$ when $`\alpha >8`$.
To conclude this section, we remark that the performance of $`HC`$, $`Max`$, and $`\kappa `$ depends on different situations of non-Gaussianity. Intuitively, $`Max`$ is designed to capture evidence of unusual behavior of the most extreme observations against the Gaussian assumption. The $`HC`$ statistic is able to capture unusual behavior of the most extreme observations, as well as unusually large amounts of moderately high observations. Thus $`HC`$ is better than $`Max`$, in general. However, when the evidence against the Gaussian assumption truly lies in the most extreme observations, $`HC`$ and $`Max`$ are almost equivalent. In contrast, $`\kappa `$ is designed to capture the evidence hidden in the 4th moment. Therefore this statistic depends on the bulk of the data, rather than on a few extreme observations or a small fraction of relatively large observations. Lastly, the $`HC`$ statistic offers an automatic way to find the area of the data accounting for the non-Gaussianity, while the $`Max`$ and $`\kappa `$ statistics do not have this capability.
## 3 WMAP First Year Data and Simulations
The data collected by the WMAP satellite during the first year of operation is available at the Legacy Archive for Microwave Background Data Analysis (LAMBDA) website<sup>2</sup><sup>2</sup>2http://lambda.gsfc.nasa.gov/. Following Komatsu et al. 2003, the analysis presented in this work is based on a data set obtained by the weighted combination of the released Foreground Cleaned Intensity Maps at bands Q,V and W. Each of these maps can be downloaded in the HEALpix<sup>3</sup><sup>3</sup>3HEALPix http://www.eso.org/science/healpix/ nested format at resolution $`nside=512`$ (total number of pixels being $`12\times nside^2`$). The co-added map is obtained by the following combination
$$T(i)=\frac{\underset{r=3}{\overset{10}{}}T_r(i)w_r(i)}{_{r=3}^{10}w_r(i)}.$$
The temperature at pixel $`i`$, $`T(i)`$ results from the ratio of the weighted sum of temperatures at pixel $`i`$ at each radiometer divided by the sum of the weights of each radiometer at pixel $`i`$. The radiometers Q1, Q2, V1, V2, W1, W2, W3 and W4 are sequentially numbered from 3 to 10. The weights at each pixel, for each radiometer $`w_r(i)`$, are the ratio of the number of observations $`N_r(i)`$ divided by the square of the receiver noise dispersion $`\sigma _{o,r}`$.This results in a co-added map at resolution $`nside=512`$. This map is downgraded to resolution $`nside=256`$ before the analyses are performed. In order to remove Galactic and point source emission, we applied the most conservative mask provided by the WMAP team, the Kp0 mask. This follows the procedure first presented by Vielva et al. 2004.
In order to detect any possible deviations from the predictions of the Standard Inflationary model we compared the values of several statistics in the WMAP data set described above, with those obtained from 5000 Monte Carlo simulations. The temperature at a certain pixel $`i`$ (pointing towards a direction characterized by polar angles $`\theta _i`$ and $`\varphi _i`$), can be expressed as an expansion in spherical harmonics $`Y_{lm}(\theta _i,\varphi _i)`$
$$T(\theta _i,\varphi _i)=\underset{l,m}{}a_{lm}Y_{lm}(\theta _i,\varphi _i)$$
The Monte Carlo simulations were performed assuming a Gaussian distribution $`N(0,C_l)`$ for the spherical harmonic coefficients $`a_{lm}`$ where $`C_l`$ represents the power spectrum that best fits the WMAP, CBI and ACBAR CMB data, plus the 2dF and Lyman-alpha data. Beam transfer functions as well as number of observations and receiver noise dispersion were taken into account when simulating data taken by each of the receivers (all of these, as well as the best fit power spectrum are provided by the WMAP team in the LAMBDA website).
The analysis was carried out in both real space and wavelet space. We convolved the WMAP data set and the Monte Carlo simulations with the Spherical Mexican Hat Wavelet (SMHW) at fifteen scales, following the same procedure presented in Vielva et al. 2004 and Cruz et al. 2005. The scales, numbered from 1 to 15, correspond to $`13.74`$, $`25.0`$, $`50.0`$, $`75.0`$, $`100.0`$, $`150.0`$, $`200.0`$, $`250.0`$, $`300.0`$, $`400.0`$, $`500.0`$, $`600.0`$, $`750.0`$, $`900.0`$, and $`1050.0`$ arcmin. In order to avoid border effects caused by the mask included in the WMAP data set and the simulations, analyses are performed outside extended masks defined at each scale. Given a scale, the extended mask is the Kp0 mask plus pixels near the Galactic plane that are within $`2.5`$ times the scale. All of these extended masks were presented in Figure 2 of Vielva et al. 2004. As discussed in the next section, deviations are detected in wavelet space. This shows once again the value of wavelets to provide a space in which certain features that might be causing deviations from the predictions of the Standard Inflationary model are enhanced.
## 4 ANALYSIS
Based on 5,000 simulations, we calculated the $`68\%`$, $`95\%`$, and $`99\%`$ confidence regions for each of the 4 statistics ($`\kappa `$, $`Max`$, and the two HC tests) at each of the 15 wavelet scales used. We used two-sided confidence regions for $`\kappa `$, as it is symmetric about 0 under the null hypothesis that the data is Gaussian. $`Max`$ and $`HC`$ statistics are defined so that the larger the statistic, the stronger the evidence against the null hypothesis. Therefore, we used one-sided confidence regions for these statistics.
Figure $`2`$ shows the results obtained based on $`\kappa `$. Dots denote the value of this statistic for the WMAP data set. Bands outlined by dashed, dotted-dashed, and solid lines correspond to the $`68\%`$, $`95\%`$, and $`99\%`$ confidence regions respectively (symbols and lines will represent the same in the figures included in this paper unless a comment is added). Clearly, the results agree with those presented in Vielva et al. 2004. Figures $`3`$ and $`4`$ show the new results based on the $`Max`$ and $`HC`$ statistical tests. The WMAP data appears outside the $`99\%`$ confidence level obtained based on the predictions of the Standard Inflationary model at scale 9 (300 arcmin) for these three tests. In particular, $`99.46\%`$ of the 5000 simulations have $`Max`$ and $`HC/HC^+`$ values below the one obtained from the WMAP data set. Therefore in this particular case these statistics are as competitive as the kurtosis in detecting deviations from the assumed null hypothesis.
### 4.1 Location of the outlying pixels provided by the Higher Criticism statistic
As discussed above, the $`HC`$ statistical tests introduced in this paper are able to detect deviations from the Standard Inflationary model in the WMAP data at scales around 300 arcmin. Moreover, the power of this new test resides in providing a direct way to determine which pixels in the WMAP data set are generating these deviations. By comparing the values of the $`HC`$ test at each pixel, with the value that defines the $`99\%`$ confidence limit at scale 9, we were able to extract a total of $`490`$ pixels that are causing the detected deviation. Figure 5 shows the individual pixel values of the $`HC`$ statistic for the WMAP data set at scale 9. The selected pixels above the $`99\%`$ limit are plotted in the map shown in Figure 6. As can be seen, these pixels define a ring centered at position $`(l209^{},b57^{})`$. It is important to note that the correlations introduced by the convolution with the wavelet need to be taken into account in order to properly interpret this result. Some pixels within the ring could also initially (in the map in real space) deviate from Gaussianity.
### 4.2 Effects of the subtraction of pixels at the source of the detection
Removal of the pixels in the ring from the analysed WMAP data resulted in a set of data compatible with the predictions from the Standard Inflationary model (see Figure 7). The deviation observed in the $`\kappa `$ statistic at the wavelet scale of $`5`$ degrees, in this paper as well as in previous papers by Vielva et al. 2004, Cruz et al. 2005, Mukherjee & Wang 2004, McEwen et al. 2004, disappeared as the pixels in the ring $`(l209^{},b57^{})`$ were removed. These pixels are part of the cold spot pointed out by Vielva et al. 2004 and Cruz et al. 2005 as being the source of the observed deviations. The $`Max`$ and the $`HC`$ values decrease as well (after removal of the pixels in the ring) fitting now within the $`68\%`$ confidence region. The Higher Criticism statistic is useful by offering an automatic detection of the non-Gaussian portions of the data. This characteristic is not easily available for other tests, such as Kurtosis.
## 5 Conclusions and Discussion
This paper presents an analysis of the compatibility of the distribution of the CMB temperature fluctuations observed by WMAP with the predictions from Standard Inflation. The analysis is based on the recently developed $`HC`$ statistic. This statistic is especially sensitive to two types of deviations from Gaussianity: an unusually large amount of moderately significant values or a small fraction of unusually extreme values. In both of these cases the HC statistic is optimal (Donoho & Jin 2004, Jin et al. 2004). Moreover, the definition of the statistic naturally suggests a direct way to locate the possible sources of non-Gaussianity. Even if the wavelet transform is a linear process, the original distribution is slightly distorted as the scale of the wavelet, and the area covered by the mask, increase. We are working on improving the power of the $`HC`$ statistic to be optimal in the detection of deviations from Gaussianity when applied in wavelet space.
We compared the performance of the $`HC`$ statistic with that of $`\kappa `$ and $`Max`$. The comparison was based on the analysis of the first year WMAP data, in real space and in wavelet space. Distributions of the three statistics were built on 5000 simulations assuming the Standard Inflationary model predictions as well as the constraints coming from the WMAP observations. The three statistics provided comparable results, pointing to the presence of deviations at a wavelet scale of $`5`$ degrees. We made use of the $`HC`$ statistic to automatically identify the pixels in the WMAP data set that were causing such deviations. A ring centered at $`(l209^{},b57^{})`$ containing $`490`$ pixels was shown to be the cause behind the detected deviations. Removal of this set of pixels from the WMAP data made the data compatible with the predictions from the Standard Inflationary model. One should be cautious when interpreting this result. The detection is achieved in wavelet space. Even if only the pixels in the ring appear as outliers of the HC distribution at wavelet scale of 5 degrees, some other pixels whithin the ring could be involved in the observed deviation. Convolution with the wavelet introduces correlations between pixels that need to be properly taken into account in the definition of the statistic to allow a correct interpretation. It is important to note that the pixels selected as the source of the observed deviations are part of the cold spot pointed out by Vielva et al. 2004 and Cruz et al. 2005. A careful study of the possible influence of foregrounds, noise, and beam distortion and assumption of a certain power spectrum was carried out in those two papers.
To conclude we would like to remark on the power of the $`HC`$ statistic for detecting and locating possible sources of non-Gaussianity. In particular, regarding analysis of the CMB, the $`HC`$ statistic can be useful at several steps in the data processing and final analysis, from the study of distortions caused by systematic effects in the time-ordered data to the study of compatibility of the statistical distribution of the observed temperature fluctuations with predictions from theories of the very early universe.
## Acknowledgements
We thank P. Vielva for providing the extended masks used at different wavelet scales. We would like to thank Douglas Crabill for helping with some of the figures. We acknowledge the use of the Legacy Archive for Microwave Background Data Analysis (LAMBDA). Support for LAMBDA is provided by the NASA Office of Space Science. Some of the results in this paper have been derived using the HEALPix (Górski, Hivon, and Wandelt 1999) package. L.C. and A.T. acknowledge the Indiana Space Grant Consortium for financial support. A.T. acknowledges the Spira Fellowship (Summer 2004 and Summer 2005). |
warning/0507/cond-mat0507565.html | ar5iv | text | # Dissipative Effects in the Electronic Transport through DNA Molecular Wires
## I Introduction
The idea that conduction pathways in DNA molecules may be built up as a result of the hybridization of the $`\pi `$ orbital stack along consecutive base pairs can be traced back to the 1960’s.Eley and Spivey (1962) It was not, however, till recently that a revival of interest on DNA as a potential conductor occurred. This was mainly triggered by the observation of long-range electron transfer between intercalated donor and acceptor centers in DNA molecules in solution.Murphy et al. (1993) Subsequent experimental results Treadway et al. (2002); Meggers et al. (1998a, b); Lewis et al. (1999); Brun and Harriman (1994); Kelley and Barton (1999); Grozema et al. (2000) were controversial as they showed different functional dependences of electron transfer rates on the donor-acceptor separation. Thus, strong exponential fall-offMeggers et al. (1998a, b) typical for superexchange mediated transfer as well as a weak, algebraic dependence Treadway et al. (2002); Kelley and Barton (1999) characteristic of sequential hopping processes were reported. Meanwhile, theoretical work has led to an emerging picture where different mechanisms may coexist depending on base-pair sequence and energetics.Schuster (2004); Jortner et al. (1998)
In parallel to these developments in the chemical physics community, DC transport experiments on $`\lambda `$-DNA as well as on poly(dG)-poly(dC) and poly(dA)-poly(dT) molecules between metal electrodes have been performed.Yoo et al. (2001); Porath et al. (2000); Braun et al. (1998); Kasumov et al. (2001); Storm et al. (2001); Tran et al. (2000); Porath et al. (2004); Endres et al. (2004) Several fundamental difficulties have to be surmounted in this kind of experiments: (i) how to create good contacts to the metal electrodes, (ii) how to control charge injection into the molecule, (iii) single molecule vs. bundles of molecules and (iv) dry vs. aqueous environments, among others. Consequently, sample preparation and the specific experimental conditions turn out to be very critical for DNA transport measurements. Thus, experiments have yielded contradictory results as to the conduction properties of DNA and are rather difficult to analyze. DNA has been characterized as a pure insulator, Braun et al. (1998); Storm et al. (2001) as a wide-band gap semiconductor,Porath et al. (2000) and as a metallic system.Yoo et al. (2001); Xu et al. (2004) Especially interesting are recent transport measurements on single poly(dG)-poly(dC) oligomers in aqueous solution, which displayed metallic-like $`I`$-$`V`$ characteristics and an algebraic behavior in the length dependence of the conductance.Xu et al. (2004)
Notwithstanding this variety of results and the problems related to the experimental set-up, the possibility of using DNA in molecular electronics is extremely attractive since it would open a vast range of potential applications because of its self-assembling and recognition properties.Braun et al. (1998) Alternatively, DNA can also be used as a template in molecular electronic devices.Keren et al. (2003); Hazani et al. (2004); Mertig et al. (1999)
From a theoretical point of view, the knowledge of the electronic structure of the base-pairs, the sugar/phosphate mantle and their mutual interactions is required in order to clarify the transport processes that may be effective in DNA. First principle approaches are the most suitable tools for this goal. However, the huge complexity of this molecule makes ab initio calculations still very demanding, so that only few investigations have been performed, mainly in well-stacked periodic structures. Calzolari et al. (2002); Felice et al. (2002); Barnett et al. (2001); Gervasio et al. (2002); Artacho et al. (2003); de Pablo et al. (2000); Alexandre et al. (2003); Lewis et al. (1997); Davies and Inglesfield (2004); Starikov (2003); Wang et al. (2004) To complicate this picture, environmental effects such as the presence of water molecules and counterions which stabilize the molecular structure make ab initio calculations even more challenging. Barnett et al. (2001); Gervasio et al. (2002) Hence, Hamiltonian modelsPalmero et al. (2004); Komineas et al. (2002); Hermon et al. (1998); Yi (2003); Li and Yan (2001); Cuniberti et al. (2002); Ao et al. (1996); Zhang and Ulloa (2004a); Hänggi et al. (2005); Gutierrez et al. (2004) that isolate single factors affecting electron transport are still playing a significant role and can help to shed more light onto the above issues as well as guide first principle investigations.
Recently, Cuniberti et al. Cuniberti et al. (2002) proposed a minimal model Hamiltonian to explain the semiconducting behavior previously observed by Porath et al. Porath et al. (2000) in suspended short (up to 30 base-pairs) poly(dG)-poly(dC) molecules. Remarkably enough, this experiment was performed on single molecules, in contrast to most transport experiments involving bundles of molecules. Molecular systems like Poly(dG)-poly(dC) (or Poly(dA)-poly(dT)) are very attractive from a theoretical standpoint since, being periodic, band-like transport as a result of $`\pi `$-orbitals hybridization may be more efficient than in its strongly disordered counterparts, e.g. in $`\lambda `$-DNA. The above modelCuniberti et al. (2002) mimics the electronic structure of the complex poly(dG)-poly(dC)-backbone system by a tight-binding chain to which side chains are attached. Electrons can hop along the central chain but not along the side chains. As a result a gap in the electronic spectrum opens. The gap is obviously temperature independent and the transmission near the Fermi level would show a strong exponential dependence due to the absence of electronic states to support transport.
An immediate issue that arises is how stable this electronic structure, i.e. two electronic bands separated by a gap, is against the influence of several factors which are known to play an important role in controlling charge propagation in DNA molecules, such as static and dynamic disorderBerlin et al. (2000); Grozema et al. (2002); Roche (2003); Yamada (2004); Zhang and Ulloa (2004b); Unge and Stafstrom (2003); Zhu et al. (2004) and environment.Barnett et al. (2001); Gervasio et al. (2002); Li and Yan (2001) In particular, the environment can act as a source of decoherence for a propagating electron (or hole),Li and Yan (2001) it can induce structural fluctuations that support or restrict charge motion,Barnett et al. (2001) or it can introduce additional electronic states within the fundamental gap.Gervasio et al. (2002); Endres et al. (2004) As it has been demonstrated experimentally, a modification of the humidity causes variations of orders of magnitude in the conductivity of DNA.Jo et al. (2003); Heim et al. (2004) Moreover, the recent single-molecule experiments of Xu et al. Xu et al. (2004) suggest that the environment may strongly modified the low-bias transport properties of DNA oligomers.
In this paper we elaborate on the role played by the environment by addressing signatures of the bath in the electronic transmission spectrum of the DNA wire in different coupling regimes: the mean-field approximation as well as weak-coupling and strong-coupling limits. Anticipating some of our results, we find that the semiconducting gap closes on the mean-field level as a result of thermal fluctuations. In the weak-coupling limit, however, the gap opens with increasing temperature. In both cases the electronic gap is an “intrinsic” property of the system. On the contrary, a bath-induced pseudo-gap is formed in the strong coupling limit, i.e. an energy region with a low (but finite) density of electronic states. We have further found in this regime that the transmission at the Fermi level exponentially decreases with the wire length $`L`$, $`t(E_\text{F})\text{e}^{\gamma L}`$. The decay rate $`\gamma `$ is however rather small $`0.2`$ Å<sup>-1</sup>. This together with a noticeable dependence of $`\gamma `$ on the electron-bath coupling clearly indicates that incoherent pathways do appreciably contribute to charge transport in the strong coupling limit.
In the next section we introduce the model Hamiltonian and derive the corresponding Green functions which are required to calculate the linear conductance. In section III different approximation schemes associated with different coupling regimes to the bath are discussed. The influence of structural disorder on our results is also presented. Finally, our summary follows in section IV.
## II Hamiltonian Model
Along the lines of Refs. Cuniberti et al., 2002, we represent the DNA molecular wire containing $`N`$ base pairs by the following nearest-neighbour tight-binding Hamiltonian (see Fig. 1):
$`_{\text{el}}`$ $`=`$ $`ϵ_b{\displaystyle \underset{j}{}}b_j^{}b_jt_{||}{\displaystyle \underset{j}{}}\left[b_j^{}b_{j+1}+\text{H.c.}\right]`$ (1)
$`+`$ $`ϵ{\displaystyle \underset{j}{}}c_j^{}c_j`$
$``$ $`t_{}{\displaystyle \underset{j}{}}\left[b_j^{}c_j+\text{H.c.}\right]`$
$`=`$ $`_\text{C}+_\text{b}+_{\text{C-c}}.`$
Hereby $`_\text{C}`$ and $`_\text{b}`$ are the Hamiltonians of the central and side chains, respectively, and $`_{\text{C-b}}`$ is the coupling between them. $`t_{||}`$ and $`t_{}`$ are hopping integrals along the central chain and between the backbone sites and the central chain, respectively. If not stated otherwise, the on-site energies will be later set equal to zero. The $`_\text{C}`$ Hamiltonian can be considered as effectively modeling one of the frontier orbitals of the poly(dG)-poly(dC) system, e.g. the highest-occupied molecular orbital, which is localized on the guanine bases.Artacho et al. (2003); Gervasio et al. (2002) The side chain induces then a perturbation of the $`\pi `$-stack leading to the opening of a temperature independent semiconducting gap in the electronic spectrum, the gap being proportional to the transversal hopping integral $`t_{}`$Cuniberti et al. (2002) Since this model shows electron-hole symmetry, two electronic manifolds containing $`N`$ states each, are symmetrically situated around the Fermi level, which is taken as the zero of energy.
We focus here on the influence of the environment on the electronic structure and consequently on the transport properties of the model described by $`_{\text{el}}`$. As it has been demonstrated in the past years, correlated fluctuations of hydrated counterions strongly influence electron(hole) motion along the double-helix.Barnett et al. (2001); Endres et al. (2004) Recent Raman and neutron scattering experiments on lysozyme have shown that the protein dynamics follows the solvent dynamics over a broad temperature range. Especially, conformational changes, low-energy vibrational excitations and the corresponding temperature dependences turned out to be very sensitive to the solvents dynamics. Caliscan et al. (2004) We consider the vibrational degrees of freedom of counterions and hydration shells in DNA as a dynamical bath able to act as a dissipative environment. In this model Hamiltonian approach, we do not consider specific features of the environment but represent it by a phonon bath of $`M`$ harmonic oscillators. We further make the assumption that the bath is only directly affecting the side chain whereas the central chain is well screened by the latter. Then, the extended Hamiltonian becomes:
$`_\text{W}`$ $`=`$ $`_{\text{el}}+{\displaystyle \underset{\alpha }{}}\mathrm{\Omega }_\alpha B_\alpha ^{}B_\alpha +{\displaystyle \underset{\alpha ,j}{}}\lambda _\alpha c_j^{}c_j(B_\alpha +B_\alpha ^{})`$ (2)
$`=`$ $`_{\text{el}}+_\text{B}+_{\text{c-B}},`$
where $`_\text{B}`$ and $`_{\text{c-B}}`$ are the phonon bath Hamiltonian and the backbone-bath interaction, respectively. $`B_\alpha `$ is a bath phonon operator and $`\lambda _\alpha `$ denotes the electron-phonon coupling. Note that we assume a local coupling of the bath modes to the electronic density at the side chain. Later on, the thermodynamic limit ($`M\mathrm{})`$ in the bath degrees of freedom will be carried out and the corresponding bath spectral density introduced, so that at this stage we do not need to further specify the set of bath frequencies $`\mathrm{\Omega }_\alpha `$ and coupling constants $`\lambda _\alpha `$.
Finally, we include the coupling of the molecular wire to semiinfinite left (L) and right (R) electrodes:
$``$ $`=`$ $`_\text{W}+{\displaystyle \underset{𝐤\text{L,R},\sigma }{}}ϵ_{𝐤\sigma }d_{𝐤\sigma }^{}d_{𝐤\sigma }`$ (3)
$`+{\displaystyle \underset{𝐤\text{L},\sigma }{}}(V_{𝐤,1}d_{𝐤\sigma }^{}b_1+\text{H.c.})`$
$`+{\displaystyle \underset{𝐤R,\sigma }{}}(V_{𝐤,N}d_{𝐤\sigma }^{}b_N+\text{H.c.})`$
$`=`$ $`_\text{W}+_{\text{L/R}}+_{\text{L-C}}+_{\text{R-C}}`$
The Hamiltonian of Eq. (3) is the starting point of our investigation. Performing the Lang-FirsovMahan (2000) unitary transformation $`\overline{}=\text{e}^S\text{e}^S`$ with the generator $`S=_{\alpha ,j}(\lambda _\alpha /\mathrm{\Omega }_\alpha )c_j^{}c_j(B_\alpha B_\alpha ^{})`$ and $`S^{}=S`$, the linear coupling to the bath can be eliminated. In the resulting effective Hamiltonian only the backbone part is modified since the central chain operators $`b_{\mathrm{}}`$ as well as the leads’ operators $`d_{𝐤\sigma }`$ are invariant with respect to the above transformation. The new Hamiltonian reads:
$`\overline{}`$ $`=`$ $`_\text{C}+_{\text{L/R}}+_\text{B}+_{\text{L/R-C}}`$ (4)
$`+`$ $`(ϵ\mathrm{\Delta }){\displaystyle \underset{j}{}}c_j^{}c_jt_{}{\displaystyle \underset{j}{}}\left[b_j^{}c_j𝒳+\text{H.c.}\right]`$
$`𝒳`$ $`=`$ $`\mathrm{exp}\left[{\displaystyle \underset{\alpha }{}}{\displaystyle \frac{\lambda _\alpha }{\mathrm{\Omega }_\alpha }}(B_\alpha B_\alpha ^{})\right],\mathrm{\Delta }={\displaystyle \underset{\alpha }{}}{\displaystyle \frac{\lambda _\alpha ^2}{\mathrm{\Omega }_\alpha }}.`$
Let’s define two kinds of retarded thermal Green functions related to the central chain $`G_j\mathrm{}(t)`$ and to the backbones $`P_j\mathrm{}(t)`$, respectively ($`\mathrm{}=1`$):
$`G_j\mathrm{}(t)`$ $`=`$ $`i\mathrm{\Theta }(t)[b_j(t),b_{\mathrm{}}^{}(0)]_+,`$ (5)
$`P_j\mathrm{}(t)`$ $`=`$ $`i\mathrm{\Theta }(t)[c_j(t)𝒳(t),c_{\mathrm{}}^{}(0)𝒳^{}(0)]_+,`$
where $`\mathrm{\Theta }`$ is the Heaviside function and the average is taken w.r.t. $`\overline{}`$. With the above definitions and using the equation of motion technique (see Appendix A) we arrive to an expression for the Fourier transform of the central chain Green function which reads, to lowest-order in $`t_{}`$:
$`𝐆^1(E)`$ $`=`$ $`𝐆_0^1(E)t_{}^2𝐏(E)`$ (6)
$`𝐆_0^1(E)`$ $`=`$ $`E\mathrm{𝟏}_\text{C}\mathrm{\Sigma }_\text{L}(E)\mathrm{\Sigma }_\text{R}(E).`$
In this equation $`𝐆_0(E)`$ is the Green function of a chain without backbones and connected to the left and right electrodes. The influence of the latter is comprised in the complex self-energy functions $`\mathrm{\Sigma }_{\text{L/R}}(E)`$.Datta (1995) The polaronic Green function $`𝐏(E)`$ is explicitly given by:
$`P_\mathrm{}j(E)`$ $`=`$ $`\text{i}\delta _\mathrm{}j{\displaystyle _0^{\mathrm{}}}dt\text{e}^{\text{i}(E+\text{i}\mathrm{\hspace{0.17em}0}^+)t}\text{e}^{\text{i}(ϵ\mathrm{\Delta })t}`$ (7)
$`\times \left[(1f_𝖼)\text{e}^{\mathrm{\Phi }(t)}+f_𝖼\text{e}^{\mathrm{\Phi }(t)}\right]`$
with $`\text{e}^{\mathrm{\Phi }(t)}=𝒳(t)𝒳^{}(0)_\text{B}`$ being a dynamical bath correlation function. The average $`_\text{B}`$ is performed over the bath degrees of freedom. Working to lowest order in $`t_{}`$ allows to use a zero-order Green function for the side chain in Eq. (7), i.e. $`G_{0,\mathrm{}j}^𝖼(t)\delta _\mathrm{}j\text{e}^{\text{i}(ϵ\mathrm{\Delta })t}`$. $`f_𝖼`$ is the Fermi function at the backbone sites. In what follows we consider the case of empty sites by setting $`f_𝖼=0`$. Note that $`𝐏`$ is a diagonal matrix, i.e. it only modifies the on-site energies in the Hamiltonian.
In order to get closed expressions for the bath thermal averages it is appropriate to introduce a bath spectral densityWeiss (1999) defined by :
$`J(\omega )={\displaystyle \underset{\alpha }{}}\lambda _\alpha ^2\delta (\omega \mathrm{\Omega }_\alpha )=J_0({\displaystyle \frac{\omega }{\omega _\text{c}}})^s\text{e}^{\omega /\omega _\text{c}}\mathrm{\Theta }(\omega ),`$ (8)
where $`\omega _\text{c}`$ is a cut-off frequency related to the bath memory time $`\tau _\text{c}\omega _\text{c}^1`$. It is easy to show that the limit $`\omega _\text{c}\mathrm{}`$ corresponds to a Markovian bath, i.e. $`J(t)J_0\delta (t)`$. Using this Ansatz, $`\mathrm{\Phi }(t)`$ can be written in the usual way:Weiss (1999)
$`\mathrm{\Phi }(t)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}d\omega {\displaystyle \frac{J(\omega )}{\omega ^2}}\left[1\text{e}^{\text{i}\omega t}+2{\displaystyle \frac{1\mathrm{cos}\omega t}{\text{e}^{\beta \omega }1}}\right].`$ (9)
Although the integral can be performed analyticallyWeiss (1999), we consider $`\mathrm{\Phi }(t)`$ in some limiting cases where it is easier to work directly with Eq. (9).
In the transport calculations, we limit ourselves to treat the low voltage regime, thus neglecting non-equilibrium effects as well as the inelastic part of the total current. As a result, one can still define a linear conductance $`g`$ as follows: Imry et al. (2004)
$`g(E)`$ $`=`$ $`{\displaystyle \frac{2e^2}{h}}{\displaystyle 𝑑E\left(\frac{f}{E}\right)t(E)},`$ (10)
$`t(E)`$ $`=`$ $`\text{Tr}\left\{𝚫_\mathrm{L}𝐆𝚫_\mathrm{R}𝐆^{}\right\},`$
where $`𝚫_{\mathrm{L},\mathrm{R}}=\mathrm{i}\left(𝚺_{\mathrm{L},\mathrm{R}}𝚺_{\mathrm{L},\mathrm{R}}^{}\right)`$ are spectral densities of the leads. Although the foregoing expression is similar to Landauer’s formula, we stress that the influence of the phonon bath does implicitly appear via the Green function $`𝐆`$. Hence, both coherent and incoherent pathways for charge transport mediated by phonon processes are included in Eq. (10). We concentrate our discussion on the temperature and length dependence of $`t(E)`$. In what follows we always plot $`t(E)`$ rather than $`g`$ to filter out temperature effects arising from the derivative of the Fermi function in Eq. (10). For completeness the current as given by $`I(V)=(2e/h)dE(f(EeV/2)f(E+eV/2))t(E)`$ is also shown. We remark however, that this expression neglects non-equilibrium effects, which are beyond the scope of this investigation.
## III Limiting cases
We use now the results of the foregoing section to discuss the electronic transport properties of our model in some limiting cases for which analytic expressions can be derived. In all cases, we use the wide-band limit in the electrode selfenergies, i.e. $`\mathrm{\Sigma }_{\text{L},\mathrm{}j}(E)=\text{i}\mathrm{\Gamma }_\text{L}\delta _1\mathrm{}\delta _{1j}`$ and $`\mathrm{\Sigma }_{\text{R},\mathrm{}j}(E)=\text{i}\mathrm{\Gamma }_\text{R}\delta _N\mathrm{}\delta _{Nj}`$. We discuss the mean-field approximation and the weak-coupling regime in the electron-bath interaction as well as the strong-coupling limit. Farther, the cases of ohmic ($`s=1`$) and superohmic ($`s=3`$) spectral densities are treated.
### III.1 Mean-field approximation (MFA)
Within the mean-field approximation bath fluctuations contained in $`P(E)`$ are neglected. The MFA can be introduced by writing the phonon operator $`𝒳`$ as $`𝒳_\mathrm{B}+\delta 𝒳`$ in $`_{\text{C-c}}`$ in Eq. (4), i.e. $`_{\text{C-c}}^{\text{MF}}=t_{}_j\left[b_j^{}c_j𝒳_\text{B}+\text{H.c.}\right]+O(\delta 𝒳)`$. As a result a real, static and temperature dependent term in Eq. (6) is found:
$`𝐆^1(E)=𝐆_0^1(E)t_{}^2{\displaystyle \frac{|𝒳_\text{B}|^2}{Eϵ+\mathrm{\Delta }+\text{i}\mathrm{\hspace{0.17em}0}^+}}\mathrm{𝟏},`$ (11)
where $`\left|𝒳_\text{B}\right|^2=\text{e}^{2\kappa (T)}`$ and $`\kappa (T)`$ is given by:
$`\kappa (T)={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\mathrm{d}\omega }{\omega ^2}}J(\omega )\mathrm{coth}{\displaystyle \frac{\omega }{2k_\text{B}T}}.`$ (12)
The effect of the MF term is thus to scale the bare transversal hopping $`t_{}`$ by the exponential temperature dependent factor $`\text{e}^{\kappa (T)}`$.
In the case of an ohmic bath, $`s=1`$, the integrand in $`\kappa (T)`$ scales as $`1/\omega ^p,p=1,2`$ and has thus a logarithmic divergence at the lower integration limit, see Eqs. (8) and (12). Thus, the MF contribution would vanish. In other words, no gap would exist on this approximation level.
In the superohmic case ($`s=3`$) all integrals are regular. One obtains $`\mathrm{\Delta }=d\omega \omega ^1J(\omega )=\mathrm{\Gamma }(s1)J_0=2J_0`$, with $`\mathrm{\Gamma }(s)`$ being the Gamma function and $`\kappa (T)`$ reads:
$`\kappa (T)=`$ (13)
$`{\displaystyle \frac{2J_0}{\omega _\text{c}}}\left[\mathrm{\hspace{0.17em}2}\left({\displaystyle \frac{k_\text{B}T}{\omega _\text{c}}}\right)^2\zeta _\text{H}(2,{\displaystyle \frac{k_\text{B}T}{\omega _\text{c}}})1\right].`$
$`\zeta _\text{H}(s,z)=_{n=0}^{\mathrm{}}(n+z)^s`$ is the Hurwitz $`\zeta `$-function, a generalization of the Riemann $`\zeta `$-function.Gradshteyn and Ryzhik (2000)
It follows from Eq. (13) that $`\kappa (T)`$ behaves like a constant for low temperatures ($`k_\text{B}T/\omega _\text{c}<1`$), $`\kappa (T)J_0/\omega _\text{c}`$, while it scales linear with $`T`$ in the high-temperature limit ($`k_\text{B}T/\omega _\text{c}>1`$), $`\kappa (T)J_0/\omega _\text{c}(1+2k_\text{B}T/\omega _\text{c})`$.
If $`J_0`$ vanishes, $`\mathrm{\Delta }`$ is zero and $`𝒳_\text{B}=1`$. Thus we recover the original model of Ref. Cuniberti et al., 2002 which has a gap proportional to $`t_{}`$. For $`J_00`$ and at zero temperature the hopping integral is roughly reduced to $`t_{}\text{e}^{\frac{J_0}{\omega _\text{c}}}`$ which is similar to the renormalization of the hopping in Holstein’s polaron model,Holstein (1959) though here it is $`t_{}`$ rather than $`t_{||}`$ the term that is rescaled. At high temperatures $`t_{}`$ is further reduced ($`\kappa (T)T`$) so that the gap in the electronic spectrum finally collapses and the system becomes metallic, see Fig. 2. An appreciable temperature dependence can only be observed in the limit $`J_0/\omega _\text{c}<1`$; otherwise the gap would collapse already at zero temperature due to the exponential dependence on $`J_0`$. We further remark that the MFA may be only valid in the regime $`J_0/\omega _\text{c}<1`$, $`k_\text{B}T/\omega _\text{c}1`$, otherwise multiphonon processes in the bath, which are not considered at this stage, become increasingly relevant.
### III.2 Beyond MF: weak-coupling limit
As a first step beyond the mean-field approach let’s first consider the weak-coupling limit in $`𝐏(E)`$. For $`J_0/\omega _\text{c}<1`$ and not too high temperatures ($`k_\text{B}T/\omega _\text{c}1`$) the main contribution to the integral in Eq. (7) comes from long times $`t\omega _\text{c}^1`$.
With the change of variables $`z=\omega t`$, $`\mathrm{\Phi }(t)`$ can be written as:
$`\mathrm{\Phi }(t)`$ $`=`$ $`J_0\omega _\text{c}^st^{1s}{\displaystyle _0^{\mathrm{}}}dzz^{s2}\text{e}^{\frac{z}{\omega _\text{c}t}}`$ (14)
$`\times \left(1\text{e}^{\text{i}z}+2{\displaystyle \frac{1\mathrm{cos}z}{\text{e}^{z\frac{\beta \omega _\text{c}}{\omega _\text{c}t}}1}}\right).`$
As far as $`\omega _\text{c}t\beta \omega _\text{c}`$ this can be simplified to:
$`\mathrm{\Phi }(t)`$ $``$ $`J_0\omega _\text{c}^st^{1s}{\displaystyle _0^{\mathrm{}}}dxz^{s2}\text{e}^{\frac{z}{\omega _\text{c}t}}`$ (15)
$`\times \left(1\text{e}^{\text{i}z}+2{\displaystyle \frac{\beta \omega _\text{c}}{\omega _\text{c}t}}{\displaystyle \frac{1\mathrm{cos}z}{z}}\right).`$
Since in the long-time limit the low-frequency bath modes are giving the most important contribution we may expect some qualitative differences in the ohmic and superohmic regimes. For $`s=1`$ we obtain $`\mathrm{\Phi }(t)\pi \frac{J_0}{\omega _\text{c}}\frac{k_\text{B}T}{\omega _\text{c}}(\omega _\text{c}t)`$ which leads to (using $`\mathrm{\Delta }(s=1)=J_0`$):
$`𝐆^1(E)`$ $`=`$ $`𝐆_0^1(E)t_{}^2{\displaystyle \frac{1}{E+J_0+\text{i}\pi \frac{J_0}{\omega _\text{c}}k_\text{B}T}}\mathrm{𝟏},`$ (16)
i.e. there is only a pure imaginary contribution from the bath. For the simple case of a single site coupled to a backbone one can easily see that the gap approximately scales as $`\sqrt{k_\text{B}T}`$; thus it grows with increasing temperature. This is shown in Fig. 3, where we also see that the intensity of the transmission resonances strongly goes down with increasing temperature. The gap enhancement is induced by the suppression of the transmission peaks of the frontier orbitals, i. e. those closest to the Fermi energy.
For $`s=3`$ and $`k_\text{B}T/\omega _\text{c}1`$, $`\mathrm{\Phi }(t)`$ takes a nearly temperature independent value proportional to $`J_0/\omega _\text{c}`$. As a result the gap is slightly reduced ($`t_{}t_{}\text{e}^{J_0/\omega _\text{c}}`$) but, because of the weak-coupling condition, the effect is rather small.
From this discussion we can conclude that in the weak-coupling limit ohmic dissipation in the bath induces an enhancement of the electronic gap while superohmic dissipation does not appreciably affect it. In the high-temperature limit $`k_\text{B}T/\omega _\text{c}>1`$ a short-time expansion can be performed which yields similar results to those of the strong-coupling limit (see next section),Ao et al. (1996) so that we do not need to discuss them here. Note farther that the gap obtained in the weak-coupling and mean-field limits is an “intrinsic” property of the electronic system; it is only quantitatively modified by the interaction with the bath degrees of freedom. We thus trivially expect a strong exponential dependence of $`t(E=E_\mathrm{F})`$ on the wire length, typical of virtual tunneling through a gap. Indeed, we find $`t(E=E_\mathrm{F})\mathrm{exp}(\beta L)`$ with $`\beta 23\AA ^1`$.
### III.3 Beyond MF: strong coupling limit (SCL)
In this section we discuss the strong-coupling regime, as defined by the condition $`J_0/\omega _\text{c}>1`$. This may be the regime to be found in presence of an aqueous environment, as recent theoretical estimations using the classical Onsager model for solvation processes have shown. Gilmore et al. (2005) In the SCL the main contribution to the time integral in Eq. (7) arises from short times. Hence a short-time expansion of $`\mathrm{\Phi }(t)`$ may already give reasonable results and it allows, additionally, to find an analytical expression for $`𝐏(E)`$. At $`t\omega _\text{c}^1`$ we find,
$`\mathrm{\Phi }(t)`$ $``$ $`\text{i}\mathrm{\Delta }t+(\omega _\text{c}t)^2\kappa _0(T)`$ (17)
$`P_\mathrm{}j(E)`$ $`=`$ $`\text{i}\delta _\mathrm{}j{\displaystyle _0^{\mathrm{}}}dt\text{e}^{\text{i}(Eϵ+\text{i}\mathrm{\hspace{0.17em}0}^+)t}\text{e}^{(\omega _\text{c}t)^2\kappa _0(T)}`$
$`=`$ $`\text{i}\delta _\mathrm{}j{\displaystyle \frac{\sqrt{\pi }}{2}}{\displaystyle \frac{1}{\omega _\text{c}\sqrt{\kappa _0(T)}}}\mathrm{exp}\left({\displaystyle \frac{(Eϵ+\text{i}\mathrm{\hspace{0.17em}0}^+)^2}{4\omega _\text{c}^2\kappa _0(T)}}\right)`$
$`\times \left(1+\text{erf}\left[{\displaystyle \frac{\text{i}(Eϵ+\text{i}\mathrm{\hspace{0.17em}0}^+)}{2\omega _\text{c}\sqrt{\kappa _0(T)}}}\right]\right),`$
$`\kappa _0(T)`$ $`=`$ $`{\displaystyle \frac{1}{2\omega _\text{c}^2}}{\displaystyle _0^{\mathrm{}}}d\omega J(\omega )\mathrm{coth}{\displaystyle \frac{\omega }{2k_\text{B}T}}.`$
Before presenting the results for the electronic transmission, it is useful to first consider the dependence of the real and imaginary parts of $`𝐏(E)`$ on temperature and on the reduced coupling constant $`J_0/\omega _\text{c}`$. Both functions are shown in Fig. 4. We see that around the Fermi level at $`E=0`$ the real part is approximately linear, $`\mathrm{Re}P(E)E`$ while the imaginary part shows a Gaussian-like behavior. The imaginary part loses intensity and becomes broadened with increasing temperature or $`J_0`$, while the slope in the real part decreases when $`k_\text{B}T`$ or $`J_0`$ are increased.
If we neglect for the moment the imaginary part (the dissipative influence of the bath), we can understand the consequences of the real part being nonzero around the Fermi energy, i.e. in the gap region. The solutions of the non-linear equation $`\mathrm{det}|(Et_{}^2\mathrm{Re}P(E))\mathrm{𝟏}_\text{C}|=0`$ give the new poles of the Green function of the system in presence of the phonon bath. For comparison, the equation determining the eigenstates without the bath is simply $`\mathrm{det}|(Et_{}^2/E)\mathrm{𝟏}_\text{C}|=0`$. It is just the $`1/E`$ dependence near $`E=0`$ that induces the appearance of two electronic bands of states separated by a gap.Cuniberti et al. (2002) In our present study, however, $`\mathrm{Re}P(E0)`$ has no singular behavior and additional poles of the Green function may be expected to appear in the low-energy sector. This is indeed the case, as shown in Fig. 5 (upper panel). We find a third band of states around the Fermi energy, which we call a polaronic band because it results from the strong interaction between an electron and the bath modes. The intensity of this band as well as its band width strongly depend on temperature and on $`J_0`$. When $`k_\text{B}T`$ (or $`J_0`$) become large enough, these states spread out and eventually merge with the two other side bands. This would result in a transmission spectrum similar of a metallic system.
This picture is nevertheless not complete since the imaginary component of $`P(E)`$ has been neglected. Its inclusion leads to a dramatic modification of the spectrum, as shown in Fig. 5 (middle panel). We now only see two bands separated by a gap which basically resembles the semiconducting-type behavior of the original model. The origin of this gap or rather pseudo-gap (see below) is however quite different. It turns out that the imaginary part of $`P(E)`$, being peaked around $`E=0`$, strongly suppresses the transmission resonances belonging to the central band. Additionally, the frontier orbitals on the side bands, i.e. orbitals closest to the gap region, are also strongly damped, this effect becoming stronger with increasing temperature ($`\mathrm{Im}P(E)`$ broadens). This latter effect has some similarities with the previously discussed weak-coupling regime. Note, however, that the new electronic manifold around the Fermi energy does not appear in the weak-coupling regime. We further stress that the density of states around the Fermi level is not exactly zero (hence the term pseudo-gap); the states on the polaronic manifold, although strongly damped, contribute nevertheless with a finite, temperature dependent incoherent background to the transmission. As a result, with increasing temperature, a crossover from tunneling to activated behavior in the low-voltage region of the $`I`$-$`V`$ characteristics takes place, see Fig. 5 (lower panel). The slope in the $`I`$-$`V`$ plot becomes larger when $`t_{}`$ is reduced, since the side bands approach each other and the effect of $`\mathrm{Im}P(E)`$ is reinforced.
In Fig. 6 an Arrhenius plot of the transmission at the Fermi energy is shown for different strengths of the transversal hopping integral and the electron-bath coupling. After a nearly $`T`$-independent region, the transmission strongly grows up following approximately a $`\text{e}^{1/T}`$ law. Increasing the coupling to the phonon bath makes the suppression of the polaronic band around $`E=0`$ less effective ($`\mathrm{Im}P(E0)`$ decreases) so that the density of states around this energy becomes larger. Hence the absolute value of the transmission also increases. Similar $`T`$-dependences have been experimentally observed in poly(dG)-poly(dC)Yoo et al. (2001) as well as in $`\lambda `$-DNA.Tran et al. (2000) On the other side, increasing $`t_{}`$ leads to a reduction of the transmission at the Fermi level, since the energetic separation of the side bands increases with $`t_{}`$.
We have further investigated the length dependence of the transmission at the Fermi energy. This is a very important aspect that helps to identify the influence of different transport mechanisms.Jortner et al. (1998); Segal et al. (2000) The results are displayed in Fig. 7 for different values of the reduced coupling $`J_0/\omega _\text{c}`$. For a homogeneous chain (on-site energies are set to zero) an exponential dependence on the chain length $`t(E_\text{F})\text{e}^{\gamma L}`$ was found. In this expression $`L=Na_0`$, where $`N`$ is the number of sites on the molecular wire and $`a_03.4\mathrm{\AA }`$ is the average distance between consecutive base pairs. Note that the inverse decay lengths $`\gamma `$ are rather small $`0.10.3\mathrm{\AA }^1`$. An exponential dependence usually indicates virtual tunneling through a gap. Inverse decay lengths, as extracted e. g. from complex band structure calculations, Wang et al. (2004); Fagas et al. (2004) are however much larger that those obtained in the present investigation. So have recent DFT-based calculations found values of $`\gamma 1.5\mathrm{\AA }^1`$ for gap tunneling in dry Poly(dG)-Poly(dC) oligomers.Wang et al. (2004) With increasing bath coupling the exponential dependence farther weakens and eventually becomes algebraic $`t(E_\text{F})N^\alpha `$. The introduction of a tunnel barrier as realized e.g. through insertion of (AT)<sub>n</sub> groups, by shifting the on-site energies along a finite segment of the chain increases the inverse decay length $`\gamma `$ by a factor of 2, approximately. Obviously, this model cannot describe the crossover from superexchange mediated electron transfer (strong exponential behavior) to sequential hopping-mediated transport (algebraic dependence) as a function of the wire length $`N`$, as discussed in other works.Jortner et al. (1998); Segal et al. (2000) We guess that vibrational excitations inside the central chain, which renormalize the longitudinal hopping integral $`t_{||}`$, have to be included to get this non-monotonic transition.
From the previous discussion we may conclude that electron transport on the low-energy sector of the transmission spectrum is supported by the formation of polaronic states. Though strongly damped, these states manifest nonetheless with a finite density of states inside the bandgap.
It has been meanwhile demonstrated Roche (2003); Yamada (2004); Barnett et al. (2001); Grozema et al. (2002); Zhang and Ulloa (2004b); Unge and Stafstrom (2003); Zhu et al. (2004) that electron (or hole) motion in DNA is extremely sensitive to different kinds of disorder: static disorder (random base-pair sequences), structural fluctuations and inhomogeneities of the counterions distribution along the backbones. These factors may strongly distort the base pair stacking along the double helix and eventually affect the electronic transport properties. They deserve a separate study. However, as a test for the stability of our results we have randomly varied the on-site energies along the central chain by extracting them from a Gaussian distribution with variance $`\sigma _0`$. In this way we are simulating some kind of structural disorder induced, e.g. by thermal fluctuations inside the central chain. In Fig. 8 the cases of weak ($`\sigma _00.12t_{||}`$) and strong disorder ($`\sigma _0t_{||}`$) are shown. Two main features can be seen: (i) the transmission resonances on the side bands are strongly washed out and lose in intensity, and (ii) the pseudo-gap is slightly reduced with increasing disorder. However, the suppression of the central band due to $`\mathrm{Im}P(E)`$ and hence, the pseudo-gap formation is not affected by this kind of disorder. As soon as electronic states shift from the side bands into the region with nonzero $`\mathrm{Im}P(E)`$ they are strongly damped and thus the pseudo-gap structure of the spectrum is conserved. A similar effect of disorder is expected in the other coupling regimes to the bath degrees of freedom discussed above.
## IV Summary
Charge propagation in DNA molecules is extremely sensitive to disorder and environmental effects. We have focused in this paper on the influence of a dissipative environment on the electronic transport properties of a model Hamiltonian which mimics some basic features of the electronic structure of DNA oligomers. Although we have chosen Poly(dG)-Poly(dC) molecules as a reference point, we believe that our model is quite generic and may be useful for a large class of $`\pi `$-conjugated systems.
We have shown that a mean-field approximation cannot fully catch the action of a dissipative environment on charge transport, because it only gives a real, energy independent contribution. Indeed, while the mean-field approach leads to gap reduction with increasing temperature, bath fluctuations eventually lead to gap opening in the weak-coupling limit. We have further shown that a bath-induced pseudo-gap in the electronic spectrum can appear for strong electron-bath coupling giving a temperature-dependent background around the Fermi energy. As a result the system may show with increasing temperature a transition from a tunneling to an activated behavior in the low-bias region when coupled to an external dissipative bath. An Arrhenius-like temperature dependence of the transmission at the Fermi level and a rather weak exponential dependence on the wire length were additionally found, indicating a strong contribution of incoherent pathways of the charge carriers.
A natural extension of this investigation would be the inclusion of non-equilibrium effects at large bias and consequently of inelastic components of the current. This issue is although interesting from a formal point of view, since the Lang-Firsov transformation introduces polaronic rather than pure electronic propagators, see Eq. (5). For the former the appropriate Keldysh Green functions should be derived in order to deal with the non-equilibrium regime. This problem deserves a separate investigation which is now in progress.
## V Acknowledgments
We would like to thank M. Hartung and J. Keller for fruitful discussions. This work has been supported by the Volkswagen foundation and by the EU under contract IST-2001-38951.
## Appendix A Derivation of Eq. (6)
The equation of motion for the retarded Green function in Eq. (5) in the frequency representation reads:
$`EG_\mathrm{}j(E)=[b_j,b_{\mathrm{}}]_++(([b_j,\overline{}]|b_{\mathrm{}})).`$
Using it we get for the Hamiltonian of Eq. (3) :
$`{\displaystyle \underset{n}{}}\left[G_0^1(E)\right]_\mathrm{}nG_{nj}(E)=\delta _\mathrm{}jt_{}((c_{\mathrm{}}𝒳|b_j^{}))`$ (18)
$`\left[G_0^1(E)\right]_\mathrm{}n`$ $`=`$ $`(Eϵ_b)\delta _n\mathrm{}+t_{||}(\delta _{n,\mathrm{}+1}+\delta _{n,\mathrm{}1})`$
$`\mathrm{\Sigma }_\text{L}\delta _\mathrm{}1\delta _{n1}\mathrm{\Sigma }_\text{R}\delta _\mathrm{}N\delta _{nN}`$
$`\mathrm{\Sigma }_{\text{L(R)}}`$ $`=`$ $`{\displaystyle \underset{𝐤L(R)}{}}{\displaystyle \frac{|V_{𝐤,1(N)}|^2}{Eϵ_𝐤+\text{i}\mathrm{\hspace{0.17em}0}^+}}`$
Now, equations of motion from the “right” may be written for the Green function $`Z_\mathrm{}j^𝒳(E)=((c_{\mathrm{}}𝒳|b_j^{}))`$,
leading to :
$`{\displaystyle \underset{m}{}}Z_\mathrm{}m^𝒳(E)\left[G_0^1(E)\right]_{mj}`$ $`=`$ $`t_{}((c_{\mathrm{}}𝒳|c_j^{}𝒳^{}))`$ (19)
$`=`$ $`t_{}P_\mathrm{}j(E)`$
In the former equations we have neglected cross-terms of the form $`((c_{\mathrm{}}𝒳|c_j^{}))`$, since they will give contribution of $`O(t_{}^3)`$. Inserting Eq. (19) into Eq. (18) we arrive at the matrix equation:
$`𝐆(E)=𝐆_0(E)+𝐆_0(E)𝚺_\text{B}(E)𝐆_0(E),`$
which can be transformed into a Dyson-like equation when introducing the irreducible part $`𝚺_\text{B}(E)=𝚺_\text{B}^{\text{irr}}(E)+𝚺_\text{B}^{\text{irr}}(E)𝐆_0(E)𝚺_\text{B}^{\text{irr}}(E)+\mathrm{}`$:
$`𝐆(E)=𝐆_0(E)+𝐆_0(E)𝚺_\text{B}^{\text{irr}}(E)𝐆(E).`$ (20)
From Eq. (20) it immediately follows Eq. (6) with $`𝚺_\text{B}^{\text{irr}}(E)=t_{}^2𝐏(E)`$. We emphasize that these expressions are exact only to lowest-order in the transversal hopping $`t_{}`$. This approximation may be justified in the low-voltage limit we are dealing with. |
warning/0507/cond-mat0507203.html | ar5iv | text | # Dimer and trimer fluctuations in the 𝑠=1/2 transverse 𝑋𝑋 chain
## 1 Introduction
The theoretical and computational study of frequency-resolved quantum fluctuations and thermal fluctuations in many-body model systems is an important area of research for several reasons. Such fluctuations are observable directly or indirectly by a host of measuring techniques used in condensed matter physics and materials science. The shape of the spectrum (dispersions, bandwidths, gaps, soft modes, etc.), the spectral-weight distributions, and the singularity structure of the dynamic structure factors as measured or calculated for specific fluctuation operators yield detailed insights into the state of the material and reveal important clues about the susceptibility of the system to phase transitions with order parameters modeled after the fluctuation operators at hand .
An increasing number of exactly solvable quantum many-body model systems turn out to be relevant for situations where certain degrees of freedom of a material are kinematically constrained to one spatial dimension. Many such situations offer the most detailed comparisons in many-body dynamics of exact theoretical results with direct experimental observation .
Various properties of dynamic spin structure factors of quantum spin chain models are observable in quasi-one-dimensional magnetic insulators, for example, via magnetic neutron scattering . The properties of dynamic dimer and trimer structure factors, on the other hand, are important indicators of structural phase transitions driven by magnetic interactions such as in spin-Peierls compounds. Dimer fluctuations are key participants in phonon-assisted optical absorption processes of magnetic chain compounds and are thus observable in optical conductivity measurements .
Here we consider the exactly solvable $`s=\frac{1}{2}`$ $`XX`$ chain with a magnetic field in the direction transverse to the spin coupling (in spin space) . The Hamiltonian reads
$`H={\displaystyle \underset{n=1}{\overset{N}{}}}\mathrm{\Omega }s_n^z+{\displaystyle \underset{n=1}{\overset{N}{}}}J\left(s_n^xs_{n+1}^x+s_n^ys_{n+1}^y\right).`$ (1.1)
We set $`g\mu _B=1,\mathrm{}=1`$, use the exchange constant as the energy unit, and measure the magnetic field $`\mathrm{\Omega }`$ in units of $`J`$. The Jordan-Wigner transformation to spinless lattice fermions maintains the bilinear operator structure,
$`H={\displaystyle \underset{n=1}{\overset{N}{}}}\mathrm{\Omega }\left(c_n^{}c_n{\displaystyle \frac{1}{2}}\right)+{\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{N}{}}}J\left(c_n^{}c_{n+1}c_nc_{n+1}^{}\right).`$ (1.2)
A Fourier transform, $`c_\kappa =N^{1/2}_{n=1}^N\mathrm{exp}\left(\text{i}\kappa n\right)c_n`$, brings (1.2) into diagonal form:
$`H={\displaystyle \underset{\kappa }{}}\mathrm{\Lambda }_\kappa \left(c_\kappa ^{}c_\kappa {\displaystyle \frac{1}{2}}\right),\mathrm{\Lambda }_\kappa =\mathrm{\Omega }+J\mathrm{cos}\kappa .`$ (1.3)
For periodic boundary conditions in (1.1) the allowed values of the fermion momenta $`\kappa _i`$ depend on whether the number $`N_f`$ of fermions in the system is even or odd: $`\kappa _i\left\{\left(2\pi /N\right)\left(n+\frac{1}{2}\right)\right\}`$ if $`N_f`$ is even or $`\kappa _i\left\{\left(2\pi /N\right)n\right\}`$ if $`N_f`$ is odd. Fermion momenta within the first Brillouin zone are specified by integers $`n=N/2,N/2+1,\mathrm{},N/21`$ (if $`N`$ is even) or $`n=(N1)/2,(N1)/2+1,\mathrm{},(N1)/2`$ (if $`N`$ is odd). The Fermi level in the band $`\mathrm{\Lambda }_\kappa `$ is controlled by the magnetic field $`\mathrm{\Omega }`$. The number of fermions can vary between an empty band $`(N_f=0)`$ and a full band $`(N_f=N)`$ and is related to the quantum number $`S^z`$ (the $`z`$-component of the total spin) of the same model in the spin representation: $`S^z=N_fN/2`$.
The dimer and trimer fluctuation operators for the $`XX`$ model (1.1) will be introduced in Sec. 2. The two-fermion and four-fermion dynamic structure factors associated with these fluctuation operators will be discussed in Secs. 3 and 4, respectively. Finally, in Sec. 5 we give the conclusions and perspectives for future work.
## 2 Fluctuation operators and dynamic structure factors
Most fluctuation operators of interest are constructed from local operators $`A_n`$ of the model system under consideration:
$`A_\kappa ={\displaystyle \frac{1}{\sqrt{N}}}{\displaystyle \underset{n=1}{\overset{N}{}}}\mathrm{exp}\left(\text{i}\kappa n\right)A_n.`$ (2.1)
Associated with each fluctuation operator (2.1) is a dynamic structure factor,
$`S_{AA}(\kappa ,\omega )=2\pi {\displaystyle \underset{\lambda \lambda ^{}}{}}{\displaystyle \frac{\mathrm{exp}\left(\beta E_\lambda ^{}\right)}{Z}}\left|\lambda ^{}|A_\kappa |\lambda \right|^2\delta \left(\omega E_\lambda +E_\lambda ^{}\right),`$ (2.2)
which describes fluctuations of a specific kind. Here $`E_\lambda `$ and $`|\lambda `$ are the eigenvalues and the eigenvectors of $`H`$, and $`Z`$ is the partition function. Of particular interest is the zero-temperature limit $`T=0`$ (i.e. $`\beta \mathrm{}`$), where the thermal fluctuations fade away leaving pure quantum fluctuations in the wake. What remains in (2.2) are transitions between the ground state and all excited states that can be reached by the fluctuation operator $`A_\kappa `$:
$`S_{AA}(\kappa ,\omega )\stackrel{(\beta \mathrm{})}{=}2\pi {\displaystyle \underset{\lambda }{}}\left|\mathrm{GS}|A_\kappa |\lambda \right|^2\delta \left(\omega \omega _\lambda \right),\omega _\lambda =E_\lambda E_{\mathrm{GS}}.`$ (2.3)
The dynamically relevant spectrum observable in (2.2) or (2.3) may vary considerably between fluctuation operators. Among other things, the spectrum is sensitive to their symmetry properties.
For the $`s=\frac{1}{2}`$ transverse $`XX`$ chain (1.1) the most important and most widely studied dynamic structure factors are those for the local spin operators
$$s_n^z=c_n^{}c_n\frac{1}{2},s_n^+=s_n^x+\text{i}s_n^y=c_n^{}\mathrm{exp}\left(\text{i}\pi \underset{j=1}{\overset{n1}{}}c_j^{}c_j\right),s_n^{}=s_n^x\text{i}s_n^y=\mathrm{exp}\left(\text{i}\pi \underset{j=1}{\overset{n1}{}}c_j^{}c_j\right)c_n.$$
(2.4)
At zero temperature the dynamic spin structure factor $`S_{zz}(\kappa ,\omega )`$ is known to couple exclusively to the continuum of particle-hole excitations in the fermion representation, whereas $`S_{xx}(\kappa ,\omega )=S_{yy}(\kappa ,\omega )`$ couples to excitations involving an arbitrarily high number of fermion excitations from the ground state .
The fluctuation operators considered here are constructed from local spin operators on nearest and next-nearest neighbor sites. The dimer fluctuation operator $`D_\kappa `$ and trimer fluctuation operator $`T_\kappa `$ are obtained via (2.1) from
$$D_n=s_n^xs_{n+1}^x+s_n^ys_{n+1}^y=\frac{1}{2}\left(c_n^{}c_{n+1}c_nc_{n+1}^{}\right)$$
(2.5)
and
$$T_n=s_n^xs_{n+2}^x+s_n^ys_{n+2}^y=\frac{1}{2}\left(c_n^{}c_{n+2}c_nc_{n+2}^{}2c_n^{}c_{n+1}^{}c_{n+1}c_{n+2}+2c_nc_{n+1}^{}c_{n+1}c_{n+2}^{}\right),$$
(2.6)
respectively. There is no unique way of defining dimer and trimer fluctuation operators. The most suitable choice depends on the nature and the symmetry of the model system at hand. The operators (2.5) and (2.6) have the advantage that the associated dynamic structure factors $`S_{DD}(\kappa ,\omega )`$ and $`S_{TT}(\kappa ,\omega )`$ can be analyzed exactly for the $`s=\frac{1}{2}`$ transverse $`XX`$ chain (1.1) in the fermion representation.
As a motivation for the dimer and trimer operators used in this study, we offer a twofold argument. For a completely dimerized state, where nearest-neighbor spin correlations alternate between zero and a nonzero value along the chain, the operator $`\sqrt{N}D_\pi `$ plays the role of dimer order parameter. Likewise, for a completely trimerized state, where next-nearest neighbor spin correlations assume a period-three sequence of values zero, zero, nonzero, the operator $`\sqrt{N}T_{2\pi /3}`$ plays the role of trimer order parameter.
Conversely, if we perturb the uniform $`XX`$ Hamiltonian (1.1) by interactions of the form
$`H_D=\epsilon {\displaystyle \underset{n=1}{\overset{N}{}}}\mathrm{cos}\left(\pi n\right)D_n`$ (2.7)
or
$`H_T=\epsilon {\displaystyle \underset{n=1}{\overset{N}{}}}\mathrm{cos}\left({\displaystyle \frac{2\pi }{3}}n\right)T_n,`$ (2.8)
the ground state becomes dimerized or trimerized, respectively. In the former case, nearest neighbor correlations are modified by period-two perturbative corrections of order $`\epsilon ,+\epsilon `$ and in the latter case by period-three corrections of order $`\frac{1}{2}\epsilon ,\frac{1}{2}\epsilon ,+\epsilon `$.
We may formally introduce the polymer fluctuation operator of order $`l`$, $`𝒫_\kappa ^{(l)}`$, via (2.1) from
$`𝒫_n^{(l)}=s_n^xs_{n+l}^x+s_n^ys_{n+l}^y.`$ (2.9)
It includes the dimer and trimer operators for $`l=1,2`$, respectively: $`𝒫_n^{(1)}=D_n`$, $`𝒫_n^{(2)}=T_n`$. From the fermion representation of the polymer operator (2.9) as carried out explicitly in (2.5) and (2.6) for the lowest two orders, it is evident that the dynamic polymer structure factor $`S_{𝒫𝒫}(\kappa ,\omega )`$ at zero temperature will involve $`2m`$-fermion excitations with $`m=1,2,\mathrm{},l`$ from the ground state. For an infinitely long chain $`(N\mathrm{})`$ the polymer fluctuation operator and the function $`S_{𝒫𝒫}(\kappa ,\omega )`$ may thus serve useful roles in attempts to understand the enormously complex dynamic spin structure factors $`S_{xx}(\kappa ,\omega )=S_{yy}(\kappa ,\omega )`$. Such tools by which the complexity of the dynamically relevant excitation spectrum can be gradually and systematically increased are not only useful for the calculations in the fermion representation as performed here but also for the recently developed techniques of calculating transition rates for the $`XX`$ model in the framework of the Bethe ansatz . The time-dependent polymer correlation function is related, in the limit $`l\mathrm{}`$, to time-dependent spin correlation functions as follows:
$`𝒫_n^{(l)}(t)𝒫_{n+m}^{(l)}(0)\stackrel{(l\mathrm{})}{}2s_n^x(t)s_{n+m}^x(0)^2+2s_n^x(t)s_{n+m}^y(0)^2.`$ (2.10)
Note that $`s_n^x(t)s_{n+m}^y(0)`$ is nonzero only if $`\mathrm{\Omega }0`$.
## 3 Two-fermion dynamic structure factors
We start with the dynamic quantities which are governed by particle-hole excitations. The equilibrium time-dependent correlation functions for the operators $`s_n^z(t)`$ and $`D_n(t)`$ can be evaluated directly:
$`s_n^z(t)s_{n+l}^z(0)s^z^2={\displaystyle \frac{1}{N^2}}{\displaystyle \underset{\kappa _1,\kappa _2}{}}e^{\text{i}\left(\kappa _1\kappa _2\right)l}\mathrm{exp}\left(\text{i}\left[\mathrm{\Lambda }_{\kappa _1}\mathrm{\Lambda }_{\kappa _2}\right]t\right)n_{\kappa _1}\left(1n_{\kappa _2}\right),`$ (3.1)
$`D_n(t)D_{n+l}(0)D^2={\displaystyle \frac{1}{N^2}}{\displaystyle \underset{\kappa _1,\kappa _2}{}}\mathrm{cos}^2{\displaystyle \frac{\kappa _1+\kappa _2}{2}}e^{\text{i}\left(\kappa _1\kappa _2\right)l}\mathrm{exp}\left(\text{i}\left[\mathrm{\Lambda }_{\kappa _1}\mathrm{\Lambda }_{\kappa _2}\right]t\right)n_{\kappa _1}\left(1n_{\kappa _2}\right),`$ (3.2)
where $`n_\kappa =1/\left(1+\mathrm{exp}\left(\beta \mathrm{\Lambda }_\kappa \right)\right)`$ is the Fermi function and
$`s^z={\displaystyle \frac{1}{N}}{\displaystyle \underset{n=1}{\overset{N}{}}}s_n^z={\displaystyle \frac{1}{2N}}{\displaystyle \underset{\kappa }{}}\mathrm{tanh}{\displaystyle \frac{\beta \mathrm{\Lambda }_\kappa }{2}},`$ (3.3)
$`D={\displaystyle \frac{1}{N}}{\displaystyle \underset{n=1}{\overset{N}{}}}D_n={\displaystyle \frac{1}{2N}}{\displaystyle \underset{\kappa }{}}\mathrm{cos}\kappa \mathrm{tanh}{\displaystyle \frac{\beta \mathrm{\Lambda }_\kappa }{2}}.`$ (3.4)
The associated dynamic structure factors,
$`S_{AA}(\kappa ,\omega )={\displaystyle \underset{l=1}{\overset{N}{}}}\mathrm{exp}\left(\text{i}\kappa l\right){\displaystyle _{\mathrm{}}^{\mathrm{}}}\text{d}t\mathrm{exp}\left(\text{i}\omega t\right)\left(A_n(t)A\right)\left(A_{n+l}(0)A\right),`$ (3.5)
all of which involve two-fermion transitions, are obtained by Fourier transform. The resulting expressions for $`N\mathrm{}`$ can be brought into the form
$`S_{zz}(\kappa ,\omega )={\displaystyle _\pi ^\pi }\text{d}\kappa _1n_{\kappa _1}\left(1n_{\kappa _1+\kappa }\right)\delta \left(\omega +\mathrm{\Lambda }_{\kappa _1}\mathrm{\Lambda }_{\kappa _1+\kappa }\right)={\displaystyle \underset{\kappa ^{}}{}}{\displaystyle \frac{n_\kappa ^{}\left(1n_{\kappa +\kappa ^{}}\right)}{2|J\mathrm{sin}\frac{\kappa }{2}\mathrm{cos}\left(\frac{\kappa }{2}+\kappa ^{}\right)|}},`$ (3.6)
$`S_{DD}(\kappa ,\omega )={\displaystyle \underset{\kappa ^{}}{}}{\displaystyle \frac{\mathrm{cos}^2\left(\frac{\kappa }{2}+\kappa ^{}\right)n_\kappa ^{}\left(1n_{\kappa +\kappa ^{}}\right)}{2|J\mathrm{sin}\frac{\kappa }{2}\mathrm{cos}\left(\frac{\kappa }{2}+\kappa ^{}\right)|}},`$ (3.7)
where $`\pi \kappa ^{}\pi `$ are the solutions of the equation
$`\omega =2J\mathrm{sin}{\displaystyle \frac{\kappa }{2}}\mathrm{sin}\left({\displaystyle \frac{\kappa }{2}}+\kappa ^{}\right).`$ (3.8)
The dynamic structure factors (3.6), (3.7) are governed by the two-fermion (particle-hole) excitation continuum the properties of which were examined in Refs. . This continuum is well visible in Figs. LABEL:fig01 and LABEL:fig02 below, which we include mainly for comparison with new results. At zero temperature, $`T=0`$, the two-fermion excitation continuum has the following lower, middle and upper boundaries in the $`(\kappa ,\omega )`$-plane<sup>1</sup><sup>1</sup>1We assume that $`0\kappa \pi `$ in the rest equations of this Section; these equations are valid also for $`\pi \kappa 0`$ after the change $`\kappa \kappa `$.
$`{\displaystyle \frac{\omega _l}{|J|}}=2\mathrm{sin}{\displaystyle \frac{\kappa }{2}}\left|\mathrm{sin}\left({\displaystyle \frac{\kappa }{2}}\alpha \right)\right|,`$ (3.9)
$`{\displaystyle \frac{\omega _m}{|J|}}=2\mathrm{sin}{\displaystyle \frac{\kappa }{2}}\mathrm{sin}\left({\displaystyle \frac{\kappa }{2}}+\alpha \right),`$ (3.10)
$`{\displaystyle \frac{\omega _u}{|J|}}=\{\begin{array}{cc}2\mathrm{sin}{\displaystyle \frac{\kappa }{2}}\mathrm{sin}\left({\displaystyle \frac{\kappa }{2}}+\alpha \right),\hfill & \text{if}\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0}\kappa \pi 2\alpha ,\hfill \\ 2\mathrm{sin}{\displaystyle \frac{\kappa }{2}},\hfill & \text{if}\pi 2\alpha \kappa \pi ,\hfill \end{array}`$ (3.13)
respectively. The parameter $`\alpha =\mathrm{arccos}\left(\mathrm{\Omega }/|J|\right)`$ varies from $`\pi `$ when $`\mathrm{\Omega }=|J|`$ to $`0`$ when $`\mathrm{\Omega }=|J|`$. The $`\omega `$-profiles at fixed $`\kappa `$ of the two-fermion dynamic structure factors may exhibit square-root divergences (a common density-of-states effect in one dimension) when $`\omega 2|J|\mathrm{sin}(\kappa /2)`$. At $`T>0`$ the lower boundary of two-fermion excitation continuum is smeared out. The spectral weight in (3.6) and (3.7) is now confined to $`|\omega |2|J|\mathrm{sin}(\kappa /2)`$.
Closed-form expressions for the two-fermion dynamic structure factors (3.6), (3.7) exist in the low-temperature and high-temperature limits. At $`T=0`$ we have
$`S_{zz}(\kappa ,\omega )={\displaystyle \frac{1}{\sqrt{4J^2\mathrm{sin}^2\frac{\kappa }{2}\omega ^2}}}\{\begin{array}{cc}\mathrm{\Theta }(\omega \omega _l)\mathrm{\Theta }(\omega _u\omega ),\hfill & \text{if}\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0}\kappa \pi 2\alpha ,\hfill \\ \left(\mathrm{\Theta }(\omega \omega _l)+\mathrm{\Theta }(\omega \omega _m)\right)\mathrm{\Theta }(\omega _u\omega ),\hfill & \text{if}\pi 2\alpha \kappa \pi ,\hfill \end{array}`$ (3.16)
$`S_{DD}(\kappa ,\omega )={\displaystyle \frac{\sqrt{4J^2\mathrm{sin}^2\frac{\kappa }{2}\omega ^2}}{4J^2\mathrm{sin}^2\frac{\kappa }{2}}}\{\begin{array}{cc}\mathrm{\Theta }(\omega \omega _l)\mathrm{\Theta }(\omega _u\omega ),\hfill & \text{if}\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0}\kappa \pi 2\alpha ,\hfill \\ \left(\mathrm{\Theta }(\omega \omega _l)+\mathrm{\Theta }(\omega \omega _m)\right)\mathrm{\Theta }(\omega _u\omega ),\hfill & \text{if}\pi 2\alpha \kappa \pi ,\hfill \end{array}`$ (3.19)
and at $`T\mathrm{}`$ we have
$`S_{zz}(\kappa ,\omega )={\displaystyle \frac{1}{2\sqrt{4J^2\mathrm{sin}^2\frac{\kappa }{2}\omega ^2}}}\mathrm{\Theta }\left(2|J|\mathrm{sin}{\displaystyle \frac{\kappa }{2}}|\omega |\right),`$ (3.20)
$`S_{DD}(\kappa ,\omega )={\displaystyle \frac{\sqrt{4J^2\mathrm{sin}^2\frac{\kappa }{2}\omega ^2}}{8J^2\mathrm{sin}^2\frac{\kappa }{2}}}\mathrm{\Theta }\left(2|J|\mathrm{sin}{\displaystyle \frac{\kappa }{2}}|\omega |\right).`$ (3.21)
The zero-temperature results for $`S_{zz}(\kappa ,\omega )`$ can be found in Eq. (2.3) of Ref. and for $`S_{DD}(\kappa ,\omega )`$ at $`\mathrm{\Omega }=0`$ in Eq. (3.2) of Ref. .
In Figs. LABEL:fig01 and LABEL:fig02 we show the dynamic structure factors (3.6), (3.7) at zero temperature $`T=0`$ and different values of the transverse field (panels a, b, c), and at $`T\mathrm{}`$ (panel d). The results for $`T\mathrm{}`$ are independent of $`\mathrm{\Omega }`$. As we can see, the two-fermion dynamic structure factors are nonzero within the two-fermion excitation continuum in the $`(\kappa ,\omega )`$-plane. Their spectral-weight distributions are controlled by the Fermi functions, the multiplicity of the solution of Eq. (3.8), the singularities in the density of one-particle states, and the explicit form of the rest of the integrand in (3.6), (3.7). Another two-fermion dynamic quantity will be presented in Sec. 4, namely the two-fermion contribution to the dynamic trimer structure factor, $`S_{TT}^{(2)}(\kappa ,\omega )`$.
## 4 Four-fermion dynamic structure factor
Next we consider the dynamics of the trimer fluctuations. The method remains the same but its execution is more tedious. In addition to two-fermion transitions also four-fermion transitions contribute to the trimer fluctuations. The expression for the equilibrium time-dependent trimer-trimer correlation function reads
$`T_n(t)T_{n+l}(0)T^2={\displaystyle \frac{1}{N^2}}{\displaystyle \underset{k_1,k_2}{}}C^{(2)}(\kappa _1,\kappa _2)e^{\text{i}\left(\kappa _1\kappa _2\right)l}\mathrm{exp}\left(\text{i}\left[\mathrm{\Lambda }_{\kappa _1}\mathrm{\Lambda }_{\kappa _2}\right]t\right)n_{\kappa _1}\left(1n_{\kappa _2}\right)`$
$`+{\displaystyle \frac{1}{N^4}}{\displaystyle \underset{\kappa _1,\kappa _2,\kappa _3,\kappa _4}{}}C^{(4)}(\kappa _1,\kappa _2,\kappa _3,\kappa _4)e^{\text{i}\left(\kappa _1+\kappa _2\kappa _3\kappa _4\right)l}\mathrm{exp}\left(\text{i}\left[\mathrm{\Lambda }_{\kappa _1}+\mathrm{\Lambda }_{\kappa _2}\mathrm{\Lambda }_{\kappa _3}\mathrm{\Lambda }_{\kappa _4}\right]t\right)`$
$`\times n_{\kappa _1}n_{\kappa _2}\left(1n_{\kappa _3}\right)\left(1n_{\kappa _4}\right),`$ (4.1)
where
$`T={\displaystyle \frac{1}{N}}{\displaystyle \underset{n=1}{\overset{N}{}}}T_n=c_2+2c_1^22c_0c_2,`$ (4.2)
$`C^{(2)}(\kappa _1,\kappa _2)=\left(12c_0\right)^2\mathrm{cos}^2\left(\kappa _1+\kappa _2\right)+4c_1\left(12c_0\right)\left(\mathrm{cos}^2\left(\kappa _1+{\displaystyle \frac{\kappa _2}{2}}\right)+\mathrm{cos}^2\left({\displaystyle \frac{\kappa _1}{2}}+\kappa _2\right)\right)`$
$`+4c_1^2\left(\mathrm{cos}^2\kappa _1+\mathrm{cos}^2\kappa _2\right)+8\left(c_2+c_1^2+2c_0c_2\right)\mathrm{cos}^2{\displaystyle \frac{\kappa _1+\kappa _2}{2}}+8c_1^2\mathrm{cos}^2{\displaystyle \frac{\kappa _1\kappa _2}{2}}`$
$`+4c_1\left(12c_04c_2\right)\left(\mathrm{cos}^2{\displaystyle \frac{\kappa _1}{2}}+\mathrm{cos}^2{\displaystyle \frac{\kappa _2}{2}}\right)`$
$`+4c_28c_18c_1^2+4c_2^2+16c_0c_18c_0c_2+16c_1c_2,`$ (4.3)
$`C^{(4)}(\kappa _1,\kappa _2,\kappa _3,\kappa _4)=16\mathrm{sin}^2{\displaystyle \frac{\kappa _1\kappa _2}{2}}\mathrm{sin}^2{\displaystyle \frac{\kappa _3\kappa _4}{2}}\mathrm{cos}^2{\displaystyle \frac{\kappa _1+\kappa _2+\kappa _3+\kappa _4}{2}}.`$ (4.4)
Here we have introduced the function $`c_p=\left(1/N\right)_\kappa \mathrm{cos}\left(p\kappa \right)n_\kappa `$. For $`N\mathrm{}`$ and at zero temperature we have $`c_0=1`$ if $`\mathrm{\Omega }|J|`$, $`c_0=\alpha /\pi `$ if $`|J|\mathrm{\Omega }|J|`$, $`c_0=0`$ if $`|J|\mathrm{\Omega }`$, $`c_p=\left(\text{sgn}(J)\right)^p\mathrm{sin}\left(p\alpha \right)/\left(p\pi \right)`$ if $`|\mathrm{\Omega }||J|`$ and $`c_p=0`$ otherwise ($`p=1,2,\mathrm{}`$ ). At $`T\mathrm{}`$ we have $`c_p=\frac{1}{2}\delta _{p,0}`$. The resulting dynamic trimer structure factor (3.5) for $`N\mathrm{}`$ then has the following form:
$`S_{TT}(\kappa ,\omega )=S_{TT}^{(2)}(\kappa ,\omega )+S_{TT}^{(4)}(\kappa ,\omega ),`$ (4.5)
where
$$S_{TT}^{(2)}(\kappa ,\omega )=_\pi ^\pi \text{d}\kappa _1C^{(2)}(\kappa _1,\kappa _1+\kappa )n_{\kappa _1}\left(1n_{\kappa _1+\kappa }\right)\delta \left(\omega +\mathrm{\Lambda }_{\kappa _1}\mathrm{\Lambda }_{\kappa _1+\kappa }\right),$$
(4.6)
$`S_{TT}^{(4)}(\kappa ,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{4\pi ^2}}{\displaystyle _\pi ^\pi }\text{d}\kappa _1{\displaystyle _\pi ^\pi }\text{d}\kappa _2{\displaystyle \text{d}\kappa _3C^{(4)}(\kappa _1,\kappa _2,\kappa _3,\kappa _1+\kappa _2\kappa _3+\kappa )}`$ (4.7)
$`\times `$ $`n_{\kappa _1}n_{\kappa _2}\left(1n_{\kappa _3}\right)\left(1n_{\kappa _1+\kappa _2\kappa _3+\kappa }\right)\delta \left(\omega +\mathrm{\Lambda }_{\kappa _1}+\mathrm{\Lambda }_{\kappa _2}\mathrm{\Lambda }_{\kappa _3}\mathrm{\Lambda }_{\kappa _1+\kappa _2\kappa _3+\kappa }\right).`$
The spectral weight in this quantity comes from both the two-fermion (one particle and one hole) excitation continuum and the four-fermion (two particles and two holes) excitation continuum. Let us first discuss the properties of the four-fermion excitation continuum and then the properties of $`S_{TT}(\kappa ,\omega )`$. At $`T=0`$ the four-fermion excitation continuum (for $`J=|J|<0`$) is determined by the conditions
$`{\displaystyle \frac{\omega }{|J|}}=\mathrm{cos}\kappa _1+\mathrm{cos}\kappa _2\mathrm{cos}\kappa _3\mathrm{cos}\kappa _4,\kappa =\kappa _1\kappa _2+\kappa _3+\kappa _4(\text{mod}(2\pi )),`$
$`\mathrm{cos}\kappa _1{\displaystyle \frac{\mathrm{\Omega }}{|J|}},\mathrm{cos}\kappa _2{\displaystyle \frac{\mathrm{\Omega }}{|J|}},\mathrm{cos}\kappa _3{\displaystyle \frac{\mathrm{\Omega }}{|J|}},\mathrm{cos}\kappa _4{\displaystyle \frac{\mathrm{\Omega }}{|J|}},`$ (4.8)
$`\pi \kappa _{1,2,3}\pi `$, $`\pi \kappa \pi `$. Equations (4) imply that the four-fermion excitation continuum (like the two-fermion excitation continuum) exists only if the magnetic field does not exceed the saturation field: $`|\mathrm{\Omega }||J|`$.
An analytical expression for the lower boundary of the four-fermion excitation continuum in the $`(\kappa ,\omega )`$-plane depends on $`\mathrm{\Omega }`$ and $`\kappa `$ and is given by one of the following expressions (see Appendix for additional details):
$`{\displaystyle \frac{\omega _l^{(1)}}{|J|}}=2\mathrm{sin}{\displaystyle \frac{|\kappa |}{2}}\mathrm{sin}\left(\alpha {\displaystyle \frac{|\kappa |}{2}}\right),`$ (4.9)
$`{\displaystyle \frac{\omega _l^{(2)}}{|J|}}=4\mathrm{cos}{\displaystyle \frac{\kappa }{4}}\mathrm{cos}\left(\alpha +{\displaystyle \frac{|\kappa |}{4}}\right),`$ (4.10)
$`{\displaystyle \frac{\omega _l^{(3)}}{|J|}}=2\mathrm{sin}\left(\alpha +{\displaystyle \frac{|\kappa |}{2}}\right)\mathrm{sin}\left(2\alpha +{\displaystyle \frac{|\kappa |}{2}}\right),`$ (4.11)
$`{\displaystyle \frac{\omega _l^{(4)}}{|J|}}=2\mathrm{sin}\left(\alpha {\displaystyle \frac{|\kappa |}{2}}\right)\mathrm{sin}\left(2\alpha {\displaystyle \frac{|\kappa |}{2}}\right),`$ (4.12)
$`{\displaystyle \frac{\omega _l^{(5)}}{|J|}}=4\mathrm{sin}{\displaystyle \frac{|\kappa |}{4}}\mathrm{sin}\left(\alpha {\displaystyle \frac{|\kappa |}{4}}\right).`$ (4.13)
The range in $`\kappa `$ over which for a given $`\mathrm{\Omega }`$ one of the expressions (4.9) – (4.13) forms the lower boundary of the four-fermion continuum can be read off Fig. LABEL:fig03. The darkness in this gray-scale plot is a measure of the size of the energy threshold of the four-fermion continuum (white means zero excitation energy i.e. a soft mode). The boundary between the region $`i`$ (where $`\omega _l^{(i)}`$ is the lower boundary) and the region $`j`$ (where $`\omega _l^{(j)}`$ is the lower boundary) follows from the matching condition $`\omega _l^{(i)}=\omega _l^{(j)}`$ and is given by the formula $`|\kappa |=l_{ij}(\alpha )`$ where
$`l_{12}(\alpha )=4\mathrm{arctan}{\displaystyle \frac{\mathrm{tan}\alpha \sqrt{\mathrm{tan}^2\alpha 3}}{3}},|\kappa |{\displaystyle \frac{2\pi }{3}},`$ (4.14)
$`l_{13}(\alpha )=\pi \alpha ,{\displaystyle \frac{\pi }{2}}|\kappa |{\displaystyle \frac{2\pi }{3}},`$ (4.15)
$`l_{14}(\alpha )=2\alpha ,`$ (4.16)
$`l_{23}(\alpha )=2\pi 4\alpha ,`$ (4.17)
$`l_{34}(\alpha )=|\kappa |+\mathrm{cos}\alpha {\displaystyle \frac{1}{2}},{\displaystyle \frac{2\pi }{3}}|\kappa |\pi ,`$ (4.18)
$`l_{45}(\alpha )=4\alpha .`$ (4.19)
The boundary between regions 2 and 4 is determined by the cubic equation
$`\left(\mathrm{sin}|\kappa |2\mathrm{sin}{\displaystyle \frac{|\kappa |}{2}}\right)\mathrm{tan}^3\alpha +\left(3+2\mathrm{cos}{\displaystyle \frac{\kappa }{2}}+3\mathrm{cos}\kappa \right)\mathrm{tan}^2\alpha `$
$`\left(2\mathrm{sin}{\displaystyle \frac{|\kappa |}{2}}+3\mathrm{sin}|\kappa |\right)\mathrm{tan}\alpha +3+2\mathrm{cos}{\displaystyle \frac{\kappa }{2}}\mathrm{cos}\kappa =0.`$ (4.20)
Typical lower boundaries of the four-fermion continuum for several values of $`\mathrm{\Omega }`$ can be seen in Fig. LABEL:fig04. The soft modes according to (4.9) – (4.13) are given by
$`|\kappa _0|=\{0,2\pi 4\alpha ,2\alpha ,4\alpha \}.`$ (4.21)
Alternatively, the soft modes (4.21) may be determined directly from (4). They occur when $`\mathrm{cos}\kappa _1=\mathrm{cos}\kappa _2=\mathrm{cos}\kappa _3=\mathrm{cos}\kappa _4=\mathrm{cos}\alpha `$.
The upper boundary of the four-fermion continuum for $`0\mathrm{\Omega }|J|/\sqrt{2}`$ is given by
$`{\displaystyle \frac{\omega _u^{(1)}}{|J|}}=4\mathrm{cos}{\displaystyle \frac{\kappa }{4}}.`$ (4.22)
For $`|J|/\sqrt{2}\mathrm{\Omega }|J|`$ the upper boundary is given by (4.22) only if $`|\kappa |4\alpha `$ whereas, if $`4\alpha |\kappa |\pi `$, it is given by
$`{\displaystyle \frac{\omega _u^{(2)}}{|J|}}=4\mathrm{cos}{\displaystyle \frac{\kappa }{4}}\mathrm{cos}\left(\alpha {\displaystyle \frac{|\kappa |}{4}}\right)`$ (4.23)
(see Appendix). The upper boundaries of the four-fermion continuum for several values of $`\mathrm{\Omega }`$ can be seen in Fig. LABEL:fig04 in comparison with the corresponding two-fermion continuum.
The four-fermion continuum always contains the two-fermion continuum. The lower boundaries coincide in part. The upper boundaries are different. In the zero field case we have $`\omega _l=|J|\mathrm{sin}|\kappa |`$ for both continua. The upper boundaries are $`\omega _u=4|J|\mathrm{cos}\frac{\kappa }{4}`$ and $`\omega _u=2|J|\mathrm{sin}\frac{|\kappa |}{2}`$ for the two-fermion and four-fermion continua, respectively. As the saturation field $`\mathrm{\Omega }=|J|`$ is approached from below, the two-fermion continuum narrows to a branch and then disappears whereas the four-fermion continuum remains an extended region, bounded by $`\omega _l=4|J|\mathrm{sin}^2\frac{\kappa }{4}`$ and $`\omega _u=4|J|\mathrm{cos}^2\frac{\kappa }{4}`$, and then disappears abruptly.
Now consider the equation
$`\sqrt{{\displaystyle \underset{j=1}{\overset{3}{}}}\left({\displaystyle \frac{}{\kappa _j}}\left[\mathrm{cos}\kappa _1+\mathrm{cos}\kappa _2\mathrm{cos}\kappa _3\mathrm{cos}\left(\kappa +\kappa _1+\kappa _2\kappa _3\right)\right]\right)^2}=0.`$ (4.24)
It is satisfied for
$`{\displaystyle \frac{\omega _s^{(1)}}{|J|}}=2\mathrm{sin}{\displaystyle \frac{|\kappa |}{2}},`$ (4.25)
$`{\displaystyle \frac{\omega _s^{(2)}}{|J|}}=4\mathrm{sin}{\displaystyle \frac{|\kappa |}{4}},`$ (4.26)
$`{\displaystyle \frac{\omega _s^{(3)}}{|J|}}=4\mathrm{cos}{\displaystyle \frac{\kappa }{4}}.`$ (4.27)
Thus, for $`\kappa `$ or $`\omega `$ approaching the curves (4.25) – (4.27) in $`(\kappa ,\omega )`$-space, the quantity
$`S(\kappa ,\omega )`$ $`=`$ $`{\displaystyle _\pi ^\pi }\text{d}\kappa _1{\displaystyle _\pi ^\pi }\text{d}\kappa _2{\displaystyle _\pi ^\pi }\text{d}\kappa _3S(\kappa _1,\kappa _2,\kappa _3,\kappa )`$ (4.28)
$`\times `$ $`\delta \left(\omega |J|\mathrm{cos}\kappa _1|J|\mathrm{cos}\kappa _2+|J|\mathrm{cos}\kappa _3+|J|\mathrm{cos}\left(\kappa +\kappa _1+\kappa _2\kappa _3\right)\right)`$
exhibits cusp singularities (akin to density-of-states effects in three dimensions). The exact nature of the cusps also depends on the factor $`S(\kappa _1,\kappa _2,\kappa _3,\kappa )`$, which varies between different dynamic structure factors with a four-fermion part. It always includes the factor $`n_{\kappa _1}n_{\kappa _2}\left(1n_{\kappa _3}\right)\left(1n_{\kappa _1+\kappa _2\kappa _3+\kappa }\right)`$ as can be seen in expression (4.7). In Fig. LABEL:fig05 we show the $`\omega `$-dependence of $`S(\kappa ,\omega )`$ as given by (4.28) at $`\kappa =0,\frac{\pi }{2},\frac{2\pi }{3},\pi `$ when $`S(\kappa _1,\kappa _2,\kappa _3,\kappa )=1`$ and when $`S(\kappa _1,\kappa _2,\kappa _3,\kappa )=n_{\kappa _1}n_{\kappa _2}\left(1n_{\kappa _3}\right)\left(1n_{\kappa _1+\kappa _2\kappa _3+\kappa }\right)`$ for several values of $`\mathrm{\Omega }`$.
At $`T>0`$ the lower boundary of the four-fermion excitation continuum is smeared out and the upper boundary becomes $`\omega _u=4|J|\mathrm{cos}\frac{\kappa }{4}`$.
The properties of the multimagnon continua of quantum spin chains have been examined in some detail in the recent paper of Barnes . In particular, the lower/upper boundary of the two- and higher magnon continua were determined. It was shown that the boundary curves under certain conditions may exhibit discontinuous changes in composition and cusps. Moreover, a behavior of the density of (two- and higher magnon) states on the continuum boundaries and within the continuum was considered and the existence of discontinuities was pointed out. These features of one-dimensional quantum spin systems are expected to become accessible experimentally in high-resolution inelastic neutron scattering on alternating chain and ladder materials.
Interestingly, the $`s=\frac{1}{2}`$ transverse $`XX`$ chain which can be mapped onto noninteracting spinless fermions provides an excellent example of a system whose dynamic properties are governed by continua of multifermion excitations. In particular, the dynamics of trimer fluctuations provides a direct motivation for analyzing the four-fermion excitation continuum. Unlike in the analysis reported in Ref. , where the statistics of the elementary excitations (magnons) is not known, here the quasiparticles are known to be fermions and the consequences are fully taken into account.
Finally, let us examine the explicit expression for the dynamic trimer structure factor $`S_{TT}(\kappa ,\omega )`$ (4.5). In Figs. LABEL:fig06, LABEL:fig07, LABEL:fig08, LABEL:fig09 we present the zero-temperature dynamic trimer structure factor at different values of $`\mathrm{\Omega }`$. In Fig. LABEL:fig10 we present the same quantity at infinite temperature. We show separately the two-fermion contribution (panels a) and the four-fermion contribution (panels b) as well as the sum of these contributions (panels c). We observe how the spectral weight is spread across the four-fermion continuum. We also see that the two-fermion contribution stands out in terms of spectral weight. The two-fermion and four-fermion contributions are comparable in intensity at $`T=0`$ and small $`\mathrm{\Omega }`$. As $`\mathrm{\Omega }`$ increases the two-fermion contribution becomes more important and it completely dominates as $`\mathrm{\Omega }|J|`$. In the high-temperature limit the two-fermion contribution is very dominant but the four-fermion continuum is still in evidence.
## 5 Conclusions
In summary, we have investigated some aspects of the dynamics of the $`s=\frac{1}{2}`$ transverse $`XX`$ chain examining, in particular, the dynamics of dimer and trimer operators. For this purpose we have calculated several dynamic structure factors on a rigorous basis within the Jordan-Wigner representation. While the dynamic dimer structure factor $`S_{DD}(\kappa ,\omega )`$ and the dynamic transverse spin structure factor $`S_{zz}(\kappa ,\omega )`$ are governed by fermionic one-particle-one-hole excitations, the dynamic trimer structure factor $`S_{TT}(\kappa ,\omega )`$ also contains contributions from two-particle-two-hole excitations. We have described the structure of the two-fermion and four-fermion excitation spectra in some detail and investigated the distribution of spectral weight in $`S_{TT}(\kappa ,\omega )`$ across these continua at zero and nonzero temperature. In particular, we have established the boundaries of the four-fermion spectral range, the locations of soft modes, and the singularity structure, which includes one-dimensional and three-dimensional density-of-states features (van Hove singularities).
An alternative technique to evaluate dynamic structure factor of quantum spin chains is based on the Bethe ansatz solutions . Recently such an approach has been applied to the $`s=\frac{1}{2}`$ $`XX`$ chain . Moreover, the relation between spinons or magnon quasiparticles and Jordan-Wigner fermions was discussed in some detail. It will be interesting to interpret the two-fermion and four-fermion excitations discussed here in terms of the Bethe ansatz solution as studied in Ref. .
## Acknowledgments
This study was performed within the framework of the STCU project #1673. O. D., J. S. and G. M. thank the Wilhelm und Else Heraeus-Stiftung for the kind hospitality during the 288. WE-Heraeus-Seminar on the Theme of “Quantum Magnetism: Microscopic Techniques for Novel States of Matter” (Bad Honnef, 2002) when the present study was launched. O. D. acknowledges the kind hospitality of the University of Dortmund in the spring of 2003 when part of the work was done. O. D. expresses gratitude to the Max-Planck-Institut für Physik komplexer Systeme (Dresden) for its hospitality in the spring of 2004.
## Appendix
To find the lower and upper boundaries of the four-fermion excitation continuum (4) at fixed $`\mathrm{\Omega }`$ and $`\pi \kappa \pi `$ we search for the extrema of $`\omega /|J|`$ as given in (4) and the values of $`\kappa _1`$, $`\kappa _2`$, $`\kappa _3`$, $`\kappa _4`$ at which such extrema occur. Typical results are reported in Fig. LABEL:fig11 (lower boundary) and Fig. LABEL:fig12 (upper boundary).
In Fig. LABEL:fig11 we show the dependence on $`|\kappa |`$ of $`\kappa _1`$, $`\kappa _2`$, $`\kappa _3`$, $`\kappa _4`$ where $`\omega /|J|`$ assumes a local minimum. We distinguish five different regions. The global minimum yields the lower continuum boundary. If $`0\kappa \kappa _a`$
$`\kappa _1=\alpha \kappa ,\kappa _2=\kappa _3=\kappa _4=\alpha `$ (A.1)
and
$`{\displaystyle \frac{\omega _l}{|J|}}=\mathrm{cos}\left(\alpha \kappa \right)\mathrm{cos}\alpha =2\mathrm{sin}{\displaystyle \frac{\kappa }{2}}\mathrm{sin}\left(\alpha {\displaystyle \frac{\kappa }{2}}\right)={\displaystyle \frac{\omega _l^{(1)}}{|J|}};`$ (A.2)
if $`\kappa _a\kappa \kappa _b`$
$`\kappa _1=\kappa _2=\alpha ,\kappa _3=\kappa _4={\displaystyle \frac{\kappa }{2}}+\alpha \pi `$ (A.3)
and
$`{\displaystyle \frac{\omega _l}{|J|}}=2\mathrm{cos}\alpha 2\mathrm{cos}\left({\displaystyle \frac{\kappa }{2}}+\alpha \pi \right)=4\mathrm{cos}{\displaystyle \frac{\kappa }{4}}\mathrm{cos}\left(\alpha +{\displaystyle \frac{\kappa }{4}}\right)={\displaystyle \frac{\omega _l^{(2)}}{|J|}}`$ (A.4)
etc. The values of $`\kappa _a`$, $`\kappa _b`$, $`\kappa _c`$, $`\kappa _d`$ follow from the matching conditions.
In Fig. LABEL:fig12 we show the dependence on $`|\kappa |`$ of $`\kappa _1`$, $`\kappa _2`$, $`\kappa _3`$, $`\kappa _4`$ where $`\omega /|J|`$ assumes a local maximum. We distinguish two different regions. The global maximum yields the upper continuum boundary. If $`0\kappa \kappa _A`$
$`\kappa _1=\kappa _2={\displaystyle \frac{\kappa }{4}},\kappa _3=\kappa _4=\pi +{\displaystyle \frac{\kappa }{4}}`$ (A.5)
and
$`{\displaystyle \frac{\omega _u}{|J|}}=4\mathrm{cos}{\displaystyle \frac{\kappa }{4}}={\displaystyle \frac{\omega _u^{(1)}}{|J|}};`$ (A.6)
if $`\kappa _A\kappa \pi `$
$`\kappa _1=\kappa _2=\alpha ,\kappa _3=\kappa _4=\alpha \pi +{\displaystyle \frac{\kappa }{2}}`$ (A.7)
and
$`{\displaystyle \frac{\omega _u}{|J|}}=2\mathrm{cos}\alpha +2\mathrm{cos}\left(\alpha {\displaystyle \frac{\kappa }{2}}\right)=4\mathrm{cos}{\displaystyle \frac{\kappa }{4}}\mathrm{cos}\left(\alpha {\displaystyle \frac{\kappa }{4}}\right)={\displaystyle \frac{\omega _u^{(2)}}{|J|}}`$ (A.8)
etc. From the matching condition we find $`\kappa _A=4\alpha `$.
FIGURE CAPTIONS
Fig. 1. $`S_{zz}(\kappa ,\omega )`$ at $`T=0`$ and (a) $`\mathrm{\Omega }=0`$, (b) $`\mathrm{\Omega }=0.3`$, (c) $`\mathrm{\Omega }=0.9`$, and (d) at $`T\mathrm{}`$ (independent of $`\mathrm{\Omega }`$; only $`\omega 0`$ is shown).
Fig. 2. $`S_{DD}(\kappa ,\omega )`$ at $`T=0`$ and (a) $`\mathrm{\Omega }=0`$, (b) $`\mathrm{\Omega }=0.3`$, (c) $`\mathrm{\Omega }=0.9`$, and (d) at $`T\mathrm{}`$ (independent of $`\mathrm{\Omega }`$; only $`\omega 0`$ is shown).
Fig. 3. Lower boundary $`\omega _l=\mathrm{min}\left(\omega _l^{(j)}\right)`$, $`j=1,\mathrm{},5`$ of the four-fermion excitation continuum versus wave number $`\kappa `$ and transverse field $`\mathrm{\Omega }`$ (for $`|J|=1`$).
Fig. 4. Lower and upper boundaries of the two-fermion and four-fermion continua for $`|J|=1`$ and $`\mathrm{\Omega }=0`$, $`\mathrm{\Omega }=0.3`$, $`\mathrm{\Omega }=0.6`$, $`\mathrm{\Omega }=0.9`$. The two-fermion continuum is shown shaded.
Fig. 5. $`S(\kappa ,\omega )`$ as given in (4.28) versus $`\omega `$ at $`\kappa =0,\frac{\pi }{2},\frac{2\pi }{3},\pi `$ with $`S(\kappa _1,\kappa _2,\kappa _3,\kappa )=1`$ (bold curves), and $`S(\kappa _1,\kappa _2,\kappa _3,\kappa )=n_{\kappa _1}n_{\kappa _2}\left(1n_{\kappa _3}\right)\left(1n_{\kappa _1+\kappa _2\kappa _3+\kappa }\right)`$ for $`\mathrm{\Omega }=0`$ (solid curves), $`\mathrm{\Omega }=0.3`$ (long-dashed curves), $`\mathrm{\Omega }=0.6`$ (short-dashed curves), $`\mathrm{\Omega }=0.9`$ (dotted curves). Vertical lines mark the values of $`\omega _s^{(j)}`$, $`j=1,2,3`$ as given in (4.25) – (4.27). Note the different vertical scales left and right.
Fig. 6. $`S_{TT}(\kappa ,\omega )`$ at $`T=0`$ and $`\mathrm{\Omega }=0`$. Separate plots are shown for (a) $`S_{TT}^{(2)}(\kappa ,\omega )`$, (b) $`S_{TT}^{(4)}(\kappa ,\omega )`$, and (c) the sum $`S_{TT}(\kappa ,\omega )`$.
Fig. 7. $`S_{TT}(\kappa ,\omega )`$ at $`T=0`$ and $`\mathrm{\Omega }=0.3`$. Separate plots are shown for (a) $`S_{TT}^{(2)}(\kappa ,\omega )`$, (b) $`S_{TT}^{(4)}(\kappa ,\omega )`$, and (c) the sum $`S_{TT}(\kappa ,\omega )`$.
Fig. 8. $`S_{TT}(\kappa ,\omega )`$ at $`T=0`$ and $`\mathrm{\Omega }=0.6`$. Separate plots are shown for (a) $`S_{TT}^{(2)}(\kappa ,\omega )`$, (b) $`S_{TT}^{(4)}(\kappa ,\omega )`$, and (c) the sum $`S_{TT}(\kappa ,\omega )`$.
Fig. 9. $`S_{TT}(\kappa ,\omega )`$ at $`T=0`$ and $`\mathrm{\Omega }=0.9`$. Separate plots are shown for (a) $`S_{TT}^{(2)}(\kappa ,\omega )`$, (b) $`S_{TT}^{(4)}(\kappa ,\omega )`$, and (c) the sum $`S_{TT}(\kappa ,\omega )`$.
Fig. 10. $`S_{TT}(\kappa ,\omega )`$ at $`T\mathrm{}`$ (independent of $`\mathrm{\Omega }`$; only $`\omega 0`$ is shown). Separate plots are shown for (a) $`S_{TT}^{(2)}(\kappa ,\omega )`$, (b) $`S_{TT}^{(4)}(\kappa ,\omega )`$, and (c) the sum $`S_{TT}(\kappa ,\omega )`$.
Fig. 11. Search for the lower boundary of the four-fermion excitation continuum. Shown are the values of $`\kappa _1`$, $`\kappa _2`$, $`\kappa _3`$, $`\kappa _4`$ at which a minimum of $`\omega /|J|`$ as given in (4) occurs at $`\mathrm{\Omega }=0.3|J|`$ and $`\pi \kappa \pi `$. The dependences $`\kappa _1`$ and $`\kappa _3`$ on $`\kappa `$ are shown by dashed curves, the dependences $`\kappa _2`$ and $`\kappa _4`$ on $`\kappa `$ are shown by dotted curves, the dependence of the minimal value of $`\omega /|J|`$ on $`\kappa `$ is shown by solid curves.
Fig. 12. Search for the upper boundary of the four-fermion excitation continuum. Shown are the values of $`\kappa _1`$, $`\kappa _2`$, $`\kappa _3`$, $`\kappa _4`$ at which a maximum of $`\omega /|J|`$ as given in (4) occurs at $`\mathrm{\Omega }=0.9|J|`$ and $`\pi \kappa \pi `$. The dependences of $`\kappa _1(=\kappa _2)`$ and $`\kappa _3(=\kappa _4)`$ on $`\kappa `$ are shown by dashed curves, the dependence of the maximal value of $`\omega /|J|`$ on $`\kappa `$ is shown by solid curves. |
warning/0507/astro-ph0507070.html | ar5iv | text | # Molecular Gas in NUclei of GAlaxies (NUGA)
## 1 Introduction
The phenomenon of nuclear activity is understood to be a result of the feeding of supermassive black holes (SMBHs) in galactic nuclei. Observational evidence accumulated over the last decade indicates that SMBH exist in most galactic bulges (e.g., Kormendy & Richstone kor95 (1995); Magorrian et al. mag98 (1998); Ferrarese & Merritt fer00 (2000); Gebhardt et al. geb00 (2000)). Among all these massive black holes, very few are highly active. AGN are found in 10$`\%`$ of the local galaxies (Ho et al. ho97 (1997)); however, this percentage is increased up to $``$44$`\%`$ if LINERs are taken into account. One of the present challenges is to understand how AGN can be fed during their lifetime. In the feeding problem the gas supply must come from the whole disk of the host galaxy at large distances compared to the radius of gravitational influence of the central engine. Therefore it is expected that a hierarchy of mechanisms combine to drive virtually all the gas from the large $``$kpc scales down to the inner $``$pc scales. The different spatial scales involved suggest that the various mechanisms at work have very different time scales (Shlosman et al. shl89 (1989, 1990); Combes com01 (2001, 2003); Jogee jog04 (2004)). Recent observational and theoretical evidence indicates that the AGN lifetimes may be as short as $``$a few 10<sup>7</sup>–10<sup>8</sup>yr (Ho et al. ho03 (2003); Martini mar04 (2004); Wada wad04 (2004); Merloni mer04 (2004)). Moreover, Wada (wad04 (2004)) finds evidence that mass accretion may not be constant even during the nominal duty cycle of 10<sup>8</sup>yr, but composed of several shorter episodes with a duration of 10<sup>4-5</sup>yr. This time-scale conspiracy could explain the lack of success of observers in finding any correlation between the presence of $``$kpc scale non-axisymmetric perturbations (e.g., large-scale bars and interactions) and the onset of activity in galaxies, except for very high luminosity objects(QSOs) (Moles et al. mol95 (1995); Mulchaey & Regan mul97 (1997); Knapen et al. kna00 (2000); Krongold et al. kro01 (2001); Schmitt sch01 (2001)). On these spatial scales, Hunt & Malkan (hun99 (1999)) have found a significant correlation between the detection rate of outer rings and the onset of activity.
The search for a universal feeding mechanism has been pursued by looking for morphological features in the central kpc of nearby AGN with high spatial resolution ($``$a few 100 pc), though with limited success. Nuclear stellar bars seem to be as common in AGN as in non-AGN (Regan & Mulchaey reg99 (1999); Laine et al. lai02 (2002)). Furthermore, a similar scenario holds for nuclear spirals: while Martini & Pogge (mar99 (1999)) and Pogge & Martini (pog02 (2002)) initially found a high frequency of dusty spirals in the HST enhanced color images of Seyfert nuclei of their sample, more recently, Martini et al. (mar03 (2003)) have shown on a firmer statistical basis that these nuclear features are not preferentially found in AGN. Similarly there is only weak statistical evidence that nuclear rings are more frequently found in Seyferts (Knapen kna05 (2005)). Hunt & Malkan (huntmalkan04 (2004)) found evidence for an excess in type 2 Seyferts of kpc-scale twisted isophotes, though the cause of such twists is unclear.
On the modeling front, significant progress has been made in recent years on the study of the feeding efficiency of different types of gravitational instabilities: nested bars (e.g., Shlosman et al. shl89 (1989); Friedli & Martinet fri93 (1993); Maciejewski & Sparke mac00 (2000); Englmaier et al. eng04 (2004)), gas spiral waves (e.g., Englmaier & Shlosman eng00 (2000); Maciejewski et al. mac02 (2002); Maciejewski 2004a ; 2004b ), $`m=1`$ perturbations (e.g., Shu et al. shu90 (1990); Junqueira & Combes jun96 (1996); García-Burillo et al. gb00 (2000)) and nuclear warps (e.g., Schinnerer et al. sch00 (2000)). However, the lack of high quality multi-wavelength observational constraints on the different models has thus far made the choice of a single optimal scenario rather difficult.
The study of interstellar gas in AGN is essential to understand the phenomenon of nuclear activity in galaxies and its possible link to circumnuclear star formation. As most of the neutral gas in galactic nuclei is in the molecular phase, CO lines are best suited to undertake high-resolution mapping of AGN hosts, with interferometer resolution of $`<`$100 pc, i.e., the scales on which secondary modes embedded in kpc-scale perturbations are expected to take over. CO lines better trace the total gas column densities than dust extinction probes obtained from HST NIR/optical color images. Most importantly, CO maps provide the gas kinematics (velocity fields and velocity dispersions). This information is essential to characterize gravitational instabilities and to constrain the models. The NUclei of GAlaxies–NUGA–project, fully described by García-Burillo et al. (2003a ; 2003b ), is the first high-resolution ($``$0.5<sup>′′</sup>–1<sup>′′</sup>) CO survey of 12 low luminosity AGN (LLAGN) including the full sequence of activity types (Seyferts, LINERs and transition objects from HII to LINER). In the case of LLAGN, the required mass accretion rates derived from the typical bolometric luminosities of these objects range from 10<sup>-2</sup> to 10<sup>-5</sup>Myr<sup>-1</sup> (from Seyferts to LINERs; e.g., see compilation by Jogee jog04 (2004)). Observations, carried out with the IRAM Plateau de Bure Interferometer (PdBI), have been completed early 2004. NUGA surpasses in both spatial resolution and sensitivity ongoing surveys of nearby AGN conducted at OVRO (MAIN: Jogee et al. jog01 (2001)) and at NRO (Kohno et al. koh01 (2001)).
In this paper we focus on the study of gravitational torques in a subset of NUGA galaxies, which span the range of the different activity classes within our sample: NGC 4321 (transition object: HII/LINER), NGC 4826 (transition object: HII/LINER), NGC 4579 (LINER 1.9/Seyfert 1.9) and NGC 6951(Seyfert 2). Information on the stellar potentials, obtained through available HST and ground-based optical/NIR images of the sample, is used to determine the gravitational torques exerted by the derived stellar potentials on the gaseous disk. The efficiency with which gravitational torques drain the angular momentum of the gas depends first on the strength of the non-axisymmetric perturbations of the potential (m$`>`$0) but, also, on the existence of significant phase shifts between the gas and the stellar distributions. The estimate of these phase shifts necessarily requires the availability of images of comparably high spatial resolution ($``$0.5<sup>′′</sup> in our case) showing the distribution of the stars and the gas. In this paper we purposely neglect the role of gas self-gravity as a source of non-axisymmetry in the gravitational potential. Nuclear galaxy disks with a high gas surface density and a mostly axisymmetric stellar potential can be prone to develop this kind of gas self-gravitating perturbation. We purposely defer the study of pure gas instabilities and their ability to drive gas inflow to a forthcoming publication.
We describe in Sect. 2 the observations used, including high-resolution CO maps and NIR images of NGC 4321, NGC 4826, NGC 4579 and NGC 6951. Sect. 3 interprets them in terms of AGN feeding. Sect. 4 computes from NIR images the gravitational potentials and forces, and deduces from the CO maps the effective torques applied to the gas. From these torques, it is possible to derive time-scales for gas flows and discuss whether gravity torques alone are efficient enough to feed the AGNs. The general implications of these results for the current understanding of AGN feeding are presented in Sect. 5.
## 2 Observations
### 2.1 CO NUGA observations
Observations of the circumnuclear disks of NGC 4826, NGC 4579 and NGC 6951 were carried out as part of the NUGA survey with the PdBI between December 2000 and March 2003. We used the ABCD set of configurations of the array (Guilloteau et al. gui92 (1992)). This assures high spatial resolution ($`<`$1<sup>′′</sup> at the highest frequency) but also an optimum sensitivity to all spatial frequencies in the maps. A previous set of observations of NGC 4826, using data taken with the BCD configurations, has been discussed in García-Burillo et al. (2003b ) (hereafter called paper I). We have observed simultaneously the J=1–0 and J=2–1 lines of <sup>12</sup>CO in single fields. Table 1 lists the spatial resolution of these CO observations as well as other source related parameters. The primary beam size is 42<sup>′′</sup> (21<sup>′′</sup>) in all the 1–0 (2–1) line observations. As the bulk of the relevant nuclear disk emission arises well inside the central 15<sup>′′</sup> in the four galaxies discussed in this paper, observations have not been corrected for primary beam attenuation. During the observations the spectral correlator was split in two halves centered at the transition rest frequencies corrected for the assumed recession velocities. The correlator configuration covers a bandwidth of 580 MHz for each line, using four 160 MHz-wide units; this is equivalent to 1510 km s<sup>-1</sup>(755 km s<sup>-1</sup>) at 115 GHz (230 GHz). Visibilities were obtained using on-source integration times of 20 minutes framed by short ($`2`$ min) phase and amplitude calibrations on nearby quasars. The absolute flux scale in our maps was derived to a 10% accuracy based on the observations of primary calibrators whose fluxes were determined from a combined set of measurements obtained at the 30m telescope and the PdBI array. Image reconstruction was done using standard IRAM/GAG software (Guilloteau & Lucas gui00 (2000)). In this work we use prior observations of the 1–0 line of <sup>12</sup>CO of NGC 4321 made using the BCD configurations of the PdBI and previously published by García-Burillo et al. (gb98 (1998)) (see this paper and Table 1 for details).
Depending on the resolution/sensitivity requirements, we use either naturally or uniformly weighted line maps, as indicated throughout the paper. Uniform weighting in the 2–1 line enables us to achieve subarcsecond spatial resolution in the maps of NGC 4826, NGC 4579 and NGC 6951. By default, all velocities are referred to the systemic velocities (v<sub>sys</sub>, listed in Table 1), as determined from this work and from García-Burillo et al. (gb98 (1998)). Similarly ($`\mathrm{\Delta }\alpha `$, $`\mathrm{\Delta }\delta `$) offsets are relative to the AGN loci derived from our own estimates (Table 1).
Molecular gas masses are derived from the CO(1–0) integrated intensities assuming a CO–to–H<sub>2</sub> conversion factor X=N(H<sub>2</sub>)/I<sub>CO(1-0)</sub>=2.2$`\times `$10<sup>20</sup>cm<sup>-2</sup> K<sup>-1</sup> km<sup>-1</sup> s (Solomon & Barrett sol91 (1991)). When required, molecular gas column densities are inferred from the CO(2–1) integrated intensity maps. In this case CO(2–1) intensities are first corrected by the 2–1/1–0 ratio measured within the equivalent 1–0 beam at each position, and subsequently we convert them into N(H<sub>2</sub>) assuming the X factor referred to above.
We have estimated the percentage of CO(1–0) flux recovered in the PdBI maps by comparing the single-point fluxes detected by the 30m telescope (García-Burillo & Krips, private communication) towards the nuclei of the four galaxies discussed in this paper with the fluxes recovered in the PdBI maps, corrected by primary beam attenuation and convolved to the 30m resolution at this frequency (21$`\mathrm{}`$). The fraction of the flux recovered inside the 21$`\mathrm{}`$ field-of-view ranges from 60$`\%`$-65$`\%`$ in NGC 6951 and NGC 4579 to 75$`\%`$ in NGC 4321 and 90$`\%`$ in NGC 4826. The corresponding values for the 2–1 line of CO are comparable. The missing zero spacing flux in these maps is expected to be found in low-level emission arising in the shape of smooth extended components. As has been shown by Helfer et al. (hel03 (2003)), who estimated the percentage of flux recovery in BIMA SONG galaxies using the 12m NRAO telescope, this percentage is always very high in the central regions of galaxies (i.e., the domain of NUGA maps). The reason is that the velocity gradient is largest at the nuclear regions. CO emission is thus confined to much smaller areas in individual channel maps compared to the outer disk regions where the percentage of flux filtered out can be higher. Considering that the percentage of flux actually present in the PdBI maps over the total single-dish estimate is moderate–to–large in the galaxies studied here (see above), we do not expect that the morphology of the maps will significantly change by the addition of a plateau-like component. Moreover, the gravity torque calculation developed in this paper is based on azimuthal averages made to infer time-scales for the gas flows in these galaxy nuclei. The approach followed makes our results virtually insensitive to the presence of a weak extended component: gravity torques on this type of source distribution will be zero. Therefore we do not expect the derived torque budget to be significantly biased.
### 2.2 Near-infrared and optical observations
We acquired from the HST archive<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute. STScI is operated by the Association of Universities for Research in Astronomy, Inc. under NASA contract NAS 5-26555. broadband images of NGC 4826, NGC 4579 and NGC 6951, including three NICMOS images (F110W and F160W for NGC 6951; F160W for NGC 4826) and four WFPC2 images (F450W and F814W for NGC 4826; F555W and F814W for NGC 4579). The optical images were combined using (crreject) to eliminate cosmic rays, and calibrated according to Holtzman et al. (holtzman (1995)). The NICMOS images were re-reduced with the STSDAS task calnica using the best reference files, and the images were calibrated in the standard way. The “pedestal” effect (see Böker et al. bok99 (1999)) was removed with the van der Marel algorithm<sup>2</sup><sup>2</sup>2http://www.stsci.edu/$`_\stackrel{~}{}`$marel/software/pedestal.html. Sky values were assumed to be zero since the galaxy filled the WFPC2/NICMOS frames, an assumption which makes an error of $``$ 0.1 mag at most, in the corner of the images. A J-H color image of NGC 6951 was constructed from F110W-F160W according to the transformations by Origlia & Leitherer (origlia (2000)). For NGC 4321, we have adopted the (ground-based) K-band image presented by Knapen et al. (kna95 (1995)).
## 3 Observational evidence of ongoing feeding
### 3.1 NGC 4321
The inner r$``$1.5 kpc of this galaxy was mapped by García-Burillo et al. (gb98 (1998)) with the PdBI at moderate ($``$2”) spatial resolution in the 1–0 line emission of <sup>12</sup>CO. NGC 4321 has been classified as a transition object (i.e., HII/LINER) by Ho et al. (ho97 (1997)). As can be seen in Fig. 1a, molecular gas in the nucleus of NGC 4321 is concentrated in a two spiral arm structure that starts at r$``$550 pc, near the end points of a prominent nuclear bar (detected in the K band by Knapen et al. kna95 (1995); see also Fig. 5) and extends out to r$``$1.2 kpc. There is also a central CO source coinciding with the AGN which is marginally resolved (r$``$150 pc) by the $``$2” beam. This central source contains a molecular gas mass of $``$10<sup>8</sup>M. The CO spiral arms mostly lie at the trailing edges of the nuclear bar. This particular geometry determines the feeding budget for the gas in this region (see Sect. 4.2.1). The spiral arms and the central source are connected by a molecular gas bridge which is spatially resolved NW of the nucleus, i.e., at the leading edge of the nuclear bar (component B in Fig. 1b). García-Burillo et al. (gb98 (1998)) interpreted these CO observations and their relation with other gaseous and stellar tracers using numerical simulations of the cloud hydrodynamics. They found that the best fit for the gas flow corresponds to the nuclear bar being decoupled from the large-scale bar. The nuclear bar is fast, with a pattern speed $`\mathrm{\Omega }_p`$150 km s<sup>-1</sup>kpc<sup>-1</sup>. This pushes corotation of the nuclear bar inward, likely inside the outer edges of the CO spiral arms.
The kinematics of molecular gas are characterized by streaming motions detected in the CO spiral arms (see discussion in García-Burillo et al. gb98 (1998)). Closer to the AGN (r$`<`$500 pc), we find that molecular gas also displays significant departures from circular rotation at the location of B. Fig. 1b shows the position-velocity (p-v) plot along a strip at PA=127, purposely oriented to illustrate the CO kinematics at B. The discontinuity in radial velocities at $`\mathrm{\Delta }`$x$``$–4” indicates that the molecular gas flow is decelerated, exactly as expected if gas flows along the leading edges of the bar. Details on the kinematics of molecular gas in the central component are hidden due to the insufficient spatial resolution of the CO maps. We notice however that molecular gas emission is asymmetric with respect to v<sub>sys</sub>: CO emission is preferentially blue-shifted.
As is the case for the other transition object analyzed in this paper (NGC 4826), we have found a large molecular gas concentration ($``$10<sup>8</sup>M) near the AGN in NGC 4321 (r$`<`$150 pc). The analysis of gas kinematics provides evidence for gas fueling at present on intermediate scales: at r$``$200-300 pc from the AGN. However, the spatial resolution of the CO maps does not allow us to probe closer than 100 pc from the AGN.
### 3.2 NGC 4826
The first NUGA maps of the transition object NGC 4826 were published in paper I. The CO images showed already a large concentration of molecular gas ($``$1600M/pc<sup>2</sup> ) in the $``$160 pc-diameter circumnuclear disk (CND) of this galaxy. The distribution of molecular gas in the inner CND is significantly lopsided with respect to the position of the AGN; this suggests that m=1 instabilities may be at work at radial distances of r$``$50–60 pc from the central engine. With the newest 0.5<sup>′′</sup> (10 pc) resolution <sup>12</sup>CO(2–1) observations the distribution of molecular gas in the CND is fully resolved; the CND appears in the new maps as an off-center ringed disk (Fig. 2a). The dynamical center of the galaxy determined from these observations coincides within the errors with a blue point source identified in the B-I HST color map of Fig. 2a. This confirms our earlier findings that the putative super massive black hole lies on the southeastern inner side of the off-center ringed disk.
The gas kinematics in the CND are characterized by the presence of streaming motions. A first analysis of the 2D kinematics performed on the lowest resolution data of paper I indicated that the instabilities identified in the CND (and those of the inner m=1 spiral) may not favor AGN feeding. The information contained in the newest images confirms this result. Fig. 2b shows the departures from circular motion of gas velocities, identified at the northern and southern crossings of the CND ring along the minor axis. Gas velocities become systematically redder (bluer) when we approach the nucleus from the southern (northern) side of the ring along the minor axis. Deprojected onto the galaxy plane (North is the near side), this pattern indicates that the radial velocity component changes from an inflow signature (outside the ring) into an outflow signature (inside the ring). As it is fully discussed in paper I, this measured change of sign across the minor axis is compatible with the pattern expected for a trailing wave outside corotation (fast wave), i.e., the type of perturbation that would not help to drain the gas angular momentum. The driving agent of these ’mainly’ gaseous $`m=1`$ instabilities may not be related to the stellar potential which is essentially featureless and mostly axisymmetric (see Sect. 4.2.2 and discussion of paper I).
Based on the CO maps of NGC 4826 we find little evidence of ongoing AGN feeding at scales r$`<`$150 pc from the AGN. While the molecular gas reservoir of NGC 4826 is abundant ($`>`$3$`\times `$10<sup>7</sup>M) close to its central engine (r$`<`$150 pc), the analysis of gas kinematics provides no evidence that ongoing AGN feeding is at work in this transition object.
### 3.3 NGC 4579
The <sup>12</sup>CO(2–1) emission in the S1.9/L1.9 galaxy NGC 4579 has been mapped at $``$0.5<sup>′′</sup> resolution (see Fig. 3a adapted from García-Burillo et al. 2005 in prep.). The molecular gas distribution in the central $``$1 kpc of NGC 4579 suggests that the gas flow responds to the 9 kpc-diameter stellar bar identified in all the NIR images of this galaxy (e.g. Jarret et al. jar03 (2003)). The <sup>12</sup>CO(2–1) map of Fig. 3a reveals a mass of 3.2$`\times `$10<sup>8</sup>M of molecular gas piled up in two highly contrasted spiral arcs. The CO lanes lie at the leading edges of the stellar bar, which is oriented along PA=58. The northern spiral is more continuous and better delineated than its southern counterpart; it is also very well correlated with the red lane seen North in the V–I color HST image of the galaxy, shown in Fig. 3a. There is little molecular gas at r$`<`$100 pc distance from the central engine of NGC 4579: the closest gas complex of $``$10<sup>6</sup>M lies East of the nucleus at r$``$150 pc (complex E in Fig. 3a, b). There is no molecular gas emission coincident with the AGN itself to a 3$`\sigma `$ detection limit of $``$a few 10<sup>5</sup>M. The V–I color HST image of the galaxy helps to identify the eastern molecular complex as part of a structured disk of 150 pc-diameter and ring-like shape (Fig. 3a, c). The position of the AGN is well defined by its radio continuum emission detected at both 1mm and 3mm coming from a point source that lies close to the southwestern edge of the central disk. This indicates that the $`m=2`$ point-symmetry of the gas flow driven by the bar of NGC 4579 breaks up at r$`<`$200 pc and lets lopsidedness take over. An independent confirmation of this picture comes from the new HST image of the nuclear region of NGC 4579 obtained with the ACS camera at 3300 Å; this image resolves the central disk into a winding $`m=1`$ spiral instability that mimics a ring (Contini con04 (2004)).
The kinematics of molecular gas in the central 1 kpc of NGC 4579 are characterized by the presence of highly non-circular motions detected over the spiral arms and most notably over the central disk. Fig. 3b shows the p-v plot along the kinematic major axis of NGC 4579 (PA=95). The radial velocities of the gas in the central gas disk depart by $`>`$100kms<sup>-1</sup> from the expected pattern of circular rotation: emission of the Eastern gas complex (at $`\mathrm{\Delta }`$x$``$1.5<sup>′′</sup>) appears at highly forbidden negative velocities (i.e., v$`<`$v<sub>sys</sub>). Assuming that the gas flows inside the galaxy plane, the reported velocity deviations measured at E would imply that gas is apparently counter-rotating at a speed of v$``$150kms<sup>-1</sup>. As it is discussed in García-Burillo et al. 2005, this could be qualitatively explained by very eccentric $`m=1`$ orbits. Alternatively, this velocity pattern could be accounted for assuming that gas is flowing out of the galaxy plane, possibly entrained by an expanding shell. The expanding shell scenario is supported by the observed kinematics of the gas close to the N side of the central disk. Fig. 3d shows the p-v plot taken along the declination axis, i.e., very close to the orientation of the minor axis. Molecular gas kinematics at the N complex (where v-v$`{}_{sys}{}^{}<`$–75 kms<sup>-1</sup>) can be interpreted either as gas flowing outward inside the plane along the minor axis (North is the near side), or as a signature of out of the plane motions (similarly to the case of the E complex). In either case this implies that AGN fueling is presently thwarted on these scales.
In summary, most of the molecular gas content of the central 1 kpc of NGC 4579 is trapped in a two arm spiral structure that can be traced from r$``$1 kpc down to r$``$200 pc. Some molecular gas (10<sup>6</sup>M) is detected at r$``$150 pc from the central engine of NGC 4579, but not on the position of the central engine itself ($`<`$a few 10<sup>5</sup>M). The first-order interpretation of the complex gas kinematics at r$`<`$150 pc provides no evidence of ongoing inflow from these scales down to the AGN, but on the contrary, it indicates outflow motions.
### 3.4 NGC 6951
NGC 6951 is a prototypical Seyfert 2 galaxy for which sub-arcsecond resolution <sup>12</sup>CO(2–1) maps have been completed within the NUGA project (Fig. 4a adapted from Schinnerer et al. 2005, in prep.). The molecular gas distribution in the central 1 kpc consists of two nuclear spiral arms that can be traced for over 180; the winding spiral arms end up as a highly contrasted pseudo-ring at r$``$350 pc. The spiral arms can be identified by their red color in the J–H HST image of the galaxy, shown in Fig. 4a. This $`m=2`$ gas instability contains a significant gas reservoir of 3$`\times `$10<sup>8</sup>M which is presently feeding a nuclear starburst, also identified by its intense radio continuum and H$`\alpha `$ emissions (Ho & Ulvestad ho01 (2001), Rozas et al. roz02 (2002)). The geometry of the molecular gas ridges likely reflects the crowding of molecular clouds along the x<sub>2</sub> family of orbits of the prominent stellar bar, detected in all NIR images of NGC 6951 (Márquez & Moles mar93 (1993); Friedli et al. fri96 (1996); Pérez et al. per00 (2000)). As shown in Fig. 4a, only a small amount of molecular gas has succeeded in making its way down to the AGN: most of the molecular gas mass is trapped in the nuclear spiral arms, quite similar to the case of NGC 4579 (though in this galaxy molecular gas likely populates both x<sub>1</sub> and x<sub>2</sub> orbits). A compact unresolved molecular complex (denoted as G in Fig. 4a, b) of $``$a few 10<sup>6</sup>M is detected at the position of the central engine. This component could correspond to a molecular torus (of $``$40-50 pc size). As seen in Fig. 4a, low level CO emission has been tentatively detected inside the ring, bridging the ’apparent’ gap between the N spiral arm (running North from West) and the central source. The bridge is better identified in the p–v plot of Fig. 4b. The molecular gas mass of this emission bridge is more accurately estimated using the natural weighted map: the gas mass amounts to $``$10<sup>7</sup>M (Schinnerer et al. 2005). This could be the northern molecular counterpart of the filamentary spiral structure identified in the J–H HST map (Fig. 4a).
The kinematics of molecular gas in the nuclear spiral arms reveal streaming motions (Schinnerer et al. 2005) also identified in the previous lower resolution CO and HCN maps of NGC 6951 (Kohno et al. koh99 (1999)). Inside the ring (50 pc$`<`$r$`<`$350 pc), however, gas kinematics are compatible with regular rotation (Fig. 4b). Of particular note, this is in clear contrast to the case of NGC 4579. The kinematics of the G–component, although compatible with circular motions, cannot be studied in detail due to insufficient spatial resolution.
While most of the 3$`\times `$10<sup>8</sup>M molecular gas disk is feeding a starburst episode in the nuclear spiral arms at r$``$350 pc, a small amount of molecular gas ($``$a few 10<sup>6</sup>M) has been detected on the central engine. This central component reveals a prior accretion episode down to scales of r$``$50 pc. It is unclear whether molecular gas detected in the bridge component is falling into the nucleus or migrating outward: the gas flow is fully compatible with regular circular motions from r=50 pc to r=350 pc.
## 4 Gravitational torques and AGN fueling
To explore more precisely the efficiency of feeding, we have estimated the gravitational torques exerted by the stellar potentials (derived from the NIR images) on their molecular circumnuclear disks (as given by the NUGA CO maps). After estimating the role of stellar gravitational torques, we will investigate whether other mechanisms are required to explain the low level of nuclear activity in these galaxies.
We first explain the general methodology employed and the basic assumptions in Sect. 4.1. The different steps are described in detail in Sect. 4.1.1 and Sect. 4.1.2. We discuss in Sect. 4.2 the results obtained from the application of this procedure to the four NUGA targets examined in this work.
### 4.1 General methodology
Gravitational forces are computed at each location in the plane of the galaxy, using near-infrared images to derive the underlying gravitational potential. We assume that the total mass budget is dominated by the stellar contribution and thus neglect the effect of gas self-gravity. We also assume a constant M/L ratio, and determine its best value by fitting the rotation curve constrained by the CO observations. From the 2D force field ($`Fx`$,$`Fy`$) we derive the torques per unit mass at each location ($`t(x,y)`$=$`xF_yyF_x`$). This torque field, by definition, is independent of the present gas distribution in the plane. The crucial step consists of using the torque field to derive the angular momentum variations and the associated flow time-scales. As explained below, the link is made through the observed distribution of the gas.
With this aim, we assume that the measured gas column density ($`N(x,y)`$) derived from a CO intensity map at each offset in the galaxy plane is a fair estimate of the probability of finding gas at this location at present. In this statistical approach, we implicitly average over all possible orbits of gaseous particles and take into account the time spent by the gas clouds along the orbit paths. We assume that CO is a good tracer of the total gas column density, since HI mass in the nuclei of galaxies is typically a very small fraction of the total gas mass. The torque field is then weighted by $`N(x,y)`$ at each location to derive the time derivative of the local angular momentum surface density $`dL_s(x,y)/dt`$=$`N(x,y)\times t(x,y)`$.
In order to estimate the gas flows induced by these angular momentum variations we produce azimuthal averages of $`dL_s(x,y)/dt`$ at each radius. The azimuthal average at each radius, using $`N(x,y)`$ as the actual weighting function, represents the global variation of the specific gas angular momentum occurring at this radius ($`dL/dt|_\theta `$). Finally, the time-scales for gas inflow/outflow can be derived by estimating the average fraction of angular momentum transferred in one rotation.
The validity of our estimate of the efficiency of stellar gravity torques to drive angular momentum transfer in the gas is based on the following simple hypothesis: we assume that the gas response to the stellar potential is roughly stationary with respect to the potential reference frame during a few rotation periods. We would like to stress that even in the particular case of nuclear bars, which might decouple from the outer stellar bars under certain circumstances, our assumption is still valid. When there are several stellar pattern speeds at different radii in a galaxy disk, numerical simulations show that the gas response tends to be coupled with the stellar potential pattern (Friedli & Benz fri95 (1995); García-Burillo et al. gb98 (1998); Bournaud & Combes bou02 (2002), 2005 in prep.). The gas response adjusts its pattern speed to that of the dominant stellar pattern at a given radius.
A different case is represented by non-axisymmetric perturbations which can be driven by gas self-gravity and that are partly independent of or possibly decoupled from the stellar perturbations of the disk. Nuclear galaxy disks with a high gas surface density and a mostly axisymmetric stellar potential can be prone to develop this kind of gas self-gravitating perturbation. In this limiting case (not contemplated here), and although our calculation is still formally correct, the inclusion of gas self-gravity is required to derive the correct torque budget. In particular, if the gas disk decouples from the stellar pattern, the azimuthally averaged gravity torques exerted by the stellar pattern on the gas will very likely be close to zero. The main source for the torques if any should come from the gas instability itself.
Note also that any radial variation of the CO-to-H<sub>2</sub> conversion factor (X<sub>CO</sub>) is not expected to affect the estimated time-scales as these are independent of the global normalization factor as a function of radius. If X<sub>CO</sub> varies as a function of azimuth at a fixed radius, there could be a potential bias. However, observational evidence indicates that the dominant variation of X<sub>CO</sub> in the central regions of galaxies is radial (Solomon & Barrett sol91 (1991); Regan et al. reg01 (2001)).
In the approach followed to derive the gravitational potential we have assumed that the bulge is as flattened as the disk, and thus no attempt has been made to separate the bulge from the pure disk contribution. The implicit assumption of a highly flattened bulge is probably not wrong for some barred galaxies but for others it will overestimate the radial forces by at most a factor of 2 (e.g., Buta & Block but01 (2001)). On the other hand, the effect of bulge stretching due to deprojection can enhance the strength of bars, especially if they are aligned with the minor axis and the inclination angles are large. As this is not the case for the galaxies analyzed here, we instead expect that the value derived for the gravity torques will be typically underestimated in our case by a factor 1.5-2 although they will still have the same sign.
#### 4.1.1 Evaluation of stellar potentials
The first step is to derive the stellar potential in the nuclear disks of these galaxies, using the high-resolution NIR images described in Sect. 2.2. Our working hypothesis is that NIR images are less affected than optical ones by dust extinction or by stellar population biases (Quillen et al. quillen (1994)). The images are first deprojected according to the angles $`PA`$ and $`i`$ given in Table 1. The images are then completed in the vertical dimension by assuming an isothermal plane model with a constant scale height, equal to $``$1/12th of the radial scale-length of the image. The potential is then derived by a Fourier transform method. We also assumed a constant mass-to-light (M/L) ratio, obtained by fitting the observed CO rotation curve–$`v_{rot}`$–for each galaxy. The potential–$`\mathrm{\Phi }(R,\theta )`$– is then decomposed in the different m-modes:
$$\mathrm{\Phi }(R,\theta )=\mathrm{\Phi }_0(R)+\underset{m}{}\mathrm{\Phi }_m(R)\mathrm{cos}(m\theta \varphi _m(R))$$
(1)
where $`\mathrm{\Phi }_m(R)`$ and $`\varphi _m(R)`$ represent the amplitude and phase of the m-mode, respectively.
Following Combes and Sanders (com81 (1981)), we define the strength of the $`m`$-Fourier component, $`Q_m(R)`$ as
$$Q_m(R)=m\mathrm{\Phi }_m/R|F_0(R)|$$
(2)
The corresponding strength of the total non-axisymmetric perturbation is defined by:
$$Q_T(R)=\frac{F_T^{max}(R)}{F_0(R)}=\frac{\frac{1}{R}\left(\frac{\mathrm{\Phi }(R,\theta )}{\theta }\right)_{max}}{\frac{d\mathrm{\Phi }_0(R)}{dR}}$$
(3)
where $`F_T^{max}(R)`$ represents the maximum amplitude of the tangential force over all $`\theta `$ and $`F_0(R)`$ is the mean axisymmetric radial force.
Figs. 5a,b to 8a,b illustrate the quantitative description of gravitational potentials given by \[Q<sub>i=1,2</sub>, Q<sub>T</sub>$`\varphi _{i=1,2}`$\] for the four galaxies examined in this work.
#### 4.1.2 Efficiency of gravitational torques
After having calculated the forces per unit mass ($`F_x`$ and $`F_y`$) from the derivatives of $`\mathrm{\Phi }(R,\theta )`$ at each pixel, the torques per unit mass–$`t(x,y)`$–can be computed by:
$$t(x,y)=xF_yyF_x$$
(4)
The sense of the circulation of the gas in the galaxy plane determines the sign of $`t(x,y)`$: positive (negative) if the torque accelerates (decelerates) the gas at $`(x,y)`$. We have next obtained the gravitational torque maps weighted by the gas column densities derived from the CO 1–0 and 2–1 lines, $`N(x,y)`$, i.e., we derive $`t(x,y)\times N(x,y)`$. These represent the effective variations of angular momentum density in the galaxy plane. We show in Fig. 9a to 12a the normalized version of these maps, i.e., divided by \[$`|`$ N(x,y) $`\times `$ t(x,y)$`|`$\]<sub>max</sub> for the 2–1 line (except for NGC 4321, where we used the 1–0 line). To estimate the radial gas flow induced by the torques, we have first computed the torque per unit mass averaged over the azimuth, using $`N(x,y)`$ as the actual weighting function,i.e.:
$$t(R)=\frac{_\theta N(x,y)\times (xF_yyF_x)}{_\theta N(x,y)}$$
(5)
Results obtained for t(R) based on the CO(1–0) and CO(2–1) maps were seen to be virtually identical within the errors. Fig. 9b to 12b show the results derived from the 1–0 line maps. By definition, $`t(R)`$ represents the time derivative of the specific angular momentum–$`L`$–of the gas averaged azimuthally, i.e., $`t(R)`$=$`dL/dt|_\theta `$. To derive azimuthal averages, we assume a radial binning ($`\mathrm{\Delta }R`$) which corresponds to the original resolution of the NIR images. Similarly to the torque maps, the sign of $`t(R)`$, either + or –, defines whether the gas may gain or lose angular momentum, respectively. More precisely, we evaluate the AGN feeding efficiency by deriving the average fraction of the gas specific angular momentum transferred in one rotation (T<sub>rot</sub>) by the stellar potential, as a function of radius, i.e., by the non-dimensional function $`\mathrm{\Delta }L/L`$ defined as:
$$\frac{\mathrm{\Delta }L}{L}=\frac{dL}{dt}|_\theta \times \frac{1}{L}|_\theta \times T_{rot}=\frac{t(R)}{L_\theta }\times T_{rot}$$
(6)
where $`L_\theta `$ is assumed to be well represented by its axisymmetric average, i.e.,$`L_\theta =R\times v_{rot}`$. The absolute value of $`L/\mathrm{\Delta }L`$ determines how long will it take for the stellar potential to transfer the equivalent of the total gas angular momentum. Assuming that the gas response to the stellar potential is stationary with respect to the potential reference frame during a few rotation periods, a small value of $`\mathrm{\Delta }L/L`$ implies that the stellar potential is inefficient at present. The $`t(R)`$ and $`\mathrm{\Delta }L/L`$ curves derived from the 1–0 maps of the analyzed galaxies are displayed in Figs. 9b to 12b.
To calculate how much gas mass is involved in the transfer driven by the stellar potential we have estimated the radial trend for the mass inflow (- sign)/outflow(+ sign) rate of gas per unit length as a function of radius (in units of M yr<sup>-1</sup>pc<sup>-1</sup> in Figs. 13a to 16a) as follows:
$$\frac{d^2M}{dRdt}=\frac{dL}{dt}|_\theta \times \frac{1}{L}|_\theta \times 2\pi R\times N(x,y)|_\theta $$
(7)
where $`N(x,y)|_\theta `$ is the radial profile of N(x,y) averaged over the azimuth for a radial binning $`\mathrm{\Delta }R`$.
The inflow/outflow rates integrated out to a certain radius R can be derived as:
$$\frac{dM}{dt}=\frac{d^2M}{dRdt}\times \mathrm{\Delta }R$$
(8)
Figs. 13b to 16b display these integrated rates in units of M yr<sup>-1</sup>.
### 4.2 Tracking down gravitational torques in NUGA targets
#### 4.2.1 NGC 4321
Fig. 5a shows the K-band image of the nucleus of NGC 4321, obtained by Knapen et al. (kna95 (1995)), deprojected onto the galaxy plane (X<sub>G</sub>/Y<sub>G</sub> coordinates). The NIR image shows a nuclear bar (denoted as BAR(n) in Fig. 5a) which is roughly parallel to the large-scale 8 kpc stellar bar (oriented along PA=112 in sky coordinates). The nuclear bar strength is the main contributor to Q<sub>T</sub>, which is typically $``$0.18 for r$`<`$500 pc (Fig. 5b).
Fig. 9a, shows that the derived torques change sign following a characteristic 2D butterfly pattern. Quadrants I-to-IV will define hereafter the regions where the signs of the torques driven by the dominant perturbation of the stellar potential are expected to be constant (Figs.9-12a). Assuming that the gas circulation is counterclockwise, Fig. 9a shows that the bulk of the CO emission along the spiral arms lies at the trailing edges of the nuclear bar where torques are positive (quadrants I(+) and III(+)). While there is significant CO emission arising from the leading quadrants of the bar (quadrants II(-) and IV(-)) both on intermediate scales (r=250-550 pc) and in the outer disk r$`>`$900 pc, negative torques there are weaker than positive torques measured over the trailing quadrants (see Fig. 9b). Most remarkably, closer to the AGN (r$``$200 pc), the sign of the average torques $`t(R)`$ is also positive. Although the central CO source is marginally resolved in the CO map of García-Burillo et al. (gb98 (1998)), and, therefore, higher resolution observations are required to settle the question, positive torques seem to dominate over negative ones at r$``$200 pc.
In terms of mass inflow rates, the overall feeding budget in NGC 4321 is clearly positive at all radii as shown in Fig. 13. This would imply that the radial gas flow goes outward unless other mechanisms more efficiently transfer the gas angular momentum outwards. At r$``$100 pc, dM/dt$``$+0.4M yr<sup>-1</sup> and the radial rate is still positive at r$``$500 pc, where dM/dt$``$+0.4M yr<sup>-1</sup>. Similarly to the case of the other transition object, NGC 4826 (see Sect. 4.2.2), AGN feeding driven by stellar torques is quenched at present in NGC 4321. Although the gas concentration is large close to the AGN, the present gas configuration favors feeding only on intermediate scales (r=250-550 pc). This result is compatible with the CO-based diagnostic of the gas kinematics discussed in Sect.3.1.
#### 4.2.2 NGC 4826
Fig. 6a shows a zoom into the central $``$200 pc region of the deprojected H-band HST image of NGC 4826. We plot in Fig. 6b Q<sub>i</sub>( $`i`$=$`1,2`$), $`\varphi _i`$ ($`i`$=$`1,2`$) and Q<sub>T</sub> inside the full field-of-view of the image (r=530 pc). As shown in Fig. 6b, the deviations from axisymmetry of the stellar potential are exceedingly small at all radii (Q$`{}_{T}{}^{}`$0.02–0.05) except for r$`<`$75 pc where a weak oval perturbation is detected. The latter is denoted as OVAL(n) in Fig. 6a. This oval seems slightly off-center with respect to the AGN; this would account for the contribution from a $`m=1`$ stellar mode to Q<sub>T</sub>, identified in Fig. 6b.
As already suggested by the first calculations of paper I, we find that the weak stellar perturbations present in the inner 500 pc of NGC 4826 are inefficient at feeding the AGN at present. Fig. 10a shows how the derived torques change sign in the central 200 pc of the disk. The distribution of torques, defined by quadrants I-to-IV in Fig. 10a, shows the expected behavior if torques are mostly due to the oval distortion of NGC 4826. Moreover, stellar torques are weak and mostly positive inside r$``$150 pc (Fig. 10b). Very close to the AGN (r$``$20 pc), the sign of average torques might become marginally negative, but as we lack spatial resolution on these scales, the torque estimate becomes uncertain.
Alternatively, the lopsided gas disk may be unrelated to the stellar oval. Instead, gas self-gravity could be driving the gas response in the nucleus NGC 4826 (see discussion in paper I). In this case the gas disk may be decoupled from the stellar pattern, and therefore the latter would be viewed as a mostly axisymmetric component from the system of reference of the gas, i.e., providing zero torque on the gas disk. The main source for the torques should come from the gas instability itself. In NGC 4826 $`m=1`$ instabilities do not seem to favour gas inflow, however (paper I). In either case, this is compatible with the scenario depicted in Sect. 3.2, in which there is currently no evidence of fueling of the AGN.
In spite of the large molecular gas concentration close to the AGN (r$`<`$150 pc), the overall feeding budget in NGC 4826, quantified by the mass inflow rates shown in Fig. 14, is clearly positive inside the full field-of-view of the NIR image. This implies that the predicted radial mass flow should be dominated by outflow rather than inflow motions. At r$``$100 pc, dM/dt$``$+0.2M yr<sup>-1</sup> and the inflow rate is still positive at r$``$500 pc, where dM/dt$``$+0.45M yr<sup>-1</sup>.
#### 4.2.3 NGC 4579
Fig. 7a shows the central r$``$1 kpc region of the I-band HST image of NGC 4579 deprojected onto the plane of the galaxy. The dominant stellar perturbation in the HST image of the disk is the large-scale 9 kpc bar of NGC 4579, denoted BAR and oriented as shown in Fig. 7a. As illustrated by Fig. 7b, the potential strength of the bar is dominated by an $`m=1`$ mode for r$`<`$250 pc where Q$`{}_{T}{}^{}`$0.10.
Fig. 11a shows the 2D pattern of gravitational torques for r$``$1 kpc. Stellar torques change sign as expected if the orientation of quadrants I-to-IV can be attributed to the large-scale stellar bar of NGC 4579. Assuming that the sense of gas circulation is clockwise in the galactic plane, Fig. 11a shows that the bulk of the CO emission along the spiral arms lies where torques are strong and negative. Fig 11b shows that the azimuthally averaged torques are strong and negative from the edge of the image at r=1200 pc (and most probably from larger radii) down to r=200 pc.
Closer to the AGN (r$`<`$200 pc), the bulk of the CO emission is concentrated North of the nucleus in the trailing quadrant I(+), where stellar torques are positive. The CO complex closest to the AGN is in the leading quadrant II(-), East of the nucleus at r$``$150 pc, where torques are negative. However, Fig. 11b shows that the azimuthally averaged torques are positive for r$`<`$200 pc. As discussed in Sect. 3.3, the interpretation of the complex kinematics observed for molecular gas in this region is not straightforward but does not favor a direct feeding of the AGN. One of the possible scenarios explored by García-Burillo et al. 2005 (in prep.) invokes that gas is being entrained by expanding motions at r$`<`$200 pc. Note that in this case we do not expect the gas pattern to be stationary with respect to the stellar pattern and therefore, the torque calculation developed above may not be relevant on scales r$`<`$200 pc. Nonetheless this would not change the picture where the central engine of NGC 4579 is not being fueled at present.
In summary, the overall mass inflow budget is clearly negative down to r$``$200 pc due to the action of the large-scale bar of NGC 4579 (Fig. 15). The radial flow expected for the gas is clearly inward down to r$``$200 pc. At r$``$500 pc, dM/dt$``$–0.6M yr<sup>-1</sup>. Inside r$``$200 pc stellar torques do not currently help feeding.
#### 4.2.4 NGC 6951
Fig. 8a shows the deprojected J-band HST image of the nuclear region (r$``$900 pc) of NGC 6951. The HST image allows the identification of the central 1.8 kpc of the large-scale bar of NGC 6951 (denoted as BAR in Fig. 8a). The bar clearly extends beyond the image field-of-view (r=900 pc). Moreover we report on the detection of a small oval distortion most prominent inside r=400 pc (denoted as OVAL(n) in Fig. 8a). The $`m=2`$ mode of the oval perturbation dominates the strength of the potential at r$`<`$400 pc. While it is true that we likely underestimate the contribution of the large-scale bar to the strength of the $`m=2`$ mode at r$`<`$400 pc due to the restricted field-of-view of the HST image, we do not expect this contribution to be significant on these scales, considering the low Q<sub>2</sub> values typically shown by large-scale bars at radii of a few 100 pc. The measured strength of the oval is moderate–to–low: Q$`{}_{T}{}^{}`$0.06, i.e., quite comparable to the strength of the oval potential detected in the nucleus of NGC 4826. We note that the strength of the large-scale bar is likely much higher than that of the oval but it is relevant only for the gas which lies beyond the r$``$900 pc circumnuclear CO disk.
The NGC 6951 torque map shown in Fig. 12a suggests indeed that it is the nuclear oval which shapes the NGC 6951 torques map. The orientation of quadrants I-to-IV associated with the nuclear oval fits the pattern of the torques much better than that associated to the bar. As the sense of gas circulation is clockwise in NGC 6951, Fig. 12a shows that CO emission over the spiral arms in quadrants I(-) and III(-) lies at the leading edges of the nuclear oval where torques are negative down to the inner radius of the pseudo-ring, at r$``$300 pc. While the torques due to the large-scale bar on the inner CO ring are weak at present, torques are expected to be strong and negative at the leading edges of the bar. The offset dust lanes detected in the optical images of NGC 6951 (e.g., Pérez et al. per00 (2000)) indicate that gas inflow is at work along the bar on large scales. The combination of the bar and the nuclear oval has created a pattern of torques which are negative down to r$``$300 pc. The bulk of the molecular gas in the central 5 kpc of NGC 6951 has already fallen from the x<sub>1</sub> orbits down to the x<sub>2</sub> orbits of the bar where it feeds a nuclear starburst in the ring. In the vicinity of the AGN, however, the average stellar torques become marginally positive and AGN feeding is not favored by the present configuration of the stellar potential.
Fig. 16 shows the overall mass inflow budget for NGC 6951. The radial flow expected for the gas is clearly inward down to r$``$300 pc. At r$``$500 pc, dM/dt$``$–0.25M yr<sup>-1</sup>. Inside r$``$300 pc, however, stellar torques do not favour AGN feeding at present in this Seyfert.
## 5 Discussion and conclusions
The first results obtained from the analysis of stellar torques in the circumnuclear disks of a subset of NUGA targets reveal complex trends in the details of the feeding history along the Transition/LINER/Seyfert sequence. The overall picture emerging in the two Seyferts/LINERs (i.e., NGC 6951 and NGC 4579) suggests that AGN feeding may proceed in two steps. In a first step, gravity torques do help, by getting rid of gas angular momentum, driving the gas inwards and feeding a nuclear starburst on scales of $``$a few 100 pc. On smaller scales, however, stellar torques play no role in AGN fueling in the current epoch: torques on the gas are not negative all the way to the center, but on the contrary they become positive and quench the feeding. This may explain why molecular gas seems to ’avoid’ the inner 200-300 pc of NGC 6951 and NGC 4579 where we measured M$`{}_{gas}{}^{}<`$a few$`\times `$10<sup>5</sup>–10<sup>6</sup>M (see also the properties of NGC 7217 as discussed by Combes et al. com04 (2004)). Observational evidence indicates that in NGC 4579 molecular gas on these scales is presently flowing outward. In a second step, a mechanism complementary to gravity torques may be required to drive gas inflow on smaller scales.
Most notably, a much larger molecular gas concentration has been driven inwards in the transition objects analyzed in this paper: $``$3$`\times `$10<sup>7</sup>–10<sup>8</sup>M inside r=150 pc. However, we have found that stellar torques are also unable to drain angular momentum from the massive circumnuclear gas disks of NGC 4321 and NGC 4826. The underlying mechanism responsible for halting gas inflow inside r=150 pc seems to be different in these two transition objects. On similar spatial scales, the stellar torques created by the nuclear bar of NGC 4321 are much stronger than those created by the weak oval perturbation of NGC 4826. Furthermore, the $`m=1`$ perturbation identified in the gas disk of NGC 4826, which is probably unrelated to the oval perturbation, could be the main donor of angular momentum to the gas at these radii (paper I).
Our quest for a mechanism of AGN feeding based only on stellar torques has led to a paradoxical result for the four cases studied here: inflow is actually thwarted close to the AGNs (r$`<`$100-200 pc). One possible explanation is that the responsible agent in the stellar potential could be transient, or as short lived as an individual AGN episode (hence: $``$10<sup>7-8</sup> years). The feeding phase from 100 pc to 1 pc could be so short that the smoking gun evidence in the potential is likely to be missed. Alternatively, this temporary inability of gravitational torques could be overcome by other mechanisms that, over time, become competitive with non-axisymmetric perturbations. Among the different mechanisms usually cited in the literature, dynamical friction and viscous torques have been invoked to help AGN feeding on $``$100 pc scales. Gravitational torques can become positive close to the central engine as illustrated above, in contrast with viscous torques or dynamical friction that always favour gas inflow. This implies that stellar torques could regulate the gas flows in galactic nuclei in combination with other mechanisms that secularly drain the gas angular momentum.
We discuss below the efficiency of dynamical friction and viscosity versus gravity torques for driving AGN feeding in the four galaxies analyzed in this paper. This efficiency can be quantitatively estimated by comparing the typical time-scales of these processes with those of the gravitational torques derived above for each galaxy.
### 5.1 Assisting gravity torques
#### 5.1.1 Dynamical friction
Dynamical friction of giant molecular clouds (GMCs) in the stellar bulge of a galaxy is often invoked as a possible mechanism of fueling AGN. According to the well-known Chandrasekhar (1943) formula, the time-scale of dynamical friction for a GMC of mass M<sub>GMC</sub> at a radius $`r`$ where rotational velocity is $`V_{rot}`$ would be T$`{}_{df}{}^{}`$10<sup>8</sup> yr (r/100 pc)$`{}_{}{}^{2}\times `$(10<sup>6</sup> M/M<sub>GMC</sub>)$`{}_{}{}^{1}\times `$($`V_{rot}`$/200kms<sup>-1</sup>). We have derived the decay time-scale (T<sub>df</sub>) of a typical GMC (M$`{}_{GMC}{}^{}`$10<sup>6</sup>M) in the nuclear disks of the 4 galaxies discussed above. We derive T<sub>df</sub> at the radii where we have identified a barrier of positive gravity torques and get $`V_{rot}`$ from the CO observations. The values of T<sub>df</sub> are typically $``$30, except in NGC 4826 where we obtain formally T$`{}_{df}{}^{}`$5 (see Table 2). The comparison between T<sub>df</sub> and the corresponding time-scales for gravity torques T<sub>grav</sub> (defined as T$`{}_{grav}{}^{}L/\mathrm{\Delta }L`$ from Eq. ) implies that friction is a relatively slow process compared to gravity torques, except in NGC 4826 and NGC 6951.
Despite these results, we suspect that the values estimated for T<sub>df</sub> are a strict lower limit and so severely overestimate the efficiency of dynamical friction in all cases. The reason for this assessment is that Chandrasekhar’s formula implicitly assumes that a GMC is a bound rigid body, with a highly concentrated mass. This approximation underestimates the real value of T<sub>df</sub>. Molecular clouds are indeed a loose ensemble of clumps of different sizes, resembling fractals. The typical masses of the densest fragments can be as low as 10<sup>-3</sup> M (e.g., Pfenniger & Combes pfe94 (1994)). Therefore the scattering efficiency of a realistic GMC unit of 10–30 pc size including a complex mass spectrum of clumps inside is much less than estimated above: the Coulomb parameter of a GMC in Chandrasekhar’s formula should tend to vanish. We therefore conclude that dynamical friction is likely to be quite a slow process, which to first approximation can be neglected relative to gravity torques.
#### 5.1.2 Viscosity
Viscous torques are generally weak, and their associated time-scales (T<sub>vis</sub>) are quite long at large radii. The reason behind this is that T<sub>vis</sub> grows as the square of the radius according to the classical diffusion equation of Pringle (pri81 (1981)) which gives the time evolution of an axisymmetric disk distribution, $`\mathrm{\Sigma }`$($`r`$,$`t`$), under the action of viscous transport:
$$\frac{\mathrm{\Sigma }}{t}=r^1\frac{}{r}\left[\frac{[\nu \mathrm{\Sigma }r^3(d\mathrm{\Omega }/dr)]}{r}\times \left(\frac{d(\mathrm{\Omega }r^2)}{dr}\right)^1\right]$$
(9)
In this formula, the kinematic viscosity, $`\nu `$, is usually described as $`\lambda \times \sigma _v`$, where $`\lambda `$ is the characteristic mean-free path for viscous transport and $`\sigma _v`$ is the gas velocity dispersion. Eq. allows to derive by simple dimensional analysis the typical radial dependence of the time-scale for viscous transport in an axisymmetric disk, if we assume that the disk rotation curve can be described by a power law (i.e., V$`{}_{rot}{}^{}r^{1\alpha }`$):
$$T_{vis}\frac{2\alpha }{\alpha }\times \frac{r^2}{\nu }$$
(10)
However, as illustrated below, the exact value of $`T_{vis}`$ strongly depends on the assumed initial density distribution law for the disk. Fig. 17a shows the time evolution of a smooth disk distribution under the action of viscous torques. In this calculation we have taken $`\nu `$=100 pc km s<sup>-1</sup>, assuming $`\lambda `$=10 pc and $`\sigma _v`$=10 km s<sup>-1</sup>. After 200 Myr the radial gas flows are very small. However, the efficiency of viscous transport can be enhanced if the initial distribution of the disk is characterized by strong density gradients. This is the case of nuclear contrasted rings, especially when these are located in the inner regions of galaxies (r$``$100–500 pc) as frequently is the case for early-type barred spirals. Moreover, and contrary to the conventional wisdom (based on low resolution observations) that most galaxy rotation curves are close to rigid body for r$`<`$500 pc, $`\mathrm{\Omega }`$ can show a strong variation with radius in the inner regime. Therefore galactic shear can still be very high on these scales; this also favours viscous transport.
Fig. 17b illustrates the time evolution of a nuclear ring of r$``$100 pc under the action of viscous torques in a galaxy disk. The radial gas flows are remarkably faster than in the previous case: in typically 5$`\times `$10<sup>7</sup> yr, the ring starts to dissolve and a significant amount of gas can reach the center of the galaxy. We compare in Table 2 the characteristic time-scales of viscous torques and gravity torques in the 4 galaxies analyzed in this paper. Power-law rotation curves have been fit within the relevant radial ranges adapted to each galaxy, assuming in all cases $`\nu `$=100 pc km s<sup>-1</sup>, i.e., a viscosity prescription similar to the the axisymmetric smooth disk solution. Viscosity seems a viable mechanism to drive significant gas inflow in NGC 4826 at r=50-100 pc, while it is probably inefficient against the stronger gravity torques in NGC 4579 and NGC 4321 for r$`<`$200 pc. Moreover, gas lying inside the nuclear ring of NGC 6951 could be falling to the AGN favoured by viscous torques at r$`<`$200 pc. Note that these results are unchanged if we allow gravity torques to be a factor of 1.5–2 stronger, as argued in Sect. 4.1.
We can conclude that, in spite of all the uncertainties in the derivation of time-scales, viscous torques are the most viable mechanism to generate gas inflow on scales $``$100-200 pc if they act on a contrasted nuclear ring distribution, and do not have to fight against very strong positive torques from gravity.
### 5.2 A scenario for self-regulated activity in LLAGNs
In this work we have shown that gravity torques exerted by the stellar potential on the gas disks of four LLAGNs, purposely chosen to represent the whole range of activity classes in the NUGA sample, fail to account for the feeding of the AGN at present. We conclude that gravity torques need to be assisted in due time to drive the gas to the center. Based on these results, we develop in this section a simplified general scenario in which the onset of nuclear activity can be understood as a recurrent phase during the typical lifetime of any galaxy. In this scenario the recurrence of activity in galaxies is indirectly related to that of the bar instabilities although, as argued below, the active phases are not necessarily coincident with the maximum strength of a single bar episode. While bars can build up gas reservoirs towards the central regions of galaxies in the shape of nuclear rings in early type objects, the dynamical feed-back associated with these gas flows can destroy or considerably weaken the bars (Norman et al. nor96 (1996); Bournaud & Combes bou02 (2002)). When this happens, gravity torques are negligible and can make way for other competing mechanisms of gas transport, such as viscous torques, in an almost axisymmetric system.
The first step in this evolutionary scenario could begin with an axisymmetric disk of gas. Shortly after the galaxy disk would be prone to a bar instability. If the galaxy is early-type, i.e., with a sufficient bulge or central mass concentration, then the bar pattern is likely to have one or even two Inner Lindblad Resonances (ILRs). A nuclear ring would form either close to the single ILR or between the two ILRs. At this stage, the gas would be driven inwards from corotation to the ILR, and would accumulate in the ring where a burst of star formation can occur. However positive gravity torques would prevent the gas from flowing further in: under the bar forcing, gas inside the ring could be evacuated outwards (Combes com88 (1988)). The bar thus would restrict the amount of fuel available for AGN feeding to that lying very close to the center where it is under the dominant gravitational influence of the black hole. As the potential would be mostly axisymmetric in the central $``$10 pc, gas could remain there feeding the central engine.
The infall of gas driven by a bar is self-destructive, however: it progressively weakens and destroys the bar, so that the potential returns to axisymmetry (Norman et al. nor96 (1996); Bournaud & Combes bou02 (2002)). At this stage gas piled up in the nuclear ring could dissolve to form a smoother disk through viscosity which may have become competitive against gravity torques. The infall of gas produced by viscous torques could replenish the exhausted fuel near the central engine triggering a new phase of nuclear activity.
In this scenario, the feeding of active nuclei in early-type spirals is contemplated as a two step process: first gravity torques bring the gas of the large-scale disk to the nuclear ring, and when the dynamical feedback has destroyed the bar, the viscous torques could smooth out the ring, and bring gas into the central 10 pc, where it is under the influence of the Keplerian potential of the black hole. The disk becomes axisymmetric and the cycle can be restarted at its first step. In particular, the disk will be prone to a new bar instability if gas is accreted from the outer parts of the disk. These multiple bar phases must occur in gas-rich spiral galaxies, in order to account for the observed frequency of bars (Block et al. blo02 (2002); Bournaud & Combes bou02 (2002, 2004)).
We relate the recurrence of episodes of bar formation and destruction with the fueling of active nuclei in galaxies. The self-regulated competition between gravity torques and viscosity in galaxy nuclei may lead typically to several episodes of activity (each lasting for $``$ 10<sup>8</sup> yr) during a single cycle of bar formation/destruction (lasting for 10<sup>9</sup> yr). These activity episodes are not expected to be strongly correlated with the phases of maximum strength for the bar, but they may appear at different evolutionary stages of the bar potential, depending on the balance between gravity torques and viscosity. The observed prevalence of outer rings in Seyferts, interpreted as a sign of bar dissolution, would support this scheme (Hunt and Malkan hun99 (1999)). Furthermore, the ample variety of morphologies revealed by the CO maps of the circumnuclear disks of NUGA targets corroborates that there is no universal pattern associated with LLAGNs (García-Burillo et al. gb04 (2004)). Activity can also be found in a galaxy during its axisymmetric phase (e.g., see the case of the NUGA galaxy NGC 5953, discussed by Combes et al. 2005, in prep).
The gas directly responsible for the Seyfert activity in galaxies like NGC 4579 or NGC 6951 must have been brought to the center during the previous axisymmetric phase, while the bar at present is emptying the region inside the nuclear ring, thus regulating the amount of gas available for the active nucleus. In the case of NGC 4579, gravity torques clearly overcome viscosity, while in NGC 6951 the balance between both mechanisms might be reached soon. In the case of NGC 4826, the nucleus has a very low level activity in spite of the presence of a large amount of gas in the vicinity, since the potential is almost axisymmetric, and the nuclear ring could be undergoing dissolution by viscosity. The active phase would develop during the next phase, while the disk would soon become bar unstable again. However, the role of the $`m=1`$ gas perturbations identified in the disk of NGC 4826 might be to slow down the dissolution of the ring, as kinematics still suggest outward motions close to the AGN. The case of NGC 4321 is still early in the time evolution, as the bar is only now entering the destruction phase, through the formation of a nuclear bar inside the resonant ring. This nuclear bar efficiently prevents nuclear feeding through gravity torques.
The observed fraction of LLAGNs in the Local Universe ($``$44$`\%`$, including LINERs; Ho et al. ho97 (1997)) could impose tight constraints on the expected number of activity episodes per bar formation/destruction cycle in this scenario. Based on these statistics, a rough estimate would raise this number to $``$4. However this estimate can only be taken as mostly speculative for the time being. The principal source of uncertainty, which prevents us from making a more quantitative prediction, resides in the fact that there is no consensus on what is the typical duration of a nominal AGN duty-cycle. The different values of AGN duty-cycles discussed in the literature range from 10<sup>7</sup> to 10<sup>8</sup> yr (Ho et al. ho03 (2003); Martini mar04 (2004); Wada wad04 (2004); Merloni mer04 (2004)). Moreover, the required number of individual episodes could be lowered to $``$a few if we exclude starburst dominated LINERs from the LLAGN population or even more if we restrict the activity classification to Seyferts, the most active members in the family.
In barred galaxies with comparatively less prominent bulges, i.e., characterized by the absence of an ILR barrier, viscosity may not be needed to take over gravity torques close to the AGN. The gas flow should not be stalled at a radius of a few $``$100 pc forming a nuclear ring in this case. There are a few examples in the NUGA survey of strongly barred galaxies showing a strong nuclear concentration of gas instead of a ring. This class of ILR-free bars could be more common in the early Universe during the formation of disks and the growth of supermassive black holes (Sellwood & Shen sel04 (2004)).
Furthermore, other mechanisms different from axisymmetric viscous flow could be efficient at overcoming the ILR barrier imposed by most stellar bars in the Local Universe. As mentioned above, the ultimate agent for the feeding could still reside in the stellar potential but would be transient and their detection might be elusive. In addition, the role of nuclear $`m=1`$ and $`m=2`$ gaseous spirals in AGN feeding is still not completely elucidated. Most interestingly, a few examples of these gas instabilities are found in the available CO maps of AGNs, including some of the galaxies of the NUGA sample (García-Burillo et al. gb04 (2004)).
Finally, the role of gas self-gravity to drive gas inflow may not be negligible in those cases where the distribution of gas is a significant source of non-axisymmetry for the total gravitational potential, especially if the stellar potential itself appears as featureless or mostly axisymmetric. We can estimate the influence of gas self-gravity by calculating the typical gas mass fractions as a function of radius in the case of the galaxies studied in this paper. In NGC 6951 and NGC 4579, the gas mass fractions are very low throughout the disk from r$``$100 pc ($``$1$`\%`$) up to r$``$1000 pc ($``$3$`\%`$). Therefore we do not expect gas instabilities to have a significant influence in the total gravitational potential in these two galaxies. In contrast, gas mass fractions are comparatively larger at the same radii in NGC 4826 and NGC 4321: $``$10-15$`\%`$ at r$``$100 pc in both galaxies. This gas mass fraction is still high at r$``$1000 pc ($`>`$10$`\%`$) in NGC 4321, while it decreases down to $``$5$`\%`$ in NGC 4826 at that distance. This evidence points to a larger influence of gas self-gravity in NGC 4826 (as discussed in paper I) and NGC 4321 (as discussed by Wada et al. wad98 (1998)).
These results highlight the need for a significant increase in the number objects for which high-resolution CO maps are available together with careful case-by-case studies in the quest for evidence of AGN feeding. The scenario proposed in this work, as well as any alternative model accounting for nuclear activity in galaxies, remains to be tested using a larger sample of LLAGNs.
###### Acknowledgements.
We thank the anonymous referee whose comments helped to improve this paper. We would like to thank Andrew J. Baker for his careful reading of the paper and related discussions. We acknowledge the IRAM staff from the Plateau de Bure and from Grenoble for carrying out the observations and help provided during the data reduction. This paper has been partially funded by the Spanish MCyT under projects DGES/AYA2000-927, DGES/AYA2003-7584, ESP2001-4519-PE and ESP2002-01693, and European FEDER funds. |
warning/0507/cond-mat0507015.html | ar5iv | text | # Manipulating vortex motion by thermal and Lorentz force in high temperature superconductors
## I Introduction
One of the core issues in high temperature superconductor is the origin of the pseudogap above $`T_c`$ in the underdoped region. In order to understand the physics of the pseudogap, many theoretical models have been proposed such as spin fluctuationPines , preformed cooper pairEmery1 ; Millis , charge stripesEmery2 , d-density wave (DDW)Affleck ; Chakravarty , etc. Among many of them, the pseudogap state has been considered as the precursor to the superconducting state. In the pseudogap region above $`T_c`$, a significant in-plane Nernst signals has been discovered by Xu et al.Xu in the underdoped La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> single crystal, and they attributed this signal to vortex-like excitations above $`T_c`$. This result has been confirmed in other families of cuprate superconductors WYY1 ; WYY2 . By doing experiments with the magnetic field applied along different directions, Wen et al.Wen gave a strong indication for a 2D feature of the Nernst effect in the pseudogap region of underdoped cuprate superconductors. About the origin of the strong Nernst signal above $`T_c`$, it remains highly controversial. Wang et al. WYY2 ; WYY3 suggested that the large Nernst signal supports the scenarioEmery1 where the superconducting order parameter disappears at a much higher temperature instead of $`T_c`$. KontaniKontani suggested that the pseudogap phenomena including the strong Nernst signal can be well described in terms of the Fermi liquid with antiferromagnetic and superconducting fluctuations. Ussishkin et al.Ussiahkin proposed that the Gaussian superconducting fluctuations can sufficiently explain the Nernst signal in the optimally doped and overdoped region, but in the underdoped region the actual $`T_c`$ is suppressed from meanfield temperature by non-Gassian fluctuation. Tan et al.Tan believed that a preformed-pair alternative to the vortex scenario can lead to a strong Nernst signal. Alexandrov and ZavaritakyAZ proposed a model based on normal state carrier without superconducting fluctuation, which, as they asserted, ‘can describe the anomalous Nernst signal in high-$`T_c`$ cuprate’AZ .
It is thus strongly desired to investigate the feature of Nernst signal below and above $`T_c`$. One possibility is that the strong Nernst signal above $`T_c`$ is originated from or partly from the motion of vortex, but the structure and the feature of the vortices above $`T_c`$ is different from that below $`T_c`$. We thus measured the in-plane Nernst voltage and the resistance of the underdoped La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> single crystals and optimally doped YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> thin films in the magnetic field perpendicular to the ab-plane. Besides the temperature gradient, we applied a transverse current to manipulate the possible motion of the vortices. Especially we measured the longitudinal voltage in a new configuration which may reduce the dissipative contribution from the quasiparticle scattering and enhance the signal due to the possible vortex motion. In this new configuration, a strong signal measured from the Nernst leads above $`T_c`$ is observed in underdoped samples. The signal maybe caused by the motion of quasiparticles and vortex. Our results indicate that the Nernst signal above $`T_c`$ is contributed by vortex. The results were analyzed and we think that if strong Nernst signal in the pseudogap region is contributed by the vortex motion, vortices are different below and above $`T_c`$. Below $`T_c`$ the dissipation is induced by the motion of the Abrikosov vortices. Above $`T_c`$ the dissipation may be caused by the motion of the spontaneously generated unbinded vortex-antivortex pairs (vortex plasma).
## II Experimental techniques
The La<sub>1.89</sub>Sr<sub>0.11</sub>CuO<sub>4</sub> single crystals measured in this work were prepared by the traveling solvent floating-zone techniqueZhou . The perfect crystallinity of the single crystals were characterized by X-ray diffraction patternsZhou . The single crystal sample is shaped into a bar with the dimensions of of 4.2 mm (length) $`\times `$ 1 mm (width) $`\times `$ 0.5 mm (thickness). The crystal along the length direction is . The optimally doped YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> thin films were prepared by co-evaporation on MgO substrates with dimensions of 10 mm (length) $`\times `$ 1 mm (width), and the thickness of the film is about 5000 angstrom. We annealled the YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> film to underdoping state and measured its Nernst voltage in the mixed and normal state. The resistance characteristics of the samples measured is shown in FIG 1.
The measurement configuration in our experiment is illustrated in FIG. 2. The magnetic field is applied along the c-axis of the samples. The $`R(H)`$ curves of the samples were obtained by standard four-point method at different temperatures. A heater with a power of 1 mW (for the single crystal) or 9 mW (for the thin film) producing the thermal gradient is adhered fast on one end of the YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> thin film sample (left) and La<sub>1.89</sub>Sr<sub>0.11</sub>CuO<sub>4</sub> single crystal sample (right). The other end of the samples is adhered fast on a cold sink. Along the direction E-F (perpendicular to the thermal gradient) the Nernst voltage signal $`V_N`$ is measured. Two thermometers are attached onto the sample to detect the temperature gradient of the samples. In order to control the vortex motion by Lorentz force, we applied a direct transverse current of 3 mA along A-B and measured the longitudinal voltage between the points C and D. Since the DC current here can exert a Lorentz force to the vortices down or againt the thermal stream direction, the motion of the vortex can be manipulated by thermal and/or Lorentz force. All leads are stuck onto the samples by solidified silver paste at corresponding electrodes with the contact resistance below 0.1 $`\mathrm{\Omega }`$. All measurements are based on an Oxford cryogenic system ( Maglab-12 ) with temperature fluctuation less than 0.04% and magnetic fields up to 12 T. The Nernst voltage and the longitudinal voltage are measured by a Keithley 182-Nanovoltmeter with a resolution of about 5 nV in present case. During the measurement for Nernst and longitudinal voltage the magnetic field is applied parallel to c \- axis and swept between 12 to -12 T. The Nernst signal $`V_N`$ is obtained by subtracting the positive field value with the negative one in order to remove the Faraday signal during sweeping the field and the possible thermal electric power due to asymmetry of the electrodes. The Longitudinal signal $`V_L`$ is obtained by averaging the data obtained from positive field and the negative field.
## III Results and discussion
### III.1 Pure Nernst Signal
FIG. 3. shows the Nernst voltage measured on the La<sub>1.89</sub>Sr<sub>0.11</sub>CuO<sub>4</sub> single crystal with a thermal gradient along (length) direction at temperatures from 5 K to 180 K. In low temperature region, the Nernst signal is dominated by the motion of Abrikosov vortices. One can see that the Nernst signal is precisely zero when the vortices are frozen in the case of 5, 10 and 15 K. The flow of vortices after melting leads to a Nernst signal increasing drastically with $`H`$. A strong in-plane Nernst signal resulting from vortices flow can also be seen in the curves for 20, 25, 30 K, and can be measured far above $`T_c`$ (29.3 K for our sample). When the temperature is above 80 K, the signal becomes negative and gradually approaches a background. When the temperature is above 150 K the Nernst signal becomes insensitive to temperature and does not change anymore with $`T`$. Similar results appear in the Nernst voltage measured on the underdoped YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> thin film, though it is not as strong as that of La<sub>1.89</sub>Sr<sub>0.11</sub>CuO<sub>4</sub> single crystal sample. But no strong Nernst signal above $`T_c`$ is observed in optimally doped YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> thin film.
### III.2 Flux Flow Resistance
FIG. 4 shows the $`R(H)`$ curves of the La<sub>1.89</sub>Sr<sub>0.11</sub>CuO<sub>4</sub> single crystal at temperatures from 10 to 80 K measured by standard four-point method. A zero resistance can be seen when the flux lattice is frozen at the low temperature of 10, 15 and 20 K. The Abrikosov flux flow dissipation after vortex melting can be seen in higher temperature region. But the curvature due to the motion of Abrikosov vortices disappears when the temperature is above $`T_c`$. The dissipation becomes very weakly dependent on $`H`$ in the normal state, which can not be explained by the motion of Abrikosov vortices.
### III.3 Longitudinal Signal
The dissipation shown above may be comprised of the contributions of both the flux flow and the quasiparticle scattering. The new configuration in our experiment can reduce the dissipative contribution from the quasiparticle scattering and enhance the signal due to the possible vortex motion if we send a current between A-B and measure the voltage between C-D. As shown in FIG. 2, the current was applied in A-B direction, and the longitudinal voltage between the points on the side of the sample C and D was measured. As the C-D direction is perpendicular to the current direction and the two points are located outside the regime where the main current can reach, so the voltage induced by the motion of quasiparticles is practically reduced; while the motion of the vortex crossing C-D induced by the current pump can be detected and are enhanced relatively as a consequence. The results of the longitudinal voltage measurement for underdoped single crystal at the temperature 5$``$100 K and optimally doped YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> thin film at the temperature 72$``$100 K are shown in FIG. 5 and FIG. 6 respectively. We can see in both figures the similar dissipation manner as that of the resistance of the samples in magnetic field shown in FIG. 4. The dissipation deriving from the motion of Abrikosov vortices disappears when the temperature is above $`T_c`$, where the dissipation of the samples is independent on magnetic field at fixed temperatures.
The results from the measurement of Nernst voltage, resistance and longitudinal voltage of our samples at different temperatures may indicate different vortices below and above $`T_c`$, provided that the dissipation is attributed to the vortex motion. Apparently, below $`T_c`$ the dissipation is induced by the motion of the Abrikosov vortices. The Abrikosov vortex dissipation leads to the resistance in mixed state and, at the same time, the phase slip caused by the motion of Abrikosov vortices along the thermal gradient leads to strong Nernst signal. The longitudinal voltage below $`T_c`$ shows a non-linear field dependence which indicates the feature of Abrikosov flux flow.
Above $`T_c`$, the dissipation may be caused by the motion of the spontaneously generated vortex-antivortex pairs. According to the theory derived by BerezinskyBerezinsky and by Kosterlitz and ThoulessKT (BKT), just below a transition temperature $`T_{BKT}`$ which lies very nearby $`T_c`$, in samples there are thermally excited vortices and antivortices binding in pairs, even in the absense of an external magnetic field. Above $`T_{BKT}`$ the pairs unbind to form vortex plasma which can flow freely as pancake vortices with their cores confined inside individual superconducting layers. The possibility of this BKT transition was verified and studied in layered HTS systemTomoko ; Culbertson ; Matsuda ; Pradhan ; Miu ; Martin . In the region above $`T_c`$, the magnetic field applied on the samples influences the neutral vortices plasma, and polarizes part of the vortices. The dissipation caused by vortex and antivortex on the tolal resistance is the same, although the moving directions of them are different. In other words, the dissipation resulted from the current driven vortex motion is not related to the polarization of, but to the total number of the vortex and the antivortex in the sample, which is constant when changing $`H`$. So the resistance of the sample in magnetic field above $`T_c`$ is weakly dependent on the magnetic field as illustrated in FIG. 4. For the same reason, the longitudinal voltage induced by the motion of the vortex plasma in the normal state is also independent on the magnetic field.
In the case of Nernst effect a different consequence appears in the pseudogap region. The vortex plasma redirected by the external magnetic field moves across E-F along the thermal gradient, leading to the phase slips and a transverse Nernst electric field $`E_y`$ appears consequently. The direction of the $`E_y`$ caused by the vortex is different from that of the $`E_y`$ caused by the antivortex moving along the same direction. So the direction and strength of $`E_y`$ is not determined by the total number but by the amount of the difference between vortex and antivortex, i.e. by the sign and the amount of the net vortex. The density of the net vortex in the samples increases with increasing external magnetic field, leading to the increase of Nernst voltage with $`H`$ in the field range of our experiment. The 2D feature of the Nernst effect in the pseudogap region of underdoped cuprate superconductors has been proved by the experiments of Wen et alWen , which supports our postulation above.
### III.4 Manipulating Vortex Motion by Thermal and Lorentz Force
The Nernst signal and the longitudinal voltage as the result from the motion of vortex manipulated solely by the thermal gradient and solely by Lorentz force are discussed above. Now let’s go on to discuss the result of manipulating the motion of the vortex by the thermal gradient and the Lorentz force together. With a transverse current applied along A-B direction and a temperature gradient along the length direction of the La<sub>1.89</sub>Sr<sub>0.11</sub>CuO<sub>4</sub> single crystal at the same time, the voltage $`V_{EF}`$ between the Nernst leads E and F are measured.
FIG. 7 shows the result of $`V_{EF}`$ at 20 K for the case $`I`$ = 0, $`I`$ = 0.1 mA, and $`I`$ = -0.1 mA with and without thermal gradient. $`V_{EF}`$ with thermal gradient and without current (presented by solid triangles) is the Nernst signal at this temperature, where we can see the sign of $`V_{EF}`$ is dependent on the direction of magnetic field. For the case of $``$$`{}_{x}{}^{}T`$ = 0 and $`I`$ 0, $`V_{EF}`$ (solid and open squares) shows the characteristics of Abrikosov vortices and the sign of the $`V_{EF}`$ after vortices melting is independent on the direction of current. The features of the curves is similar to those in FIG. 4. The results manipulating vortex motion by thermal and Lorentz force is shown in the case of $``$$`{}_{x}{}^{}T`$ 0 and $`I`$ = $`\pm `$ 0.1 mA (solid and open circles) in FIG. 7. All the curves show that the transport properties in mixed state is dominated by Abrikosov vortex motion.
FIG. 8 shows the result at 80 K with the current of 0 mA, 3 mA and -3 mA. In FIG. 8, all the curves have been moved along $`V_{EF}`$ axis and hence the values of $`V_{EF}`$ is zero when $`H`$ = 0 for the convenience of comparison. The influence of the Faraday effects and the resistive components of the data are taken away by moving the curves. We can see from FIG. 8 that with positive transverse current, the $`V_{EF}`$ measured under negative magnetic field is almost equivalent to the $`V_{EF}`$ without current, but for the case under positive field, there is an obvious difference between the $`V_{EF}`$ with current and that without current. With negative transverse current, the $`V_{EF}`$ measured under positive magnetic field is almost equivalent to the $`V_{EF}`$ without current, but for the case under negative field, there is an obvious difference between the $`V_{EF}`$ with current and that without current. So we can conclude that the effects of the Lorentz force on the Nernst effect in the normal state is related to the direction of the magnetic field. This phenomenon can not be explained by the effects of quasiparticles and may be explained only by the motion of vortex. So the Nernst signal above $`T_c`$ is also contributed by flux motion, although the vortex structure and feature may be very different below and above $`T_c`$.
## IV Conclusions
The in-plane resistance, Nernst voltage and longitudinal voltage in a new experimental configuration of the underdoped La<sub>1.89</sub>Sr<sub>0.11</sub>CuO<sub>4</sub> single crystal and optimally doped YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> thin film are measured by sweeping magnetic field at fixed temperatures. The vortex motion is successfully manipulated by using thermal and / or Lorentz force in the mixed and normal states of the samples. The Nernst signal in the normal state is observed. The signal maybe caused by the motion of quasiparticles or vortex or both. Our results indicate that Nernst signal above $`T_c`$ may be contributed by vortex motion, though the contribution of the quasiparticle has not been excluded. Different vortices below and above $`T_c`$ are expected if the strong Nernst signal in the pseudogap region was attributed to the vortex motion. Below $`T_c`$ the dissipation is induced by the motion of the Abrikosov vortices. Above $`T_c`$ the dissipation may be caused partly by the motion of the spontaneously generated unbinded vortex-antivortex pairs (vortices plasma).
## V Acknowledgements
This work is supported by the National Natural Science Foundation of China (NSFC), the Ministry of Science and Technology of China, and the Chinese Academy of Sciences. |
warning/0507/hep-ex0507019.html | ar5iv | text | # Observation of a new charmonium-like state produced in association with a 𝐽/𝜓 in 𝑒⁺𝑒⁻ annihilation at √𝑠≈10.6GeV
The Belle Collaboration
## Abstract
We report the first observation of a new charmonium-like state at a mass of $`(3.943\pm 0.006\pm 0.006)\mathrm{GeV}/c^2`$. This state, which we denote as $`X(3940)`$, is observed in the spectrum of masses recoiling from the $`J/\psi `$ in the inclusive process $`e^+e^{}J/\psi +\mathrm{anything}`$. We also observe its decay into $`D^{}\overline{D}`$ and determine its intrinsic width to be less $`52\mathrm{MeV}/c^2`$ at the 90% C.L. These results are obtained from a $`357\mathrm{fb}^1`$ data sample collected with the Belle detector near the $`\mathrm{{\rm Y}}(4S)`$ resonance, at the KEKB asymmetric-energy $`e^+e^{}`$ collider.
Recently there have been a number of reports of new charmonium or charmonium-like states: $`\eta _c(2S)`$ choi1 , $`X(3872)`$ choi2 , $`Y(3940)`$ choi3 and $`Y(4260)`$ baby . The latter three states have not been assigned to any charmonium states in the conventional quark model. Moreover, charmonium production in different processes is not well understood. One striking example is the surprisingly large cross section for double charmonium production in $`e^+e^{}`$ annihilation observed by Belle 2cc and confirmed by BaBar bab . These experimental results have generated renewed theoretical interest in the spectroscopy, decays and production of charmonium.
In this paper we report the observation of a new charmonium-like state above $`D\overline{D}`$ threshold, $`X(3940)`$, produced in the process $`e^+e^{}J/\psi X(3940)`$. We also present results from searches for $`X(3940)`$ decay into $`D\overline{D}`$, $`D^{}\overline{D}`$ and $`J/\psi \omega `$. The data used for this analysis correspond to an integrated luminosity of $`357\mathrm{fb}^1`$ collected by the Belle detector at the $`\mathrm{{\rm Y}}(4S)`$ resonance and nearby continuum at the KEKB asymmetric-energy $`e^+e^{}`$ collider.
The $`J/\psi `$ reconstruction procedure is identical to our previously published analyses 2cc ; 2cc2 . Oppositely charged tracks that are positively identified as muons or electrons are used for $`J/\psi \mathrm{}^+\mathrm{}^{}`$ reconstruction. A partial correction for final state radiation and bremsstrahlung energy loss is performed by including the four-momentum of every photon detected within a $`50\mathrm{mrad}`$ cone around the electron direction in the $`e^+e^{}`$ invariant mass calculation. The two lepton candidate tracks are required to have a common vertex, with a distance from the interaction point in the plane perpendicular to the beam axis smaller than $`1\mathrm{mm}`$. The $`J/\psi \mathrm{}^+\mathrm{}^{}`$ signal region is defined by the mass window $`\left|M_\mathrm{}^+\mathrm{}^{}M_{J/\psi }\right|<30\mathrm{MeV}/c^2`$ ($`2.5\sigma `$). $`J/\psi `$ candidates in the signal window are subjected to a mass and vertex constrained fit to improve their momentum resolution. QED processes are substantially suppressed by requiring the total charged multiplicity ($`N_{\mathrm{ch}}`$) in the event to be greater than four. Background due to $`J/\psi `$ mesons produced from $`B\overline{B}`$ events is removed by requiring the center-of-mass (CM) momentum $`p_{J/\psi }^{}`$ to be greater than $`2.0\mathrm{GeV}/c`$. As in the previous analysis, we define the recoil mass as
$$M_{\mathrm{recoil}}(J/\psi )=\sqrt{(E_{\mathrm{CMS}}E_{J/\psi }^{})^2p_{J/\psi }^2},$$
(1)
where $`E_{J/\psi }^{}`$ is the $`J/\psi `$ CM energy after the mass constraint.
For the study of the $`X(3940)D^{()}\overline{D}`$, we reconstruct $`D^0`$ candidates using five decay modes: $`K^{}\pi ^+`$, $`K^{}K^+`$, $`K^{}\pi ^{}\pi ^+\pi ^+`$, $`K_S^0\pi ^+\pi ^{}`$ and $`K^{}\pi ^+\pi ^0`$; and $`D^+`$ candidates using $`K^{}\pi ^+\pi ^+`$, $`K^{}K^+\pi ^+`$ and $`K_S^0\pi ^+`$. For the $`D^0K^{}\pi ^{}\pi ^+\pi ^+`$ and $`D^0K^{}\pi ^+\pi ^0`$ modes, mass windows of $`\pm 10\mathrm{MeV}/c^2`$ and $`\pm 20\mathrm{MeV}/c^2`$ are used; a $`\pm 15\mathrm{MeV}/c^2`$ mass window is used for all other modes (approximately $`2.5\sigma `$ in each case). To improve their momentum resolution, $`D`$ candidates are refitted to the nominal $`D^0`$ or $`D^+`$ masses. To study the contribution of combinatorial background under the $`D`$ peak, we use $`D`$ sidebands with mass windows that are four times as large. For the $`X(3940)J/\psi \omega `$ search, candidate $`\omega `$ mesons are reconstructed from $`\pi ^+\pi ^{}\pi ^0`$ combinations within $`\pm 20\mathrm{MeV}/c^2`$ ($`2.5\sigma `$) of the nominal $`\omega `$ mass. The $`\omega `$ sideband region is defined by $`30<|M(\pi ^+\pi ^{}\pi ^0)M_\omega |<50\mathrm{MeV}/c^2`$.
The recoil mass spectrum for the inclusive $`J/\psi `$ event sample is shown in Fig. 1. Here, in addition to the three previously reported peaks at the $`\eta _c`$, $`\chi _{c0}`$ and $`\eta _c(2S)`$ masses, there is a fourth peak above $`D\overline{D}`$ threshold.
We perform a fit to this spectrum that includes the three previously seen charmonium states plus a fourth state with mass near $`3.94\mathrm{GeV}/c^2`$. The expected signal line-shapes are determined from Monte Carlo (MC) simulation as described in previous Belle publications 2cc ; 2cc2 . The mass values for the $`\eta _c`$, $`\chi _{c0}`$, $`\eta _c(2S)`$ and $`X(3940)`$ are free parameters in the fit, the widths of $`\eta _c`$ and $`\chi _{c0}`$ are fixed to PDG values pdg and the $`\eta _c(2S)`$ width is fixed to $`\mathrm{\Gamma }=17\mathrm{MeV}/c^2`$ bab\_eta2 . The width of the new state is a free parameter in the fit. The signal function for the $`X(3940)`$ is a convolution of the (zero-width) MC line shape with a Breit-Wigner function. The background is parameterized by a second order polynomial and a threshold term ($`\sqrt{M_{\mathrm{recoil}}(J/\psi )2M_D}`$) to account for a possible contribution from $`e^+e^{}J/\psi D^{()}\overline{D}^{()}`$.
The fit results are given in Table 1. The significance for each signal is defined as $`\sqrt{2\mathrm{ln}(_0/_{\text{max}})}`$, where $`_0`$ and $`_{\text{max}}`$ denote the likelihoods returned by the fits with the signal yield fixed at zero and at the fitted value, respectively. The significance of the $`X(3940)`$ signal is $`5.0\sigma `$. The fitted width of the $`X(3940)`$ state is consistent with zero within its large statistical error: $`\mathrm{\Gamma }=39\pm 26\mathrm{MeV}/c^2`$. The fit results are shown in Fig. 1 as the solid curve; the dashed curve is the background function.
The new state has a mass that is above both the $`D\overline{D}`$ and $`D^{}\overline{D}`$ thresholds. We therefore perform a search for $`X(3940)`$ decays into $`D\overline{D}`$ and $`D^{}\overline{D}`$ final states. Because of the small product of $`D^{()}`$ reconstruction efficiencies and branching fractions, it is not feasible to reconstruct fully the chain $`e^+e^{}J/\psi X(3940),X(3940)D^{()}\overline{D}`$. To increase the efficiency, only one $`D`$ meson in the event is reconstructed in addition to the reconstructed $`J/\psi `$ and the other $`\overline{D}`$ or $`\overline{D}^{}`$ is detected as a peak in the spectrum of masses recoiling against the $`J/\psi D`$ combination. The Monte Carlo simulation for $`e^+e^{}J/\psi D\overline{D}`$ and $`e^+e^{}J/\psi D^{}\overline{D}`$ processes indicates a $`M_{\mathrm{recoil}}`$$`(J/\psi D)`$ resolution of about $`30\mathrm{MeV}/c^2`$ and a separation between these two processes of $`2.5\sigma `$. Figure 2 shows the $`M_{\mathrm{recoil}}`$$`(J/\psi D)`$ spectrum in the $`D`$ mass window (points with error bars) and in the scaled $`D`$ mass sidebands (hatched histogram), where $`D`$ includes $`D^0`$ and $`D^+`$. Some events have multiple $`D`$ candidates. In these cases, only the candidate with invariant mass closest to the nominal $`D`$-meson mass is used. Two peaks around the nominal $`D`$ and $`D^{}`$ masses are clearly visible in this distribution. The excess of real $`D`$ events compared to the $`D`$ sidebands at masses above $`2.1\mathrm{GeV}/c^2`$ is due to $`e^+e^{}J/\psi D^{}\overline{D}^{}`$ or $`e^+e^{}J/\psi D^{()}\overline{D}{}_{}{}^{()}\pi `$ processes. A fit to this spectrum is performed using shapes fixed from MC for three processes ($`J/\psi D\overline{D}`$, $`J/\psi D^{}\overline{D}`$ and $`J/\psi D^{}\overline{D}^{}`$) and a second order polynomial.
The fit gives $`N_{D\overline{D}}=86\pm 17`$ ($`5.1\sigma `$) and $`N_{D^{}\overline{D}}=55\pm 18`$ ($`3.3\sigma `$) events in the $`D`$ and the $`D^{}`$ peaks, respectively. Selecting events from the $`M_{\mathrm{recoil}}(J/\psi D)`$ regions around the $`D`$ and $`D^{}`$ masses ($`\pm 70\mathrm{MeV}/c^2`$), we thus effectively tag the processes $`e^+e^{}J/\psi D\overline{D}`$ and $`e^+e^{}J/\psi D^{}\overline{D}`$.
We constrain $`M_{\mathrm{recoil}}(J/\psi D)`$ to the $`D^{()}`$ nominal mass, thereby improving the $`M(D^{()}\overline{D})M_{\mathrm{recoil}}(J/\psi )`$ resolution by a factor of 2.5 ($`\sigma 10\mathrm{MeV}/c^2`$ after constraint), according to the MC simulation. In the $`X(3940)D^{}\overline{D}`$ case, the reconstructed $`D`$ can be from either the $`X(3940)`$ decays or the $`D^{}`$ decay: the constraint $`M_{\mathrm{recoil}}(J/\psi D)M(D^{})`$ also works in the latter case, as both $`X(3940)D^{}\overline{D}`$ and $`D^{}`$ decays have very little available phase space.
The resulting $`M_{\mathrm{recoil}}(J/\psi )`$ distributions are shown in Figs. 3 a) ($`M(D)`$ region) and 3 b) ($`M(D^{})`$ region). The cross-hatched histograms show the scaled $`D`$ sideband distributions. For events with multiple entries, the candidate with invariant mass closest to the nominal $`D`$-meson mass (for $`D`$ signal window) or closest to the center of the sideband (for sidebands) is used.
An $`X(3940)`$ peak with a resolution that is better than that for the unconstrained $`M_{\mathrm{recoil}}(J/\psi )`$ distribution is evident in Fig. 3 b), corresponding to the decay $`X(3940)D^{}\overline{D}`$. We perform a fit to this distribution. The signal function is a convolution of a Breit-Wigner with a free width and a resolution function fixed to the MC expectation. The background function is a threshold function $`(A+BM(D^{}\overline{D}))\sqrt{M(D^{}\overline{D})M_{\mathrm{thr}}}`$ with $`M_{\mathrm{thr}}M(D^{})+M(D)`$. The fit yields the number of signal events in the peak $`N=24.5\pm 6.9`$ with a statistical significance of $`5.0\sigma `$. The width of the $`X(3940)`$ is $`\mathrm{\Gamma }=(15.4\pm 10.1)\mathrm{MeV}/c^2`$. The mass of the state is measured to be $`M=(3.943\pm 0.006)\mathrm{GeV}/c^2`$.
We perform a similar fit to the $`M_{\mathrm{recoil}}(J/\psi )`$ distribution for events tagged and constrained as $`e^+e^{}J/\psi D\overline{D}`$. Since no $`X(3940)`$ signal is seen for this mode, we fit this distribution with $`X(3940)`$ parameters fixed to the values found by the fit of tagged $`e^+e^{}J/\psi D^{}\overline{D}`$. The signal yield is found to be $`0.2_{3.5}^{+4.4}`$ events and we set an upper limit for the $`X(3940)`$ signal of 8.1 events at the 90% C.L.
An enhancement with a similar mass, $`Y(3940)`$, decaying into $`J/\psi \omega `$ has been recently observed by Belle choi3 in $`B`$ decays. We perform a search for the decay $`X(3940)J/\psi \omega `$ to see if $`X(3940)`$ and $`Y(3940)`$ could be the same particle. To increase the efficiency we reconstruct the $`\omega `$ and only one $`J/\psi `$ from the $`J/\psi J/\psi \omega `$ final state. The unreconstructed $`J/\psi `$ is identified as a peak in the spectrum of recoil masses against the reconstructed $`J/\psi \omega `$ combinations. A signal for $`X(3940)J/\psi \omega `$ would be seen as a peak around the nominal $`X(3940)`$ mass in a distribution of $`M_{\mathrm{recoil}}(J/\psi )M_{\mathrm{recoil}}(J/\psi \omega )+M(J/\psi )`$ if the reconstructed $`J/\psi `$ is prompt, and in the $`J/\psi \omega `$ invariant mass distribution if the reconstructed $`J/\psi `$ is from the $`X(3940)`$ decay. Since the first case has much larger combinatorial background and less sensitivity, we use only the second case. A scatterplot of $`M_{\mathrm{recoil}}`$($`J/\psi \omega `$) *vs.* $`J/\psi \omega `$ invariant mass in the data is shown in Fig. 4 a). An $`M(J/\psi \omega )`$ projection with the additional requirement $`|M_{\mathrm{recoil}}(J/\psi \omega )M_{J/\psi }|<100\mathrm{MeV}/c^2`$ is shown in Fig. 4 b). A fit to this distribution is done with the signal function and parameters fixed from the result of the $`D^{}\overline{D}`$ tagged fit; the background is a threshold function. The fit yields $`1.9_{2.4}^{+3.2}`$ signal events corresponding to a $`7.4`$ event upper limit at the 90% C.L.
The systematic errors for the $`e^+e^{}J/\psi X(3940)`$ Born cross section and for the $`X(3940)`$ branching fractions $`(X(3940))`$ are summarized in Table 2.
To estimate the systematic errors associated with the fitting procedure we study the difference in $`X(3940)`$ yield returned by the fit to the $`M_{\mathrm{recoil}}`$$`(J/\psi )`$ distribution under different assumptions for the signal and background parameterization. In particular, in the first fit we use a background function that includes several threshold functions corresponding to the production of $`D^{}\overline{D}`$ and $`D^{}\overline{D}^{}`$. We also use the threshold function $`(A+BM_{\mathrm{recoil}}(J/\psi ))\sqrt{M_{\mathrm{recoil}}(J/\psi )M_{\mathrm{thr}}}`$. Different angular distributions result in different $`J/\psi `$ (and $`D`$) reconstruction efficiencies. In the MC the $`J/\psi `$ production angle and $`J/\psi ,X(3940)`$ helicity angle distributions are assumed to be flat. The possible extreme angular distributions ($`1+\mathrm{cos}^2\theta `$ and $`\mathrm{sin}^2\theta `$) are considered to estimate the systematic uncertainty of this assumption. This uncertainty partially cancels out in the determination of $`(X(3940))`$ because of the common $`J/\psi `$ efficiency. Other contributions come from $`N_{\mathrm{ch}}>4`$ requirement efficiency, track reconstruction efficiency, lepton identification for reconstructed $`J/\psi `$ and kaon identification for reconstructed $`D`$.
The systematic errors in the measurement of the $`X(3940)`$ mass are dominated by the $`5\mathrm{MeV}/c^2`$ uncertainty associated with the fitting procedure. The uncertainty due to the $`J/\psi `$ momentum scale is less than $`3\mathrm{MeV}/c^2`$ 2cc2 . These contributions are added in quadrature to give $`6\mathrm{MeV}/c^2`$ total. The systematic error on the mass of the nearby $`\eta _c(2S)`$, which is still poorly known pdg , is found to be the same. From the fit to Fig. 3 b) we estimate the $`X(3940)`$ width to be smaller than $`47\mathrm{MeV}/c^2`$ at the 90% C.L.; this takes into account the fact that the likelihood function is not parabolic. When fitting systematics are taken into account, we find $`\mathrm{\Gamma }<52\mathrm{MeV}/c^2`$ at the 90% C.L.
The Born cross section for $`e^+e^{}J/\psi X(3940)`$ is calculated following the procedure used in Ref. 2cc2 . As in the previous Belle papers, because of selection criteria the result is presented in terms of the product of the cross section and the branching fraction of the $`X(3940)`$ into more than two charged tracks ($`_{>2}`$). We obtain
$`\sigma _{\text{Born}}\times _{>2}=(10.6\pm 2.5\pm 2.4)\mathrm{fb}.`$ (2)
Using the $`X(3940)`$ yields in inclusive and $`D^{}\overline{D}`$ tagged samples, we calculate the fraction of $`X(3940)`$ decays with more than two charged tracks in the final state into $`D^{}\overline{D}`$, $`_{>2}(X(3940)D^{}\overline{D})`$. To remove the correlation between the two samples, we apply a veto on $`D^{}\overline{D}`$ tagging in the inclusive sample. Correcting for the tagging and veto efficiencies obtained from MC with equal fractions of $`X(3940)D^0\overline{D}^0`$ and $`X(3940)D^+D^{}`$ assumed, we find
$`_{>2}(X(3940)D^{}\overline{D})=(96_{32}^{+45}\pm 22)\%`$
$`(>45\%\text{ at 90\% C.L.}),`$ (3)
where the systematic errors are taken into account for the lower limit. In the limit of a vanishing fraction of low charged multiplicity $`X(3940)`$ decays, the measured value of $`_{>2}`$ corresponds to $`(X(3940)D^{}\overline{D})`$.
We set upper limits on the branching fractions of decay of $`X(3940)`$ into $`D\overline{D}`$ and $`X(3940)J/\psi \omega `$ final states, taking into account the estimated systematic errors:
$`(X(3940)D\overline{D})<41\%\text{ at 90\% C.L.;}`$ (4)
$`(X(3940)J/\psi \omega )<26\%\text{ at 90\% C.L.}`$ (5)
These limits assume that low charged multiplicity $`X(3940)`$ decays are negligible and, thus, may be overestimated.
In summary, we have observed a new charmonium-like state, $`X(3940)`$, produced in the process $`e^+e^{}J/\psi X(3940)`$, both in inclusive production and via the $`X(3940)D^{}\overline{D}`$ decay mode. Both observations have a $`5\sigma `$ statistical significance. We have measured the Born cross section for the production process, the branching fraction for $`X(3940)D^{}\overline{D}`$, and set upper limits on $`X(3940)`$ decays to $`D\overline{D}`$ and $`J/\psi \omega `$.
We thank the KEKB group for the excellent operation of the accelerator, the KEK cryogenics group for the efficient operation of the solenoid, and the KEK computer group and the NII for valuable computing and Super-SINET network support. We acknowledge support from MEXT and JSPS (Japan); ARC and DEST (Australia); NSFC (contract No. 10175071, China); DST (India); the BK21 program of MOEHRD and the CHEP SRC program of KOSEF (Korea); KBN (contract No. 2P03B 01324, Poland); MIST (Russia); MHEST (Slovenia); SNSF (Switzerland); NSC and MOE (Taiwan); and DOE (USA). |
warning/0507/astro-ph0507587.html | ar5iv | text | # MOLECULAR HYDROGEN IN INFRARED CIRRUS
## 1 INTRODUCTION
In an attempt to understand the molecular content and physical characteristics of interstellar gas in the low Galactic halo, we exploit infrared and ultraviolet data from two NASA satellites: the Infrared Astronomical Satellite (IRAS) mission of 1983 and the Far Ultraviolet Spectroscopic Explorer (FUSE) satellite of 1999–2005. The combination of infrared emission and ultraviolet absorption along sight lines to 45 active galactic nuclei (AGN) allows us to correlate the infrared cirrus emission intensity with the molecular hydrogen (H<sub>2</sub>) absorption column density at select locations at high Galactic latitude.
IRAS mapped the sky in four infrared bands centered on 12, 25, 60, and 100 $`\mu `$m. Low et al. (1984) introduced one of the most surprising results from the IRAS maps: diffuse filamentary dust clouds that pervade our Galaxy at high latitudes, even in the direction of the Galactic poles. These “infrared cirrus” clouds are most prominent at long wavelengths, 100 $`\mu `$m, but some can be seen in the 60 $`\mu `$m, 25 $`\mu `$m, and 12 $`\mu `$m bands. Comparisons between IRAS dust maps and maps of 21 cm emission reveal a generally good correlation between neutral hydrogen (Hartmann & Burton 1997) and cirrus dust filaments (Figure 1). Because molecular hydrogen (H<sub>2</sub>) forms catalytically on the surface of dust grains (Hollenbach, Werner, & Salpeter 1971), with significant formation rates for grain temperatures $`T_{\mathrm{gr}}<100`$ K and gas temperatures $`T_{\mathrm{gas}}<300`$ K (Hollenbach & McKee 1979; Shull & Beckwith 1982), the cold, dusty conditions of the infrared cirrus clouds are favorable for the formation of H<sub>2</sub>. Thus, it is plausible that some fraction of the hydrogen atoms in the cirrus clouds are bound into molecules.
Previously, the presence of H<sub>2</sub> in infrared cirrus was inferred indirectly. First, under the assumption that the infrared emission and total hydrogen column density, N<sub>H</sub>, are proportional, regions of high dust/H I ratio, termed “infrared excess”, were attributed to the presence of H<sub>2</sub> (de Vries, Heithausen, & Thaddeus 1987; Desert, Bazell, & Boulanger 1988; Reach, Koo, & Heiles 1994; Moritz et al. 1998; Schlegel, Finkbeiner, & Davis 1998, hereafter SFD98). Second, the detection of CO in dense cirrus clouds suggests that the diffuse cirrus clouds should contain H<sub>2</sub> as well. Weiland et al. (1986) compared CO maps from Magnani, Blitz, & Mundy (1985) to IRAS maps of infrared cirrus. Each of the 33 CO clouds had a cirrus counterpart with similar morphology. This work established that at least some of the infrared cirrus cloud cores contains CO gas.
Unfortunately, there is currently no experiment that can map diffuse H<sub>2</sub> emission, either in the 2.12 $`\mu `$m \[(1–0) S(1)\] vibrational line or in the S(0), S(1), S(2) pure rotational lines at 28 $`\mu `$m, 17 $`\mu `$m, and 12 $`\mu `$m, respectively. Although H<sub>2</sub> is over $`10^4`$ times more abundant than CO, the ultraviolet and infrared fluorescent emission of H<sub>2</sub> is very weak. Ultraviolet absorption-line spectroscopy is therefore the primary means for detecting cold H<sub>2</sub> in diffuse clouds. However, it requires background sources with sufficient UV flux to provide adequate signal-to-noise ratio (S/N) to detect the weak H<sub>2</sub> lines. The first major project to conduct such observations was the Copernicus mission of the 1970s (Spitzer & Jenkins 1975). However, its sensitivity limited the possible background sources to early-type stars within about 500 pc of the Sun. Most OB stars that fit this criterion are at low Galactic latitude, and they suffer from confusion and dust extinction in the Galactic plane. Individual features in the infrared cirrus cannot be discerned at low Galactic latitudes, and these stellar sight lines are not effective probes of the dusty filaments.
The FUSE satellite, which has been observing the ultraviolet sky since 1999, has expanded the opportunities for detecting H<sub>2</sub>. The increased sensitivity of FUSE ($`m_{\mathrm{lim}}15.5`$ mag) over Copernicus ($`m_{\mathrm{lim}}8`$ mag) allows us to use more distant stars as well as active galactic nuclei (AGN) as background sources. Our FUSE survey of H<sub>2</sub> toward high-latitude AGN (Gillmon et al. 2005) is particularly well suited for probing H<sub>2</sub> in infrared cirrus. The high-latitude sight lines avoid the confusion of the Galactic disk, and they provide long path lengths through the Galactic halo. In addition, the random distribution of AGN on the sky samples a range of infrared cirrus emission intensities. The main limitation of using “pencil-beam” (absorption) sight lines to detect H<sub>2</sub> in infrared cirrus is the inability to determine whether the gas and dust detected along a given sight line are physically associated. Therefore, we must rely on indirect correlations between cirrus and H<sub>2</sub> absorption.
In this paper, we compare the H<sub>2</sub> column densities in the survey by Gillmon et al. (2005) with the corresponding infrared cirrus fluxes (SFD98). By establishing a correlation between the two, we assert that at least some of the detected H<sub>2</sub> resides in the cirrus clouds. In § 2 we describe the data acquisition and analysis for both IRAS and FUSE. In § 3 we compare the cirrus emission and H<sub>2</sub> absorption and discuss the correlation of the two. Exploiting this correlation and summing over the distribution of H<sub>2</sub> column densities with 100 $`\mu `$m cirrus intensity, we estimate the total H<sub>2</sub> mass ($`10^7M_{}`$) in cirrus clouds around the Milky Way. With molecular fractions ranging from 1–30%, the total cirrus mass throughout the Milky Way is $`10^8M_{}`$. In § 4 we summarize our results and the implications of finding this amount of gas in the low halo of the Milky Way.
## 2 DATA ACQUISITION AND ANALYSIS
### 2.1 FUSE Observations
The H<sub>2</sub> absorption data were taken from FUSE spectra of AGN, using standard data analysis techniques (Tumlinson et al. 2002; Gillmon et al. 2005). Studies of H<sub>2</sub> have been a major part of the FUSE science plan. The satellite, its mission, and its on-orbit performance are described in Moos et al. (2000) and Sahnow et al. (2000). Scientific results on interstellar H<sub>2</sub> have appeared in a number of papers (Shull et al. 2000; Snow et al. 2000; Rachford et al. 2002; Richter et al. 2001, 2003; Tumlinson et al. 2002; Shull et al. 2005). The resolution of FUSE varies from $`R=\lambda /\mathrm{\Delta }\lambda =15,00020,000`$ across the far-UV band. All observations were obtained in time-tag (TTAG) mode, using the $`30^{\prime \prime }\times 30^{\prime \prime }`$ LWRS aperture, with resolution $`20`$ km s<sup>-1</sup> at 1050 Å. The S/N of the co-added data ranges from 2–11 per pixel; the S/N per resolution element varies with spectral resolution, which is not fixed in our survey. Most of the data were binned by 4 pixels before analysis, with the rare case of binning by 2 or 8 pixels.
The data in this paper were taken from our high-latitude H<sub>2</sub> survey (Gillmon et al. 2005), which describes our search for H<sub>2</sub> absorption along 45 sight lines to background AGN at Galactic latitudes $`|b|>20^{}`$. The 45 AGN in the survey by Gillmon et al. (2005) were a subset of the 219 FUSE targets selected in Wakker et al. (2003) as candidates for the analysis of Galactic O VI. These targets probe diffuse gas in both the local Galactic disk and low Galactic halo. Of the available AGN at high latitude, 45 sight lines were chosen, based on an imposed S/N requirement of $`(\mathrm{S}/\mathrm{N})_{\mathrm{bin}}>4`$ or $`(\mathrm{S}/\mathrm{N})_{\mathrm{pix}}>2`$ with 4-pixel binning. The observed H<sub>2</sub> lines arise from the Lyman and Werner electronic transitions, from the ground electronic state, X $`{}_{}{}^{1}\mathrm{\Sigma }_{g}^{+}`$, to the excited states, B $`{}_{}{}^{1}\mathrm{\Sigma }_{u}^{+}`$ (Lyman bands) and C $`{}_{}{}^{1}\mathrm{\Pi }_{u}^{}`$ (Werner bands). The rotational-vibrational lines arise from the ground vibrational state and a range of rotational states. Our analysis was restricted to ten vibrational-rotational bands, Lyman (0–0) to (8–0) and Werner (0–0), which are located between 1000 $`\mathrm{\AA }`$ and 1126 $`\mathrm{\AA }`$. The vibrational state notation is ($`v_{\mathrm{upper}}v_{\mathrm{lower}}`$). In most sight lines, we observed absorption lines from rotational states $`J=0`$–3, and sometimes up to $`J=4`$.
The end product of the H<sub>2</sub> absorption-line analysis is the column density, N<sub>H2</sub> (cm<sup>-2</sup>), the physical density of H<sub>2</sub> molecules, integrated along the sight line. Each absorption line was fitted with a Voigt profile in order to determine the equivalent width, a measure of the absorbed light in the line. The equivalent widths were fitted to a curve of growth, to find the column density, N$`{}_{\mathrm{H2}}{}^{}(J)`$ in each rotational state $`J`$. The points for each $`J`$ were tied together during the fitting to produce a single, consistent column density for each rotational state, N($`J`$). The sum of the column densities in all rotational states then gives the total column density, N<sub>H2</sub>. FUSE spectra of 87% (39 of 45) of the observed AGN showed detectable H<sub>2</sub> absorption, with column densities ranging from N<sub>H2</sub> $`=10^{14.1719.82}`$ cm<sup>-2</sup>. The FUSE survey is sensitive to N<sub>H2</sub> $`>10^{13.814.6}`$ cm<sup>-2</sup>, depending on the S/N (2–11 per pixel) and spectral resolution ($`R=15,000`$–20,000).
### 2.2 IRAS Dust Maps
To obtain infrared cirrus emission intensities (MJy sr<sup>-1</sup>), we use the 100 $`\mu `$m maps presented by SFD98. These were a composite of data from the IRAS mission of 1983 and the COBE mission of 1989–1990, capitalizing on the strengths of each. Because the interstellar dust emits like a “grey body”, the emission intensity is sensitive to the dust temperature. As a result, two regions with the same dust column density but different dust temperatures will have different infrared intensities. To correct for this effect, SFD98 used the ratio of the 100 $`\mu `$m and 240 $`\mu `$m COBE maps to produce a map of dust color temperature. They used this ratio to correct the 100 $`\mu `$m IRAS map so that it is proportional to dust column density. IRAS mapped the sky in four broadband infrared channels, centered at 12, 25, 60, and 100 $`\mu `$m with a resolution of $`5^{}`$, while COBE mapped the sky in 10 broad photometric bands from 1 to 240 $`\mu `$m at a resolution of $`0.7^{}`$. Before combining the data sets, SFD98 took great care in the difficult removal of the zodiacal foreground emission and IRAS striping artifacts that arise from differences in solar elongation between scans. Confirmed point sources were also removed. The maps were combined in such a way as to preserve the COBE calibration and IRAS resolution.
Maps of the temperature-corrected 100 $`\mu `$m intensity, $`D_{100}^{(T)}`$, are presented for the northern Galactic hemisphere (Figure 2) and for the southern Galactic hemisphere (Figure 3). We overplot the locations of the 45 sight lines from the FUSE H<sub>2</sub> survey: 28 northern AGN and 17 southern AGN.
## 3 Comparisons of H<sub>2</sub> and Infrared Cirrus
### 3.1 The H<sub>2</sub> Self-Shielding Transition in N<sub>H</sub>
Before the FUSE mission, the direct detection of H<sub>2</sub> in infrared cirrus by UV absorption-line spectroscopy was prevented mainly by a lack of background sources at high Galactic latitude. Even though this problem has been alleviated by FUSE’s ability to observe bright AGN as background sources, another problem with absorption-line spectroscopy along “pencil-beam” sight lines comes to the forefront. Absorption studies lack morphological information and cannot provide distances to clouds along the line of sight. Therefore, it is difficult to identify the gas that gives rise to the detected column density. The infrared cirrus problem is a classic case. For any given sight line, it is possible that the H<sub>2</sub> absorption is not physically associated with the infrared cirrus along the beam. If this is the case, then comparing N<sub>H2</sub> with the cirrus dust column density would lead to erroneous results. It would be helpful to determine whether the diffuse H<sub>2</sub> along all sight lines resides in the cirrus. If no other component of the diffuse ISM harbored significant amounts of H<sub>2</sub>, then N<sub>H2</sub> for a given sight line could safely be associated with other cirrus properties.
In this section, we investigate this possibility, based on a property of H<sub>2</sub> called “self-shielding”. Following line absorption of UV photons from the mean interstellar radiation field, H<sub>2</sub> decays approximately 11% of the time to the dissociative continuum. As a result, the molecular fraction,
$$f_{\mathrm{H2}}=\frac{2\mathrm{N}_{\mathrm{H2}}}{\mathrm{N}_{\mathrm{HI}}+2\mathrm{N}_{\mathrm{H2}}}\frac{2N_{\mathrm{H2}}}{\mathrm{N}_\mathrm{H}},$$
(1)
is generally larger in clouds with a greater total hydrogen column density, N<sub>H</sub> = N<sub>HI</sub> \+ 2N<sub>H2</sub>. Molecules on the outside of the cloud shield those in the interior from dissociating UV (Hollenbach, Werner, & Salpeter 1971; Black & Dalgarno 1976; Browning, Tumlinson, & Shull 2003).
In optically thin clouds, the density of molecules can be approximated by the equilibrium between formation and destruction,
$$f_{\mathrm{H2}}\frac{2n_HR(T_{\mathrm{gas}},T_{\mathrm{gr}},Z)}{\beta f_{\mathrm{diss}}}(10^5)R_{17}n_{30}\left(\frac{\beta _0}{\beta }\right).$$
(2)
In this formula, the numerical value for $`f_{\mathrm{H2}}`$ is scaled to fiducial values of hydrogen density, $`n_H`$ (30 cm<sup>-3</sup>), H<sub>2</sub> formation rate coefficient, $`R`$ ($`10^{17}`$ cm<sup>3</sup> s<sup>-1</sup>), and mean H<sub>2</sub> pumping rate in the FUV Lyman and Werner bands, $`\beta _0=5\times 10^{10}`$ s<sup>-1</sup>. The H<sub>2</sub> photodissociation rate is written as $`f_{\mathrm{diss}}\beta `$, where the coefficient $`f_{\mathrm{diss}}0.11`$ is the average fraction of FUV excitations of H<sub>2</sub> that result in decays to the dissociating continuum. The H<sub>2</sub> formation rate per unit volume is written as $`n_Hn_{\mathrm{HI}}R`$, where the coefficient $`R`$ depends on the gas temperature, grain surface temperature, and gas metallicity ($`Z`$). The metallicity dependence comes from the assumed scaling of grain-surface catalysis sites with the grain/gas ratio. For sight lines in the local Galactic disk, this rate coefficient has been estimated (Jura 1974) to range from $`R=(13)\times 10^{17}`$ cm<sup>3</sup> s<sup>-1</sup> at solar metallicity. This standard value for $`R`$ is expected to apply at suitably low temperatures of gas ($`T_{\mathrm{gas}}300`$ K) and grains ($`T_{\mathrm{gr}}100`$ K) as discussed by Shull & Beckwith (1982) and Hollenbach & McKee (1979).
The effects of self-shielding were observed by plotting $`f_{\mathrm{H2}}`$ versus N<sub>H</sub>, the total column density of hydrogen. As H<sub>2</sub> absorption lines in the Lyman and Werner bands become optically thick, the rate ($`\beta `$) of UV pumping and molecular dissociation diminish. In this way, the presence of H<sub>2</sub> screens molecules in the inner portions of the cloud from dissociation. A transition from low ($`f_{\mathrm{H2}}10^5`$) to high ($`f_{\mathrm{H2}}>10^2`$) molecular fractions at N$`{}_{\mathrm{H}}{}^{}5\times 10^{20}`$ cm<sup>-2</sup> was noted in the Copernicus H<sub>2</sub> survey (Savage et al. 1977). A recent FUSE survey of Galactic disk stars (Shull et al. 2005) provided similar results (Figure 4).
In assessing the molecular content of the cirrus, we consider three generic cases. The first possibility (Case I) is that all the detected H<sub>2</sub> resides in cirrus clouds, with no contribution from foreground clouds. Under the assumption that the total hydrogen density, $`n_\mathrm{H}`$, is proportional to the number density of dust grains in any given interstellar cloud (Hollenbach & McKee 1979), the cirrus dust column density should be proportional to the amount of N<sub>H</sub> associated with H<sub>2</sub>. A plot of N<sub>H2</sub> vs. dust column density should then show the familiar self-shielding transition of molecular hydrogen. In Case II, some of the observed H<sub>2</sub> lies in the cirrus, but some resides in another component of the ISM not visible in the cirrus maps. A plot of N<sub>H2</sub> vs. dust column density would show some points that follow the self-shielding transition (those corresponding to cirrus). However, the transition would be broadened by points that are randomly distributed. If a significant portion of the detected H<sub>2</sub> along a sight line is not in cirrus, N<sub>H2</sub> will not correlate with dust column density. This effect would produce sight lines with low dust column density and significant amounts of H<sub>2</sub>. The extreme situation is Case III, in which none of the H<sub>2</sub> resides in cirrus clouds. We consider this unlikely, given the number of observed regions with “infrared excess” (high dust/N<sub>HI</sub> ratios) and the detection of CO in denser cirrus clouds. Once the N<sub>H2</sub> in cirrus exceeds $`10^{14.5}`$ cm<sup>-2</sup>, it should be detectable by FUSE. For Case III to be consistent with our survey, the covering fraction of cirrus regions with detectable H<sub>2</sub> would need to be quite small, whereas the cirrus covering factor is observed to be $`50`$% at $`b>30^{}`$ (see § 3.4).
The issue then comes down to whether Case I or Case II is a better description of the cirrus-H<sub>2</sub> correlation. Figure 5 presents a plot of N<sub>H2</sub> along sight lines to the 45 AGN in the FUSE survey vs. the temperature-corrected IR flux at each location from the SFD98 maps (blue diamonds). If the temperature-corrected flux is proportional to dust column density, this plot can be used to test the three cases mentioned above. The blue diamonds appear to show the self-shielding transition discussed in Case I, implying that much of the detected H<sub>2</sub> is in the cirrus.
### 3.2 AGN Behind Regions of Low Cirrus
In this section, we describe five additional sight lines shown in red (Figure 5) toward the regions of lowest cirrus. These additional sight lines were chosen to explore whether the correlation between N<sub>H2</sub> and 100 $`\mu `$m cirrus is universal. Table 5 of SFD98 lists the coordinates of the regions of lowest (temperature-corrected) 100 $`\mu `$m intensity, $`D_{100}^{(T)}`$, from their maps. We searched the FUSE archive for targets within $`5^{}`$ of the given coordinates and found five targets behind regions of low cirrus ($`D_{100}^{(T)}0.5`$ MJy sr<sup>-1</sup>) that also had FUSE data with sufficient (S/N)$`{}_{\mathrm{pix}}{}^{}>2`$ to conduct an H<sub>2</sub> analysis. These five targets are in addition to those in the Gillmon et al. (2005) sample of 45 AGN. They are listed in Table 1 and shown as asterisks in Figure 5.
Four of these five targets showed no evidence of H<sub>2</sub>, with a typical upper limit of N<sub>H2</sub> $`10^{14.5}`$ cm<sup>-2</sup>. This result lends further credence to the theory that most of the observed H<sub>2</sub> is in the cirrus (Case I). However, one target, UGC 5720, showed significant H<sub>2</sub> absorption lines (Figure 6). An analysis, as described in § 2.1, yielded a significant column density N<sub>H2</sub> $`=10^{18.79\pm 0.05}`$ cm<sup>-2</sup>. In § 3.3, we explore possible explanations for this anomalous sight line and for the spread in the self-shielding transition.
### 3.3 The H<sub>2</sub> Transition seen in Cirrus
Figure 5 plots N<sub>H2</sub> vs. temperature-corrected intensity, $`D_{100}^{(T)}`$ (MJy sr<sup>-1</sup>), and exhibits a clear self-shielding transition of H<sub>2</sub>. This indicates that a significant fraction of the detected H<sub>2</sub> resides in the infrared cirrus. However, with small-number statistics on 50 AGN sight lines, the transition is not sharp, occurring between log $`D_{100}^{(T)}`$ = 0.2–0.5. There is one outlying point (UGC 5720) with low cirrus intensity, log $`D_{100}^{(T)}=0.19`$, but a significant column density, log N<sub>H2</sub> $`=18.79\pm 0.05`$. This outlier suggests that not all the detected H<sub>2</sub> resides in cirrus clouds, and that other components of the diffuse gas may harbor detectable amounts of molecules.
In some sight lines that intercept non-cirrus clouds with low 100 $`\mu `$m intensity, the gas may undergo an early molecular transition and appear with higher N<sub>H2</sub>. There are other explanations that might produce the same effects, even if all the H<sub>2</sub> resides in cirrus clouds. One consideration is how well the temperature-corrected flux map traces the dust column density. To correct for temperature variations, SFD98 created a temperature map from the ratio of the COBE 100 $`\mu `$m and 240 $`\mu `$m maps. Thus, low-resolution (1.1) temperature maps were used to correct the $`5^{}`$ resolution IRAS map to produce a 100 $`\mu `$m map at $`6.^{}1`$ resolution. It is likely that the dust temperature varies on smaller scales than can be resolved by this method, and a more accurate temperature correction might tighten the self-shielding transition. It is also possible that UGC 5720 lies behind a region of cold cirrus dust unresolved by the temperature map. Thus, a significant portion of the dust column density might not be indicated by the observed flux.
A related consideration is that the IRAS map might not resolve significant variations in dust column density. The sight lines observed by FUSE absorption probe gas along a very narrow beam, while the IRAS beam is considerably larger. Thus, if there were significant structure on smaller scales than the IRAS beam, FUSE could observe a denser clump of H<sub>2</sub> while the IRAS emission would be “beam-diluted”. A higher resolution infrared map might tighten the self-shielding transition and/or bring the outlier UGC 5720 onto the correlation. Figure 7 shows a $`2^{}\times 2^{}`$ section of the SFD98 flux map centered on UGC 5720. Large fluctuations in temperature-corrected flux near the position of the AGN suggest that the sight line may be picking up a clump of cirrus with higher dust column density unresolved by these maps.
Finally, there is the possibility that there are multiple cirrus features superimposed along a line of sight. Duel & Burton (1990) compared the morphology of cirrus clouds and H I maps in various velocity intervals to show that cirrus features that appear simple are, in some cases, superpositions of kinematically distinct components. This idea of a “concatenation of clouds” was proposed for translucent H<sub>2</sub> clouds (Browning et al. 2003) to explain the high levels of H<sub>2</sub> rotational excitation in systems with large N<sub>H2</sub>. These authors suggested that, if a sight line intercepts multiple, physically distinct cloud components, the H<sub>2</sub> will be exposed to a radiation field enhanced over that expected for a single, contiguous cloud with the same total column density. The enhanced H<sub>2</sub> destruction rate from the stronger UV radiation would reduce the mean molecular fraction and produce a gradual transition to higher N<sub>H2</sub>. Some sight lines with seemingly high dust column density, but low N<sub>H2</sub>, may actually be probing multiple, superimposed filaments with high integrated dust column density.
### 3.4 Mass of the Cirrus Clouds
With the empirical correlation (Figure 5) between IR cirrus and H<sub>2</sub> column density, we can make a quantitative estimate of the H<sub>2</sub> mass in the diffuse cirrus clouds. We begin with the IRAS maps of the northern Galactic hemisphere. We assume that these cirrus clouds lie at elevation $`z(100\mathrm{pc})z_{100}`$ above the Milky Way plane, and that clouds of intensity $`D_i`$ cover a fraction $`f_c(D_i)`$ of the planar area at $`b30^{}`$. The total H<sub>2</sub> mass in this planar cloud deck is then,
$$M_{\mathrm{H2}}=(\pi R^2)(2m_{\mathrm{H2}})\underset{i}{}\mathrm{N}_{\mathrm{H2}}(D_i)f_c(D_i)\mathrm{\Delta }D_i(2400M_{})\left[\frac{\mathrm{N}_{\mathrm{H2}}}{10^{18.5}\mathrm{cm}^2}\right]\left[\frac{f_c}{0.5}\right]z_{100}^2,$$
(3)
where N<sub>H2</sub>$`(D_i)`$ is the mean H<sub>2</sub> column density corresponding to cirrus intensity $`D_i`$ (Figure 5) and where $`R=z/\mathrm{tan}(b)(173\mathrm{pc})z_{100}`$ is the radius of the cirrus disk, at elevation $`z`$ subtended by the cone at $`b=30^{}`$. For this estimate, we have assumed that the cirrus covers a fraction $`f_c0.5`$ of the sky at $`b30^{}`$, independent of intensity.
We now make a more careful calculation, summing over the actual data from IRAS and FUSE. Figure 8 shows the distribution of cirrus covering factors, $`f_c(D_i)`$, averaged over northern hemisphere regions with $`b`$ greater than $`30^{}`$, $`40^{}`$, and $`50^{}`$, respectively. For our fiducial cone at $`b30^{}`$, approximately 50% of the sky is covered by cirrus with intensity log $`D_{100}^{(T)}0.2`$, the value corresponding to the H<sub>2</sub> transition (Figure 5). Performing the full summation (eq. 3) over 8 logarithmic intensity bins above the transition, of width $`\mathrm{\Delta }`$(log $`D_i`$) = 0.1, between log $`D_i=0.2`$ and 1.0, we find a total H<sub>2</sub> mass of $`(2600M_{})z_{100}^2`$ contained in the cirrus at $`b30^{}`$.
To convert this calculation to the inner Milky Way, within the solar circle, we multiply by a factor $`2`$, for the northern and southern Galactic hemispheres, and scale by a factor $`(8.2\mathrm{kpc}/0.173\mathrm{kpc})^22250z_{100}^2`$ to account for the number of similar conical areas around the Galactic disk. Note that this total H<sub>2</sub> mass is independent of the assumed cirrus elevation, $`z`$, since the area-scaling cancels the factor $`z_{100}^2`$ in equation (3). We arrive at an extrapolated total molecular mass over the inner Milky Way, $`M_{\mathrm{H2}}^{\mathrm{cirrus}}10^7M_{}`$, assuming that the cirrus along the AGN sightlines at $`b>30^{}`$ is typical. The molecular fractions of these cirrus clouds, with N<sub>H2</sub> $`10^{18.5}`$ cm<sup>-2</sup>, range from 1–30% for log $`D_{100}^{(T)}`$ = 0.2–0.5 (see Figure 6 of Gillmon et al. 2005), with an average $`f_{\mathrm{H2}}=0.1`$ for clouds with log N<sub>H2</sub> $`18.5`$. Therefore, we estimate the total gas mass in the cirrus to be $`10^8M_{}`$.
The characteristics of the cirrus clouds along our AGN sight lines can be verified by computing the “dust-to-gas” ratio, defined as the ratio of 100 $`\mu `$m cirrus intensity, $`D_{100}^{(T)}`$ (MJy sr<sup>-1</sup>), to H I column density, N<sub>HI</sub> (cm<sup>-2</sup>). Table 2 gives these values and their ratio for 16 of our AGN sight lines observed in 21 cm (Lockman & Savage 1995). This ratio ranges from $`0.66\times 10^{20}`$ MJy sr<sup>-1</sup> cm<sup>2</sup> toward PG 0953+414 to $`4.3\times 10^{20}`$ MJy sr<sup>-1</sup> cm<sup>2</sup> toward 3C 273. The mean value and standard deviation are $`(1.43\pm 0.41)\times 10^{20}`$ MJy sr<sup>-1</sup> cm<sup>2</sup>, when we exclude the anomalous sight line to 3C 273, which lies behind Radio Loop I and the North Polar Spur. This mean ratio is in excellent agreement with the mean value, $`(1.4\pm 0.3)\times 10^{20}`$ MJy sr<sup>-1</sup> cm<sup>2</sup>, found by Boulanger, Baud, & van Albada (1985) in a $`20^{}\times 18^{}`$ field at high Galactic latitude.
## 4 SUMMARY AND FUTURE WORK
We have undertaken a comparison between the column density, N<sub>H2</sub>, and the 100 $`\mu `$m cirrus intensity for a total of 50 sight lines. For the cirrus, we used the temperature-corrected maps of Schlegel, Finkbeiner, & Davis (1998), and adopted H<sub>2</sub> column densities from our high-latitude survey (Gillmon et al. 2005). The presence of a clear correlation between UV (H<sub>2</sub>) absorption and IR (cirrus) emission indicates that a significant fraction of the H<sub>2</sub> is physically associated with the cirrus clouds. However, the self-shielding transition of H<sub>2</sub> fraction used to define the correlation is not sharp. The existence of one outlying sight line suggests either that some of the detected H<sub>2</sub> may exist in another component of the diffuse ISM, or that the limited resolution of the infrared maps is obscuring the physical conditions. Of the three possible cases for H<sub>2</sub>–cirrus connections laid out in § 3.1, Case I or II best describe the data.
Put simply, H<sub>2</sub> is contained in most, if not all diffuse cirrus clouds. At Galactic latitudes $`b>30^{}`$, approximately 50% of the sky is covered with cirrus, at temperature-corrected 100 $`\mu `$m intensities $`D_{100}^{(T)}1.5`$ MJy sr<sup>-1</sup>. With this correlation, we have found a convenient means of identifying the best extragalactic sight lines for “H<sub>2</sub>-clean” far-UV absorption studies of intergalactic or interstellar matter. Conversely, if the goal is to study H<sub>2</sub> at the disk-halo interface, the cirrus maps would be a good guide.
We also made a rough estimate of the H<sub>2</sub> mass contained in these cirrus clouds. Exploiting the H<sub>2</sub>–cirrus correlation, we summed the distributions of IR cirrus intensity and H<sub>2</sub> column density to find $`10^7M_{}`$ in cirrus H<sub>2</sub> and $`10^8M_{}`$ in total hydrogen, distributed over the Milky Way disk-halo interface, within the solar circle. Above the self-shielding transition, these diffuse halo clouds have molecular fractions ranging from 1–30%, for column densities N$`{}_{\mathrm{H}}{}^{}10^{20.4}`$ cm<sup>-2</sup> and N<sub>H2</sub> $`10^{18.5}`$ cm<sup>-2</sup> (Gillmon et al. 2005). To support such high molecular fractions by H<sub>2</sub> formation on grain surfaces, the cirrus clouds are probably compressed sheets with densities $`n_H30`$ cm<sup>-3</sup>, in which the gas and grains remain sufficiently cold to form H<sub>2</sub>. On average, the IR cirrus clouds may actually have higher molecular fractions, on average, than diffuse clouds in the Galactic disk. This point was also made by Reach et al. (1994), who estimated that the H<sub>2</sub>/H I transition in denser cirrus clouds occurs at N$`{}_{\mathrm{H}}{}^{}4\times 10^{20}`$ cm<sup>-2</sup>, on the basis of fits to the far-infrared excess. This transition column density is almost twice that found here, probably because the infrared-excess technique requires larger molecular fractions than used for the UV absorbers ($`f_{\mathrm{H2}}0.01`$).
This initial survey opens up considerable opportunities for future studies of the H<sub>2</sub>–cirrus correlation. Higher resolution infrared maps with more accurate temperature corrections, such as those obtainable with the Spitzer Space Telescope, would greatly improve the effectiveness of this method. A comparison with H I 21-cm maps (e.g., Lockman & Condon 2005) could delineate the gaseous structures associated with the cirrus and measure the dust-to-gas ratios in these diffuse clouds. Expanding the UV sample to include more sight lines to background AGN would improve the statistics of such a small sample. The 45 AGN in the survey by Gillmon et al. (2005) were a subset of the 219 FUSE targets selected in Wakker et al. (2003) as candidates for the analysis of Galactic O VI. The next 50 brightest targets have an average flux of $`6\times 10^{14}`$ erg cm<sup>-2</sup> s<sup>-1</sup> Å<sup>-1</sup>. To achieve a S/N of 3 per pixel with FUSE would require approximately 20 ksec per AGN, for a program total of 1000 ksec. An ultraviolet telescope with sensitivity greater than FUSE is probably necessary for feasible exposure times in a $`50^+`$ target survey.
Another promising avenue for future exploration is to compare the cirrus maps and H<sub>2</sub> absorption lines with other tracers, such as H I, CO and $`\gamma `$-ray emission. On the basis of such comparisons, Grenier, Casandjian, & Terrier (2005) suggest that many interstellar clouds in the solar neighborhood have extensive dark regions that bridge the dense cloud cores to atomic phases. These details are beyond the scope of our current paper. Perhaps the most direct way to investigate the connection between diffuse H<sub>2</sub> and IR cirrus would be to map H<sub>2</sub> in UV or IR fluorescent emission. This method would provide the morphological information lacking in “pencil-beam” UV-absorption sight lines. There is currently no high-resolution experiment that can map diffuse H<sub>2</sub> emission, either in the mid-infrared (28 and 17 $`\mu `$m) or the far-ultraviolet (1000–1100 Å). Intriguing results for H<sub>2</sub> ultraviolet fluorescent emission at $`10^{}`$ resolution may soon be available from the Spectroscopy of Plasma Evolution from Astrophysical Radiation (SPEAR) Mission (Edelstein et al. 2003). With appropriate IR and UV telescopes, these methods could help map the gaseous Galactic halo.
We thank Ken Sembach and Jay Lockman for useful discussions. This work was based in part on data obtained for the Guaranteed Time Team team by the NASA-CNES-CSA FUSE mission operated by the Johns Hopkins University. Financial support to U.S. participants has been provided by NASA contract NAS5-32985. The Colorado group also received FUSE support from NASA grant NAG5-10948 for studies of interstellar H<sub>2</sub>. |
warning/0507/quant-ph0507227.html | ar5iv | text | # Violating Bell Inequalities Maximally for Two 𝑑-Dimensional Systems
## Abstract
We investigate the maximal violation of Bell inequalities for two $`d`$-dimensional systems by using the method of Bell operator. The maximal violation corresponds to the maximal eigenvalue of the Bell operator matrix. The eigenvectors corresponding to these eigenvalues are described by asymmetric entangled states. We estimate the maximum value of the eigenvalue for large dimension. A family of elegant entangled states $`|\mathrm{\Psi }_{\mathrm{app}}`$ that violate Bell inequality more strongly than the maximally entangled state but are somewhat close to these eigenvectors is presented. These approximate states can potentially be useful for quantum cryptography as well as many other important fields of quantum information.
The famous Clauser-Horne-Shimony-Holt (CHSH) inequality Bell ; CHSH for two entangled spin-1/2 particles has always provided an excellent test-bed for experimental verification of quantum mechanics against the predictions of local realism exp . It is well-known that all pure entangled states of 2 dimension violate the CHSH inequality: the maximum quantum violation of $`2\sqrt{2}`$ being often called the Tsirelson’s bound Gisin . In 2000, contrary to previous study, Kaszlikowski *et al* showed numerically, based on linear optimization techniques, that the violations of local realism increase with dimensions for two maximally entangled $`d`$-dimensional systems (qudits) ($`3d9`$) Dago00 . A year later, Durt *et al* extended the analysis to $`d=16`$ under special experimental settingsDurt01 . In the same year, an analytical proof was constructed for two maximally entangled 3-dimensional systems (qutrits) JLC1 . In 2002, two research teams independently developed Bell inequalities for high-dimensional systems: the first one is a Clauser-Horne type (probability) inequality for two qutrits JLC2 ; and the second one is a CHSH type (correlation) inequality to two arbitrary $`d`$-dimensional systems CGLMP , now known as the Collins-Gisin-Linden-Masser-Popoescu (CGLMP) inequalities. The tightness of the CGLMP inequality was demonstrated in Ref. LM . The maximally entangled state of two-qudit reads $`|\mathrm{\Psi }_{\mathrm{mes}}=\frac{1}{\sqrt{d}}_{j=0}^{d1}|jj`$ where $`|j`$ is the orthonormal base in each subsystem. Collins *et al* restricted their investigation to $`|\mathrm{\Psi }_{\mathrm{mes}}`$, because on one hand $`|\mathrm{\Psi }_{\mathrm{mes}}`$ is a simple state, on the other hand there is a very natural viewpiont that maximally entangled states would maximally violate the Bell inequality, just as the CHSH inequality has worked for two-qubit. The results of CGLMP was numerically consistent with Ref. Dago00 . For $`|\mathrm{\Psi }_{\mathrm{mes}}`$, the Tsirelson’s bound or Bell expression $`I_d`$ is given by CGLMP
$`I_d(|\mathrm{\Psi }_{\mathrm{mes}})`$ $`=`$ $`4d{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}(1{\displaystyle \frac{2k}{d1}})({\displaystyle \frac{1}{2d^3\mathrm{sin}^2[\pi (k+\frac{1}{4})/d]}}`$ (1)
$`\text{ }{\displaystyle \frac{1}{2d^3\mathrm{sin}^2[\pi (k1+\frac{1}{4})/d]}}),`$
showing that the maximum quantum violation increases with the dimension $`d`$ (here $`\mathrm{}=[d/21]`$ represents the integer part of $`(d/21)`$. ). When $`d\mathrm{}`$, it is interesting to note that
$`\underset{d\mathrm{}}{lim}I_d(|\mathrm{\Psi }_{\mathrm{mes}})`$ $`=`$ $`{\displaystyle \frac{2}{\pi ^2}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}[{\displaystyle \frac{1}{(k+1/4)^2}}{\displaystyle \frac{1}{(k+3/4)^2}}]`$ (2)
$``$ $`2.96981,`$
a number related to the Catalan’s constant.
Contrary to prevalent belief, Acín *et al* studied the quantum nonlocality of two-qutrit as well as two $`d`$-dimensional systems up to $`d=8`$ and discovered another unexpected result: there existed non-maximally entangled states that lead to greater violation of the CGLMP inequalities compared with maximally entangled statesAcin . This surprising result still lacks of a good intuitive explanation. From the TABLE I of Ref. Acin , one observes that the maximal violation increases with dimension $`d`$, and it reaches 3.1013 for $`d=8`$. This gives rise to a natural question:“ What are the maximal violations of the CGLMP inequalities if one increases the dimension $`d`$, especially when $`d`$ goes to infinity?” The purpose of this Letter is to investigate the maximal violations of the CGLMP inequalities for higher dimensional systems. We also present a family of elegant entangled states $`|\mathrm{\Psi }_{\mathrm{app}}`$ whose corresponding maximal violations are approximate to the real ones.
Before computing the quantum violation, let us first estimate the upper bound under quantum mechanics. Significantly, Ref. Fu has recast the CGLMP inequality into a form that is very similar to the CHSH inequality:
$$I_d=Q_{11}+Q_{12}Q_{21}+Q_{22}2,$$
(3)
where $`Q_{ij}`$ are the correlation functions defined by probabilities in the following way
$$Q_{ij}=\frac{1}{S}\underset{m=0}{\overset{d1}{}}\underset{n=0}{\overset{d1}{}}f^{ij}(m,n)P(A_i=m,B_j=n),$$
(4)
where $`S=(d1)/2`$ is the spin of the particle for the $`d`$-dimensional system, $`f^{ij}(m,n)=SM(\epsilon (ij)(m+n),d)`$; $`\epsilon (x)=1`$ and $`1`$ for $`x0`$ and $`x<0`$ respectively; $`M(x,d)=(x\mathrm{mod}d)`$ and $`0M(x,d)d1`$. Refs. CGLMP and Fu have proved that $`I_d2`$ for hidden variable theory. From Eq. (4), one notes that the extreme values of $`Q_{ij}`$ are $`\pm 1`$ for both local realistic description and quantum mechanics, therefore it is impossible that the maximum quantum violation of $`I_d`$ is larger than 4. Furthermore $`I_d`$ cannot reach 4, because $`Q_{ij}`$ are constrained to each other, if three of them are set to be 1, the fourth must also be 1. Consequently, one can conclude easily from above analysis that the maximal quantum violation of $`I_d`$ is a number between 2 and 4.
To generalize Acín’s approach, we first note that it has been shown that unbiased multiport beam splitter Bellport can be used to test violation of local realism of two maximally entangled qudits. Unbiased $`d`$-port beam splitter is a device with the following property: if a photon enters any of the $`d`$ single input ports, its chances of exit are equally split among the $`d`$ output ports. In fact one can always build the device with the distinguishing trait that the elements of its unitary transition matrix $`T`$ are solely powers of the root of unity $`\gamma =\mathrm{exp}(i2\pi /d)`$, namely $`T_{kl}=\frac{1}{\sqrt{d}}\gamma ^{kl}`$. In front of $`i`$-th input port of the device a phase shifter is placed to change the phase of the incoming photon by $`\varphi (i)`$. These $`d`$ phase shifts, denoted for convenience as a “vector” of phase shifts $`\widehat{\varphi }=(\varphi (0),\varphi (1),\mathrm{},\varphi (d1))`$, are macroscopic local parameters that can be changed by the observer. Therefore, unbiased $`d`$-port beam splitter together with the $`d`$ phase shifters perform the unitary transformation $`U(\widehat{\varphi })`$ with the entries $`U_{kl}=T_{kl}\mathrm{exp}(i\varphi (l))`$. Devices (Bell multiports) endowed with such a matrix were proposed, and readers who are interested in it can refer to Refs. Dago00 Bellport . The approach developed in Acin is related to Bell operator. An arbitrary entangled state of two-qudit reads
$`|\mathrm{\Phi }={\displaystyle \underset{j,j^{}=0}{\overset{d1}{}}}\alpha _{jj^{}}|jj^{}.`$ (5)
The quantum prediction of the joint probability $`P(A_a=k,B_b=l)`$ when $`A_a`$ and $`B_b`$ are measured in the initial state $`|\mathrm{\Phi }`$ is given by
$`P(A_a=k,B_b=l)`$ $`=`$ $`\mathrm{Tr}[(U(\widehat{\varphi }_a)^{}U(\widehat{\phi }_b)^{})\widehat{\mathrm{\Pi }_k}\widehat{\mathrm{\Pi }_l}(U(\widehat{\varphi }_a)U(\widehat{\phi }_b))|\mathrm{\Phi }\mathrm{\Phi }|]`$ (6)
$`=`$ $`{\displaystyle \frac{1}{d^2}}{\displaystyle \underset{j,j^{},m,m^{}=0}{\overset{d1}{}}}\alpha _{jj^{}}\alpha _{mm^{}}^{}e^{i[\varphi _a(j)+\phi _b(j^{})+\frac{2\pi }{d}(jkj^{}l)\varphi _a(m)\phi _b(m^{})\frac{2\pi }{d}(mkm^{}l)]},`$
where $`\widehat{\mathrm{\Pi }_k}=|kk|`$, $`\widehat{\mathrm{\Pi }_l}=|ll|`$ are the projectors for systems $`A`$ and $`B`$, respectively. Substituting Eq. (6) into the CGLMP inequality, one gets the Bell expression for the state $`|\mathrm{\Phi }`$:
$`I_d(|\mathrm{\Phi })=`$ $`{\displaystyle \frac{1}{d^2}}{\displaystyle \underset{j,j^{},m,m^{}=0}{\overset{d1}{}}}\alpha _{jj^{}}\alpha _{mm^{}}^{}{\displaystyle \underset{l=0}{\overset{d1}{}}}e^{i\frac{2\pi }{d}[(jm)(j^{}m^{})]l}`$ (7)
$`\times \{e^{i[\varphi _1(j)\varphi _1(m)+\phi _1(j^{})\phi _1(m^{})]}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}(1{\displaystyle \frac{2k}{d1}})(e^{i\frac{2\pi }{d}k(jm)}e^{i\frac{2\pi }{d}(k+1)(j^{}m^{})})+`$
$`e^{i[\varphi _1(j)\varphi _1(m)+\phi _2(j^{})\phi _2(m^{})]}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}(1{\displaystyle \frac{2k}{d1}})(e^{i\frac{2\pi }{d}k(j^{}m^{})}e^{i\frac{2\pi }{d}(k+1)(jm)})+`$
$`e^{i[\varphi _2(j)\varphi _2(m)+\phi _1(j^{})\phi _1(m^{})]}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}(1{\displaystyle \frac{2k}{d1}})(e^{i\frac{2\pi }{d}(k+1)(j^{}m^{})}e^{i\frac{2\pi }{d}k(jm)})+`$
$`e^{i[\varphi _2(j)\varphi _2(m)+\phi _2(j^{})\phi _2(m^{})]}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}(1{\displaystyle \frac{2k}{d1}})(e^{i\frac{2\pi }{d}k(jm)}e^{i\frac{2\pi }{d}(k+1)(j^{}m^{})})\}.`$
$`I_d(|\mathrm{\Phi })`$ can be expressed as
$`I_d(|\mathrm{\Phi })=\mathrm{Tr}(\widehat{B}|\mathrm{\Phi }\mathrm{\Phi }|)=\mathrm{\Phi }|\widehat{B}|\mathrm{\Phi },`$ (8)
where $`\widehat{B}`$ is the so-called Bell operator. Starting with the CGLMP inequality and choosing suitable the experimental settings Durt01 Acin
$`\varphi _1(j)=0,\varphi _2(j)={\displaystyle \frac{j\pi }{d}},\phi _1(j)={\displaystyle \frac{j\pi }{2d}},\phi _2(j)={\displaystyle \frac{j\pi }{2d}},`$ (9)
optimal for maximal violations, we can derive the element of the Bell operator matrix as
$`B_{mm^{},jj^{}}=`$ $`{\displaystyle \frac{1}{d^2}}{\displaystyle \underset{l=0}{\overset{d1}{}}}e^{i\frac{2\pi }{d}[(jm)(j^{}m^{})]l}`$ (10)
$`\{e^{\frac{i\pi }{2d}(j^{}m^{})}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}(1{\displaystyle \frac{2k}{d1}})(e^{i\frac{2\pi }{d}k(jm)}e^{i\frac{2\pi }{d}(k+1)(j^{}m^{})})+`$
$`e^{\frac{i\pi }{2d}(j^{}m^{})}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}(1{\displaystyle \frac{2k}{d1}})(e^{i\frac{2\pi }{d}k(j^{}m^{})}e^{i\frac{2\pi }{d}(k+1)(jm)})+`$
$`e^{\frac{i\pi }{d}(jm)+\frac{i\pi }{2d}(j^{}m^{})}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}(1{\displaystyle \frac{2k}{d1}})(e^{i\frac{2\pi }{d}(k+1)(j^{}m^{})}e^{i\frac{2\pi }{d}k(jm)})+`$
$`e^{\frac{i\pi }{d}(jm)\frac{i\pi }{2d}(j^{}m^{})}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}(1{\displaystyle \frac{2k}{d1}})(e^{i\frac{2\pi }{d}k(jm)}e^{i\frac{2\pi }{d}(k+1)(j^{}m^{})})\},`$
where $`\widehat{B}`$, in general, is a $`d^2\times d^2`$ matrix. Ref. Acin found that the maximal eigenvalue of matrix $`\widehat{B}`$ is nothing but the highest quantum prediction of the Bell expression $`I_d(|\mathrm{\Phi }_{\mathrm{eig}})`$ and the corresponding eigenvector $`|\mathrm{\Phi }_{\mathrm{eig}}`$ is the state that maximally violates the Bell inequality. Thus, the problem of computing the maximal violation of $`I_d(|\mathrm{\Phi })`$ reduces to the determination of the maximal eigenvalue of matrix $`\widehat{B}`$. Due to $`_{l=0}^{d1}e^{i\frac{2\pi }{d}(pq)l}=d\delta _{pq}`$, where $`\delta _{pq}=1`$ when $`p=q`$ modulo $`d`$ and $`0`$ otherwise, the matrix $`\widehat{B}`$ can be further simplified, i.e., it can be decomposed into the sum of $`d`$ decoupled operators that act individually within the subspaces spanned by the vectors $`\{|00,|11,\mathrm{},|(d1)(d1)\}`$, $`\{|01,|12,\mathrm{},|(d1)0\}`$,$`\mathrm{}`$,$`\{|0(d1),|10,\mathrm{},|(d1)(d2)\}`$, respectively. For the first subspace spanned by the vectors $`\{|00,|11,\mathrm{},|(d1)(d1)\}`$, and with the constraint $`jm=j^{}m^{}`$, one can reduce $`B_{mm^{},jj^{}}`$ to
$`B_{mj}^{red1}`$ (11)
$`=`$ $`{\displaystyle \frac{8}{d}}\mathrm{sin}[{\displaystyle \frac{\pi }{2d}}(jm)]`$
$`\times {\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}(1{\displaystyle \frac{2k}{d1}})\mathrm{sin}[{\displaystyle \frac{2\pi }{d}}(k+{\displaystyle \frac{1}{2}})(jm)]`$
$`=`$ $`{\displaystyle \frac{2}{d1}}{\displaystyle \frac{1}{\mathrm{cos}[\frac{\pi }{2d}(jm)]}}(1\delta _{mj}),`$
where $`B_{mj}^{red1}`$ is the matrix element located at the $`m`$-th row and the $`j`$-th column of the reduced Bell operator $`B^{red1}`$. Let us denote $`B_r=B_{m,m+r}`$, since $`B_{m,j}`$ depends only on the difference $`(jm)`$. Therefore, we arrive at the following $`d\times d`$ real symmetric matrix for the reduced Bell operator:
$`\widehat{B}^{red1}`$ $`=`$ $`\left(\begin{array}{ccccccc}B_0& B_1& B_2& \mathrm{}& B_{d3}& B_{d2}& B_{d1}\\ B_1& B_0& B_1& B_2& \mathrm{}& B_{d3}& B_{d2}\\ B_2& B_1& B_0& \mathrm{}& \mathrm{}& \mathrm{}& B_{d3}\\ \mathrm{}& B_2& \mathrm{}& \mathrm{}& \mathrm{}& B_2& \mathrm{}\\ B_{d3}& \mathrm{}& \mathrm{}& \mathrm{}& B_0& B_1& B_2\\ B_{d2}& B_{d3}& \mathrm{}& B_2& B_1& B_0& B_1\\ B_{d1}& B_{d2}& B_{d3}& \mathrm{}& B_2& B_1& B_0\end{array}\right).`$ (19)
One may read from Eq.(11) that $`B_r`$ are all non-negative numbers, $`B_00`$, and interestingly $`lim_d\mathrm{}B_{d1}=4/\pi `$, $`lim_d\mathrm{}B_{dj}/B_{d1}=1/j`$, for $`j<<d`$. The Bell operator $`\widehat{B}^{red1}`$ acts on a matrix column $`(a_0,a_1,\mathrm{},a_{d1})^T`$, which corresponds to the state $`|\mathrm{\Psi }=_{j=0}^{d1}a_j|jj`$ in the first subspace. Actually, an arbitrary state $`|\mathrm{\Phi }=_{j,j^{}=0}^{d1}\alpha _{jj^{}}|jj^{}`$ can always be transformed into its Schmidt decomposition form $`|\mathrm{\Psi }=_{j=0}^{d1}a_j|jj`$ through local unitary transformations, thus it is sufficient to study the maximal violation problem in the first subspace. For a general state $`|\mathrm{\Psi }`$, one can obtain the Bell expression as
$`I_d(|\mathrm{\Psi })`$ $`=`$ $`\mathrm{Tr}(\widehat{B}^{red1}|\mathrm{\Psi }\mathrm{\Psi }|)={\displaystyle \underset{m=0}{\overset{d1}{}}}{\displaystyle \underset{j=0}{\overset{d1}{}}}B_{mj}a_ma_j`$ (21)
$`=`$ $`{\displaystyle \underset{r=1}{\overset{d1}{}}}B_r(2{\displaystyle \underset{m=0}{\overset{d1r}{}}}a_ma_{m+r}).`$
For instance, for the maximally entangled state $`|\mathrm{\Psi }_{\mathrm{mes}}`$ with $`a_j=1/\sqrt{d}`$, the summation $`(_{m,j=0}^{d1}B_{mj})/d=(_{r=1}^{d1}2(dr)B_r)/d`$ recovers the results of Eq. (1).
As the analysis shown in Acin , the maximal violation of the CGLMP inequality with the experimental settings (9) corresponds to the maximal eigenvalue of $`\widehat{B}^{red1}`$, and indeed its corresponding eigenvector is a nonmaximally entangled state of two-qudit. For instance, for $`d=3`$, one has $`B_0=0`$,$`B_1=2\sqrt{3}/3`$,$`B_2=2`$, the eigenvector $`|\mathrm{\Psi }_{\mathrm{eig}}^{d=3}=\sqrt{\frac{2}{11\sqrt{33}}}(|00+\frac{\sqrt{11}\sqrt{3}}{2}|11+|22)`$ corresponds to a maximal violation $`I_{d=3}(|\mathrm{\Psi }_{\mathrm{eig}})=1+\sqrt{11/3}2.9149`$, which is larger than $`I_{d=3}(|\mathrm{\Psi }_{\mathrm{mes}})2.8729`$; for $`d=4`$, one has $`B_0=0`$,$`B_1=2\sqrt{42\sqrt{2}}/3`$,$`B_2=2\sqrt{2}/3`$, $`B_1=2\sqrt{4+2\sqrt{2}}/3`$,the eigenvector $`|\mathrm{\Psi }_{\mathrm{eig}}^{d=4}=\frac{1}{\sqrt{2+2a^2}}(|00+a|11+a|22+|33)`$, with $`a=(\sqrt{2+\sqrt{2}}+\sqrt{83\sqrt{2}+4\sqrt{2\sqrt{2}}}\sqrt{4+2\sqrt{2}})/(\sqrt{2}+\sqrt{42\sqrt{2}})0.73937`$, corresponds to a maximal violation $`I_{d=4}(|\mathrm{\Psi }_{\mathrm{eig}})=\frac{2}{3}\sqrt{2+\sqrt{2}}+\frac{2}{3}\sqrt{83\sqrt{2}+4\sqrt{2\sqrt{2}}}2.9727`$, which is larger than $`I_{d=4}(|\mathrm{\Psi }_{\mathrm{mes}})2.8962`$.
By diagonalizing exactly the matrix $`\widehat{B}^{red1}`$, we have obtained the real maximal violations $`I_d(|\mathrm{\Psi }_{\mathrm{eig}})`$ for two entangled qudits. The highest dimension that we have calculated is $`d=8000`$. In Fig. 1, one may observe that $`I_d(|\mathrm{\Psi }_{\mathrm{eig}})`$ increases with dimension $`d`$. Based on the data of $`I_d(|\mathrm{\Psi }_{\mathrm{eig}})`$ from $`d=2`$ to $`d=8000`$, one has an empirical formula fitting $`I_d(|\mathrm{\Psi }_{\mathrm{eig}})`$ numerically to the dimension $`d`$:
$`I_d^{\mathrm{rough}}(|\mathrm{\Psi }_{\mathrm{eig}})`$ $``$ $`3.91321.2891x^{0.2226}`$ (22)
from which one can see that $`I_d^{\mathrm{rough}}(|\mathrm{\Psi }_{\mathrm{eig}})3.9132`$ is a coarse-grained limit of the maximal violation for the CGLMP inequality when $`d`$ tends to infinity.
Analysis of the eigenvectors $`|\mathrm{\Psi }_{\mathrm{eig}}`$ shows that these eigenvectors numerically satisfy some general properties: for instance, $`|\mathrm{\Psi }_{\mathrm{eig}}=_{j=0}^{d1}a_j^{\mathrm{eig}}|jj`$ with maximal eigenvalue has the following symmetric properties for the coefficients: $`a_j=a_{d1j}`$; and $`a_0:a_1:a_2:a_3:\mathrm{}1:\frac{1}{\sqrt{2}}:\frac{1}{\sqrt{3}}:\frac{1}{\sqrt{4}}:\mathrm{}`$ for large $`d`$. Thus, we may approximate a family of elegant entangled states
$`|\mathrm{\Psi }_{\mathrm{app}}`$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{d1}{}}}a_j^{\mathrm{app}}|jj,a_j^{\mathrm{app}}={\displaystyle \frac{1}{\sqrt{𝒩}}}{\displaystyle \frac{1}{\sqrt{(j+1)(dj)}}},`$ (23)
$`𝒩={\displaystyle \underset{j=0}{\overset{d1}{}}}{\displaystyle \frac{1}{(j+1)(dj)}}`$
whose corresponding Bell expressions $`I_d(|\mathrm{\Psi }_{\mathrm{app}})`$ are closed to the actual ones $`I_d(|\mathrm{\Psi }_{\mathrm{eig}})`$. For example, for $`d=8000`$, the error rate between $`I_d(|\mathrm{\Psi }_{\mathrm{eig}})`$ and $`I_d(|\mathrm{\Psi }_{\mathrm{app}})`$ is only about $`0.745\%`$. We have also listed $`I_d(|\mathrm{\Psi }_{\mathrm{app}})`$ and $`I_d(|\mathrm{\Psi }_{\mathrm{mes}})`$ in Table I and drawn the corresponding curves in Fig. 1. One may observe that for $`d=50000`$, $`I_d(|\mathrm{\Psi }_{\mathrm{mes}})2.96981`$ has almost reached the limit as shown in Eq. (2); and for $`d=600000`$, $`I_d(|\mathrm{\Psi }_{\mathrm{app}})`$ has exceeded 3.80.
In summary, we have investigated the maximal violation of the CGLMP inequalities for two entangled qudits. The maximal violation occurs at the non-maximally entangled state, which is the eigenvector of Bell operator with the maximal eigenvalue. The maximal violations increase with growing dimension $`d`$, coming close to 4 when $`d`$ approaches to infinity. Experimental test has been performed to verify the CGLMP inequalities for the first few high-dimensional systems, and indeed there exist non-maximally entangled states that violate these inequalities more strongly than the maximal entangled ones Vaziri . Bell inequalities are applicable to quantum cryptography and quantum communication complexity Kwek , previous researches are mostly based on maximally entangled states. For stronger violation, it may be useful to generate the approximate state in this paper. Nevertheless, it may be significant and interesting to apply the elegant symmetric entangled states $`|\mathrm{\Psi }_{\mathrm{app}}`$ to quantum cryptography as well as other important fields of quantum information.
This work is supported by NUS academic research Grant No. WBS: R-144-000-123-112. |
warning/0507/astro-ph0507283.html | ar5iv | text | # Discovery of X-ray eclipses from the transient source CXOGC J174540.0-290031 with XMM-Newton
## 1 Introduction
The Galactic center region is a very complex region of the Galaxy as viewed in X-ray wavelengths. This region contains numerous diffuse and point-like sources, fluorescent X-ray emission from molecular clouds, supernova remnants, compact stellar sources as well as young and evolved stellar clusters. In addition to Sgr A\*, the compact radio source which is thought to coincide with a supermassive black hole at the dynamical center of the Galaxy (Ghez et al. (2003), Schödel et al. (2003)), this region is populated by accreting compact objects, such as neutron stars and black hole binaries which are often transient in nature with peak X-ray luminosities in excess of 10<sup>36</sup> erg s<sup>-1</sup> (e.g., Churazov et al. (1997), Sidoli et al. (1999), Sakano et al. (2002), Porquet et al. (2003a)). Repeated, high sensitivity X-ray observations of this region carried out by XMM-Newton and Chandra, primarily for the purpose of monitoring Sgr A\*, have also lead to the discovery of several X-ray transient binaries in various accretion states from quiescence (L$`{}_{\mathrm{X}}{}^{}<`$ 10<sup>33</sup> erg s<sup>-1</sup>) to bright outburst (L$`{}_{\mathrm{X}}{}^{}>`$ 10<sup>36</sup> erg s<sup>-1</sup>) as reported by Sakano et al. (2005), Porquet et al. (2005), and Muno et al. (2005a). Transients with intermediate X-ray luminosities ($``$10<sup>34</sup>erg s<sup>-1</sup>) appear to be over-abundant by a factor of about 20 per unit stellar mass within 1 pc of Sgr A\* (Muno et al. (2005a)).
During the large XMM-Newton multi-wavelength project to monitor the light curve of Sgr A\*, observations were carried out in the Spring and Summer of 2004, a brightening in the 2–10 keV energy band of a factor of about two was detected within a radius of 10<sup>′′</sup> around SgrA\* (Bélanger et al. (2005)). The centroid of the excess X-ray emission was located at 2.9<sup>′′</sup> South of the radio position of Sgr A\*, and was coincident with CXOGC J174540.0-290031, a moderately bright transient, discovered by Chandra in July 2004 (Muno et al. (2005a)). VLA observations also detected bright radio transients symmetrically to CXOGC J174540.0-290031 whose flux peaked during the Spring observing campaign (Bower et al. 2005). The highest luminosity of CXOGC J174540.0-290031, in the 2–8 keV energy range, observed by Chandra was about 5$`\times `$10<sup>34</sup> erg s<sup>-1</sup> (Muno et al. (2005b)). This luminosity is too high to be explained by the X-ray emission from either the corona of an active star or the wind of a massive star, but conversely strongly favors the presence of an accreting compact object which could be either a white dwarf, or a neutron star, or a black hole. Since the typical outburst luminosity of white dwarf system is only observed up to about 10<sup>34</sup> erg s<sup>-1</sup> (e.g., GK Per, Sen & Osborne SO98 (1998)), we think that CXOGC J174540.0-290031 is more likely a X-ray binary with a neutron star or a black hole as compact object.
We report here X-ray eclipses from CXOGC J174540.0-290031 discovered during the XMM-Newton observations of Summer 2004. The observations of eclipses from X-ray binaries are rather rare with only 11 LMXB showing eclipses compared to the 80 LMXB reported in the catalog of Ritter & Kolb<sup>1</sup><sup>1</sup>1Version 7.4, http://www.mpa-garching.mpg.de/RKcat/. Indeed for very high inclination angles, the thickness of the accretion disc itself prevents us from seeing any of the X-rays at all. Eclipses are therefore observable only in a small number of systems where the inclination still allows us to see above the edge of the disk, see e.g. the review of Charles & Seward (1995). The period and duration of the eclipses can be used to constrain the mass of the secondary star and the mass of the compact object, respectively.
This paper is structured in the following manner. In section 2 we present the XMM-Newton observations and data reduction methods. Sections 3 and 4 describe the timing and spectral analysis of CXOGC J174540.0-290031, and are followed in sect. 5 and 6 by the study of the characteristics of the transient eclipses and the constraints that these impose upon the nature of the source.
## 2 XMM-Newton observations and data reduction
The XMM-Newton observations presented here are part of a large multi-wavelength campaign to monitor the flux of Sgr A\*. The observations were made from March 28 to April 1 2004 (revolutions 788 and 789), and from August 31 to September 2 2004 (revolutions 866 and 867). The total exposure time is about 500 ks. Table 1 gives the observation log of the PN camera with the beginning and the end time of the first and last Good Time Interval (GTI) of the on-axis CCD #4, respectively. We also give in this table the XMM-Newton observation of Sgr A\* of February 26, 2002, where the flaring background was very low, where no X-ray flares were observed from Sgr A\*, and where CXOGC J174540.0-290031 was in its quiescent state. This latter observation will be used as reference for the X-ray flux of the quiescent emission of Sgr A\* and the diffuse emission in the extraction region of CXOGC J174540.0-290031. The PN camera was in Prime Full extended window and Prime Full window mode, during revolutions 788 and 789, and revolutions 406, 866, and 867, respectively. The MOS cameras were in Prime Full window mode during all the observations. The data were processed with SAS (version 6.1.0). We selected only X-ray events with data quality flag equal to 0. X-ray events with patterns 0–12 and 0–4 were used for MOS and PN, respectively.
## 3 Timing analysis
The position of CXOGC J174540.0-290031 was determined in each XMM-Newton observation using its relative angular separation from three bright X-ray point sources, namely CXO J174530.0-290704, CXO J174607.5-285951, and CXO J174530.0-290704 (Muno et al. (2004)), which were observed by Chandra and XMM-Newton. We used the 5.1 ks Director Discretionary Time observation of CXOGC J174540.0-290031 made by Chandra on 28 August 2004 (ObsID 5360, PI: F. Baganoff) to measure these three angular separations. The angular separations between CXOGC J174540.0-290031 and the reference X-ray sources CXO J174530.0-290704, CXO J174607.5-285951, and CXO J174530.0-290704 are 79.5$`\pm `$0.1 ″, 362.2$`\pm `$0.4 ″, and 415.3$`\pm `$0.3 ″, respectively. The positions of these three sources were obtained in each XMM-Newton observation with the SAS command edetect\_chain, using simultaneous detections in MOS1, MOS2, and PN images in the 0.3–10 keV energy band. We extract for each of the five XMM-Newton observations, the source events from a circular region of 10<sup>′′</sup>-radius centered on the inferred XMM-Newton position of the transient source.
We build the light curve of each camera from the beginning to the end of the first and the last GTI of the (on-axis) CCD #4 of PN, respectively. We corrected the light curves for the loss of exposure due to the triggering of counting mode during high flaring background, where count rate goes beyond the detector telemetry limit. Using the GTI extension of the event file, we computed for each time bin the ratio of the bin length to exposure loss, and multiplied the count rate and corresponding error by this linear correction factor. Finally, light curves of the three detectors were added to produce the EPIC light curves for the four observations in the 2–10 keV energy range. The contribution of the quiescent X-ray emission of Sgr A\* and the diffuse X-ray emission have been removed by subtracting the averaged count rate observed in the same extraction region (0.161$`\pm `$0.002 cts s<sup>-1</sup>) during the February 2002 observation (see Table 1). The X-ray flaring background has a negligible contribution to the total flux in the 10<sup>′′</sup>-radius region around SgrA\* assuming the incident particles are distributed uniformly across the detector.
Figure 1 shows the corresponding EPIC light curve of CXOGC J174540.0-290031 in the 2–10 keV energy band for each of the four XMM-Newton revolutions. X-ray flares identified as coming from SgrA\* are clearly visible and the analysis of these events is presented in Bélanger et al. (2005). During revolution 866 we observed five periodic sharp flux troughs (indicated by thick vertical arrows in Fig. 1). These features will be discussed in sect. 5.
## 4 Spectral analysis
The spectral analysis has used only PN data, which have a higher S/N than the MOS data, and are sensitive up to 12 keV. The PN data has been cleaned from flaring background contributions by removing all time intervals where the PN background count rate was higher than 50 cts s<sup>-1</sup> and 15 cts s<sup>-1</sup> for the Spring and Summer 2004 observations, respectively. In addition we removed all time intervals with a significant contribution of X-ray flares from Sgr A\* (marked with horizontal arrows in Fig. 1). The astrophysical X-ray background (quiescent emission of Sgr A\* and diffuse emission) was obtained from the XMM-Newton observation of February 26, 2002, where both Sgr A\* and CXOGC J174540.0-290031 were in their quiescent state. With this double subtraction method, only the contribution of CXOGC J174540.0-290031 above its quiescent state is detected.
We use the updated X-ray absorption cross-sections of Wilms et al. (2000) (tbabs in xspec). We fit the data taking into account the scattering of X-rays by dust, using the scatter model (Predehl & Schmitt (1995)) assuming a visual extinction value $`A_\mathrm{V}`$=30 mag, as determined from IR observations of stars close to Sgr A\* (e.g., Rieke et al. (1989)). The fitting parameter errors correspond to 90% confidence ranges for one interesting parameter ($`\mathrm{\Delta }\chi ^2`$=2.71). Table 2 gives our best fit parameters for the four observations using an absorbed power law continuum model. Taking into account the error bars, the photon power law indices are consistent with a constant value in the four observations. However there is a hint of spectral hardening and decrease of the intrinsic luminosity during the last observation. Interestingly, there is a significant variation of the column density between Spring and Summer 2004 (see variation of the soft emission in Figure 2), with the highest value of $`N_\mathrm{H}`$ observed during the revolution 866 of Summer 2004, where several deep X-ray eclipses from CXOGC J174540.0-290031 are observed.
## 5 X-ray eclipses
As reported in sect. 3, five periodic sharp flux troughs are observed during revolution 866. Since the duration of these troughs is short compared to the trough period, we associate these features with eclipses. These eclipses are likely produced by the secondary star of the binary system rather than by optical depth effects or electron scattering in material near the magnetic pole of the compact object (as seen in some NS or WD systems), which would produced broader dips. Therefore we will assume in the following that this eclipse periodicity is associate with the orbital motion of the binary system, where a compact object emits X-rays.
We determine the eclipse ephemeris using the five eclipses of revolution 866, where the eclipses are well defined, and not affected by exposure losses (see vertical thick arrows in Fig. 1). We start from the arrival times of photons detected by MOS1+MOS2+PN, which are convolved with a Gaussian smoothing kernel to reduce the Poissonian noise and to produce a smooth shape of the light curve. Each eclipse is then fitted with a Gaussian plus a linear term to measure accurately the five eclipse epochs. The eclipse period ($`P_{\mathrm{eclipse}}`$) and the epoch of the first eclipse (T<sub>0</sub>) are obtained with corresponding errors from mean and rms of the eclipse epoch differences. The best estimate of the eclipse period is obtained by minimizing its rms error by varying the FWHM of the smoothing kernel. The minimum rms is obtained for a FWHM$``$760 s, leading to T<sub>0</sub>=05h00m18s and $`P_{\mathrm{eclipse}}=27,961\pm 5`$ s. The Julian date epochs of the eclipses are $`JD=(2453248.85856\pm 0.00004)+0.323626\pm 0.00006\times n`$, where $`n`$ is an integer. The predicted eclipses are marked in Fig. 1 by thin vertical arrows. While similar fainter eclipses, coinciding with the ephemeris predictions (see vertical thin arrows), are visible during revolution 867, no obvious deep eclipses are observed during the revolutions 788 and 789 (Spring 2004 observations) at the times predicted by the eclipse ephemeris. Based on the Chandra observations performed in July and August 2004, Muno et al. (2005a) proposed the possible presence of dips from this source with a tentative period of 7.9 h, slightly larger than our measurement. Our eclipse ephemeris predict well the epochs of the troughs observed with Chandra in July 2004.
To obtain folded light curves, the arrival times of photons detected by each instrument when Sgr A\* was not flaring were folded using the eclipse period. Then using the GTI extension and the time intervals of Sgr A\* flares, we determine the net exposure within each phase bin to calculate the count rate. Instrument folded light curves are then added to obtain the EPIC folded light curve. Figure 3 shows the resulting folded light curves. We compute the average count rates of the transient and its dispersion outside the time interval where the eclipse occurs. The average count rate level for the revolutions 788, 789, 866, and 867 are $`<`$CR$`>`$=0.16$`\pm `$0.03, 0.16$`\pm `$0.03, 0.14$`\pm `$0.03, and 0.11$`\pm `$0.02 cts s<sup>-1</sup>, respectively. During each epoch, $`<`$CR$`>`$ is constant, whereas there is a hint of a decrease between the Spring and Summer epochs. Deep eclipses are visible for the revolutions 866 and 867 during Summer 2004 (Fig. 3 lower panels). The duration of the eclipse is defined as the time interval where the count rate is lower than one sigma from the average count rate. Both eclipse duration measurements are consistent with about 1,100$`\pm `$100 s, which is a small fraction ($``$4%) of the eclipse period. The eclipses are not total as suggested by the folded curves of the revolutions 866 and 867 (Fig. 3) where the minimum flux is not compatible with zero. The eclipse of revolution 867 is less deep ($`\mathrm{\Delta }CR/<CR>60\%`$) than the eclipse of revolution 866 ($`80\%`$). This could suggest an extended X-ray emitting region, which would be only partly hidden by the secondary star. As previously suggested by the light curves of revolutions 866 and 867 there is no such deep eclipses during the Spring 2004 observations (see the upper panels of Fig. 3 upper panels). We can exclude confidently for this epoch any eclipse with depth greater than 19% (i.e. below one sigma dispersion of the light curve). The absence of these deep X-ray eclipses in Spring 2004 are possibly due to a displacement of the X-ray emitting region above the orbital plane as we will discuss in §6.3.
## 6 Nature of CXOGC J174540.0-290031
We consider that the X-ray source CXOGC J174540.0-290031 is a semidetached binary system, with a circular orbit (tidal circularization in binary system with short period), which has a sufficiently high inclination to the line-of-sight so as to produce X-ray eclipses of the compact object by the secondary star. In this framework, the eclipse period ($`P_{\mathrm{eclipse}}=27,961`$ s) is identified as the orbital period of the binary system, and will provide constraints on the mass of the secondary star, whereas the eclipse fraction ($`\mathrm{\Delta }_{\mathrm{eclipse}}=0.039`$) will then constrain the mass of the compact object (e.g., Weisskopf et al. (2004)).
### 6.1 Upper limit on the mass of the secondary star from the eclipse period
We can determine a mass-radius relation of the secondary star knowing the orbital period. The mass of the secondary star is
$$\frac{M_2}{\mathrm{M}_{}}=\frac{\overline{\rho }_2}{\overline{\rho }_{}}\left(\frac{R_2}{\mathrm{R}_{}}\right)^3,$$
(1)
where $`\overline{\rho }_2`$, $`R_2`$, and $`\overline{\rho }_{}`$ are the mean density and the radius of the secondary star, and the solar mean density (1.41 g cm<sup>-3</sup>), respectively. We consider that the secondary star’s radius is equal to the secondary star’s effective Roche lobe radius<sup>2</sup><sup>2</sup>2The effective Roche lobe radius is defined so that the sphere with this radius has the same volume that the one within the Roche lobe., $`r_\mathrm{L}`$, given by Eggleton (1983)’s approximation:
$$\frac{r_\mathrm{L}}{a}=\frac{0.49q^{2/3}}{0.6q^{2/3}+\mathrm{ln}(1+q^{1/3})},$$
(2)
where $`0<qM_2/M_1<\mathrm{}`$ (where $`M_1`$ is the mass of the primary star), and $`a`$ is the separation of the two stars, which is related to the eclipse period by
$$\frac{a}{\mathrm{R}_{}}=2.35\left(\frac{P_{\mathrm{eclipse}}}{10\mathrm{h}}\right)^{2/3}\left(\frac{M_1+M_2}{\mathrm{M}_{}}\right)^{1/3}.$$
(3)
Combining Eqs. (1)–(3) leads to:
$$\frac{\overline{\rho }_2}{\overline{\rho }_{}}=0.66\left(\frac{P_{\mathrm{eclipse}}}{10\mathrm{h}}\right)^2\frac{\{0.6q^{2/3}+\mathrm{ln}(1+q^{1/3})\}^3}{q(1+q)}.$$
(4)
Applying Eq. (1) in Eq. (4) gives the mass-radius relation of the secondary star knowing the orbital period:
$`{\displaystyle \frac{R_2}{\mathrm{R}_{}}}`$ $`=`$ $`1.15\left({\displaystyle \frac{P_{\mathrm{eclipse}}}{10\mathrm{h}}}\right)^{2/3}`$ (5)
$`\times {\displaystyle \frac{q^{1/3}(1+q)^{1/3}}{0.6q^{2/3}+\mathrm{ln}(1+q^{1/3})}}\left({\displaystyle \frac{M_2}{\mathrm{M}_{}}}\right)^{1/3}.`$
Figure 4 shows this mass-radius relation for the eclipse period, compared with the mass-radius relation of low-mass stars in the Zero-Age Main-Sequence stage (ZAMS), calculated by Siess et al. (2000). When the star evolves away from the main-sequence its stellar radius increases. Assuming that the secondary star is on the ZAMS or older, we can see that for a large range of masses of the compact star (taking $`M_1=1.4`$–3.0 M for a neutron star, or $`M_1=5`$$`30`$ M for a stellar-mass black hole), the mass donor star has mass $`\begin{array}{c}<\hfill \\ \hfill \end{array}\mathrm{\hspace{0.17em}1.05}`$ M. As shown by Kolb et al. (2001) using stellar models and allowing for different mass transfer rates and different system ages prior to mass transfer, the mass of a ZAMS star with the same spectral type is an upper limit to the donor mass. As well, Howell et al. (2001) shown that the secondary masses in CVs just above the period gap of about 2–3 h should be as much as $``$50% lower than would be inferred if one assumes ZAMS donor star. Therefore, the value of the mass donor star found here ($`\begin{array}{c}<\hfill \\ \hfill \end{array}\mathrm{\hspace{0.17em}1.05}`$ M) assuming a ZAMS star is a conservative upper limit, hence, CXOGC J174540.0-290031 is a Low-Mass X-ray Binary (LMXB).
This maximum stellar mass corresponds to a mass donor star with spectral type later than about K0. For comparison a ZAMS star later than K0 (which corresponds to an age of about 150 million years) has an absolute magnitude in the infrared $`K`$ band of $`M_\mathrm{K}>3.50`$ mag (Siess et al. (2000)), which corresponds at the galactic center distance ($`d=7.94`$ kpc; Eisenhauer et al. (2003)) where the extinction in the $`K`$ band is $`A_\mathrm{K}=2.7`$ mag (Clénet et al. (2001)), to $`K>20.7`$ mag. The detection of the mass donor late-type star with ground adaptive optics (e.g., NACO at VLT using for the infrared wavefront sensor the close $`K6.5`$ mag source IRS 7 as reference star) will allow to identify the X-ray transient binary in the near-infrared.
### 6.2 Upper limit on the mass of the compact object from the eclipse fraction
The eclipse fraction, i.e. the ratio of the eclipse duration to the orbital period, is geometrically related to the angular size of the path of the observer behind the secondary photosphere seen from the compact object by
$$\mathrm{\Delta }_{\mathrm{eclipse}}=\frac{1}{\pi }\mathrm{arccos}\left(\frac{1}{\mathrm{sin}i}\sqrt{1\left(\frac{R_2}{a}\right)^2}\right),$$
(6)
where $`i`$ is the orbital inclination. Using Eq. (2), we find that the eclipse fraction depends only on the orbital inclination and the mass ratio.
Figure 5 shows the eclipse fraction as a function of $`q^1`$ ($`M_1/M_2`$), for different orbital inclinations. To each mass ratio corresponds a critical inclination where only a grazing eclipse is possible, below this critical inclination an eclipse of the compact source by the secondary star cannot occur. More generally, the proper inclination needed to obtain a given eclipse fraction can be derived from Eq. (5):
$$i=\mathrm{arcsin}\left(\frac{\sqrt{1\left(\frac{0.49q^{2/3}}{0.6q^{2/3}+\mathrm{ln}(1+q^{1/3})}\right)^2}}{\mathrm{cos}(\pi \mathrm{\Delta }_{\mathrm{eclipse}})}\right).$$
(7)
This equation is useful to compute the eclipse visibility as a function of the mass ratio and the orbital inclination (see Figure 7).
The maximum eclipse duration is obtained in Fig. 5 for $`i=90^{}`$, i.e. when the line-of-sight is parallel to the orbital plane. Consequently, the observed eclipse fraction provides an upper limit on $`q^1`$. If the inclination is known, $`q^1`$ is obtained by solving numerically the following equation:
$$\frac{\mathrm{ln}(1+q^{1/3})}{q^{2/3}}=\frac{0.49}{\sqrt{1\mathrm{sin}^2i\mathrm{cos}^2(\pi \mathrm{\Delta }_{\mathrm{eclipse}})}}0.6.$$
(8)
The eclipse duration of 1,100 s corresponds to $`\mathrm{\Delta }_{\mathrm{eclipse}}`$=0.039, implying that the maximum value for $`q^1`$, i.e. $`M_1/M2`$, is $`56.5`$. Since $`M_2\begin{array}{c}<\hfill \\ \hfill \end{array}\mathrm{\hspace{0.17em}1.05}`$ M (see §6.1), we obtain that the mass of the compact object is $`\begin{array}{c}<\hfill \\ \hfill \end{array}\mathrm{\hspace{0.17em}59}`$ M.
Accretion disks around compact objects are likely not thin, and should then produce a permanent total X-ray eclipse when seen exactly edge-on. Therefore an orbital inclination of $`90^{}`$ is unlikely. Figure 7 shows the inferred maximum mass of the compact object as a function of the orbital inclination. Inclinations lower than $`70^{}`$ can be excluded because the luminosity of the X-ray transient exclude a white dwarf as primary star. We conclude that the possible values of the orbital inclination are between $``$$`70^{}`$ and $`90^{}`$. If the orbital inclination is $``$$`70^{}`$$``$$`74^{}`$ a black hole can be excluded.
### 6.3 Constraints from the missing eclipses in Spring 2004
The absence of eclipse in Spring 2004 cannot be explained easily in this strict geometry of the binary system, where a decrease of the secondary star radius appears rather unlikely. We investigate a vertical motion of the X-ray emitting region from the orbital plane to a region located above the compact object. In such a configuration, the critical inclination, $`i_{\mathrm{cr}}`$, which produces grazing eclipses, is related to the vertical elevation, $`h`$, of the X-ray emitting region above the orbital plane by
$$h=\frac{R_2a\sqrt{1\mathrm{sin}^2i_{\mathrm{cr}}}}{\mathrm{sin}i_{\mathrm{cr}}}.$$
(9)
Taking $`i_{\mathrm{cr}}`$ equal to the orbital inclination constrained from Eq. (7) provides the value of the minimum vertical elevation needed to suppress the eclipses. Replacing in Eq. (9), $`a`$ and $`\mathrm{sin}i`$ by Eq. (3) and (7), respectively, leads to
$`{\displaystyle \frac{h}{\mathrm{R}_{}}}`$ $`=`$ $`2.35\left({\displaystyle \frac{P_{\mathrm{eclipse}}}{10\mathrm{h}}}\right)^{2/3}\left({\displaystyle \frac{1+q}{q}}\right)^{1/3}`$ (10)
$`\times \left({\displaystyle \frac{\mathrm{cos}(\pi \mathrm{\Delta }_{\mathrm{eclipse}})}{\sqrt{(\frac{R_2}{a})^21}}}\sqrt{{\displaystyle \frac{(\frac{R_2}{a})^2\mathrm{sin}^2(\pi \mathrm{\Delta }_{\mathrm{eclipse}})}{1(\frac{R_2}{a})^2}}}\right)`$
$`\times \left({\displaystyle \frac{M_2}{\mathrm{M}_{}}}\right)^{1/3},`$
where $`R_2/a`$ is a function of the mass ratio only, see Eq. (2). Using $`M_2\begin{array}{c}<\hfill \\ \hfill \end{array}\mathrm{\hspace{0.17em}1.05}`$ M provides an upper limit on the minimum value of $`h`$ needed to produce the disappearance of the eclipses.
Figure 8 shows $`h`$ as a function of the mass ratio for $`M_2=1.05`$ M. If the compact object is a neutron star a vertical elevation of the X-ray emitting region of only $`h0.1`$ R$`{}_{}{}^{}0.1`$ R<sub>2</sub> is sufficient to suppress the eclipses by the secondary star. If the compact object is a black hole a greater vertical elevation is needed.
## 7 Discussion
From the Chandra observations (Muno et al. (2005a)), the minimum 2–8 keV luminosity of CXOGC J174540.0-290031 at 8 kpc was determined to be less than 2 $`\times `$ 10<sup>31</sup> erg s<sup>-1</sup> (Muno et al. (2005a)), and the maximum to be about 5 $`\times `$ 10<sup>34</sup> erg s<sup>-1</sup> (Muno et al. (2005b)). This large amplitude variation corresponds to a factor of at least 2,500. We infer a minimum 0.5–10 keV luminosity of less than 3.5 $`\times `$ 10<sup>31</sup> erg s<sup>-1</sup>. Such a low luminosity lends support to the interpretation that the compact object of this binary system is more likely to be a black hole than a neutron star. Indeed, the quiescent 0.5–10 keV X-ray luminosities of neutron stars are about 10<sup>32</sup>–10<sup>34</sup> erg s<sup>-1</sup>, i.e., 100 times higher than those for black holes (10<sup>30</sup>–10<sup>33</sup> erg s<sup>-1</sup>), as shown by Garcia et al. (2001). In addition, VLA observations have revealed a radio outburst coincident with the X-ray position of this source, with a peak intensity of 100 mJy at 7 mm (45 GHz) (Bower et al. (2005)). The peak radio emission occurred in March 2004 and continued decaying in its surface brightness until the start of the September 2004 campaign (Bower et al. (2005)). Therefore, the minimum luminosity and the radio outburst strongly suggest that this source is more likely to contain a black hole than a neutron star primary (e.g., Fender & Kuulkers (2001)).
We observed with XMM-Newton X-ray eclipses from CXOGC J174540.0-290031 with a period of 27,961 $`\pm `$ 5 s (about 7.76 hours) in the period from August 31 to September 2 2004 but did not however detect such a feature in the earlier observations of March 28 to April 1 2004. In the framework of eclipsing semidetached binary systems, we show that the 27,961$`\pm `$5 s eclipse period constrains the mass of the assumed main-sequence secondary star to less than 1.0 M. Therefore, we deduce that CXOGC J174540.0-290031 is a low-mass X-ray binary. Moreover the eclipse duration (1,100$`\pm `$100 s) constrains the mass of the compact object to less than about 60 M, which is consistent with a stellar mass black hole or a neutron star.
The obtained XMM-Newton spectra are well described with an absorbed power law continuum. The column density values on the line-of-sight and the 2–10 keV luminosities confirm that this source is an X-ray binary, with a black hole or a neutron star as compact object, located at or near the Galactic Center. The 2–10 keV luminosity is almost constant with 1.8–2.3 $`\times `$ 10<sup>34</sup> erg s<sup>-1</sup> over the four observations, assuming a distance of 8 kpc. While the power law indices ($`\mathrm{\Gamma }`$=1.6–2.0) are consistent within the statistical uncertainty over the four observations, we observed a significant increase of the column density during the Summer 2004 observations, i.e. when X-ray eclipses are observed.
We propose to explain the absence and the presence of deep X-ray eclipses during the Spring 2004 and Summer 2004 observations, respectively, without any significant change of the X-ray luminosity (see Table 2) by a different position of the X-ray emitting region. In Spring 2004 observations the absence of X-ray eclipses could be due to the vertical movement of the centroid of the X-ray emission from a point on the orbital plane to a position above the compact object (the offset being at least 0.1 times the radius of the secondary star for a neutron star and more for a black hole primary). In such a configuration, our line-of-sight to the X-ray emitting region does not intercept the upper layer of the accretion disk, which is consistent with the lower column density value observed in the Spring 2004 observations. We propose that this displacement of the X-ray emitting region could be linked to the presence of a jet. When the jet is launched, a part of the hot plasma accreting towards the central compact object is lift off above the accretion disk, extending consequently the X-ray emission region well above the limb of the secondary star. On the contrary in Summer 2004, the X-ray emission from the base of the jet becomes almost negligible and most of the X-ray emission comes from the mid-plane accretion disk which is eclipsed by the secondary star. This scenario is compatible with the significant increase in the column density during the Summer 2004 observations, since here the line-of-sight to the X-ray emitting region intercepts the upper layer of the accretion disk.
Radio observations of this transient object strongly support the appearance of a jet. Indeed, two new transient radio sources were discovered by the VLA at 7 mm on March 2004 symmetrically to the position of the X-ray source (Bower et al. (2005), see also fig. 1 in Muno et al. (2005a)). The peak of the radio emission occurred in March 2004 and continued decaying in its surface brightness until the start of the September 2004 campaign (Bower et al. (2005)). Also VLA observations at 6 cm show clearly the brightening of this transient radio source in the 2004 epoch compared to the 1983 epoch (Yusef-Zadeh, Cotton & Bower 2005, in preparation). These two bright radio sources (as well as X-ray reflection off the orbiting gas reported by Muno et al. (2005b)) most likely associated with the binary system CXOGC J174540.0-290031 can be the imprint of radio emission from a collimated bipolar outflow, as observed in microquasars (e.g., Spencer et al. (1997)), where a mass ejection is triggered by an instability in the accretion disk of a binary system, with the primary star being a black hole as suggested above. In conclusion, CXOGC J174540.0-290031 appears to be a LMXB with more likely a black hole as compact object, which lead to an inclination angle of the system greater than 75 according to our study. Sensitive follow up observations of this embedded X-ray binary with ground based adaptive optic system should be able to detect the mass donor star in the near-IR and allow us to impose more restrictive constraints on the parameters of this transient eclipsing binary system.
###### Acknowledgements.
This work is based on observations obtained with XMM-Newton, an ESA science mission with instruments and contributions directly funded by ESA Member States and the USA (NASA). The authors would like to thank the referee for useful comments. D.P. is supported by a MPE fellowship. |
warning/0507/cond-mat0507518.html | ar5iv | text | # Kondo Resonance Decoherence by an External Potential
## 1 Introduction
The transport of electrons through a quantum dot (QD) is a subject of increasing interest for its many applications in nanotechnology. Although there exists exact solutions for the ground state of the QD connected to metallic leads , no such methods are applicable to nonequilibrium systems. Typical cases are those of (a) a bias or (b) an electromagnetic radiation between the leads of a quantum well . It has been suggested that, in case (a), two Kondo resonances appear, each one associated with every lead ,. Some experimental evidence has been presented supporting this view . Still, it is not well understood how inelastic (decoherent) processes modify the strong-coupling regime associated with the Kondo resonance. To our knowledge, there are only a few theoretical approaches to the problem. The non-crossing approximation (NCA) ,, being essentially numerical, can give analytical results only in the limit $`Ln(V/T_K)>>1`$ (V being the applied bias and $`T_K`$ the Kondo energy), the same limit in which renormalized perturbation theory (RPT) , is valid. It is important, however, to have analytical results for the decoherence rate, not only to make theoretical predictions, but also to be able to interpret correctly the results of calculations done for more complicated situations, i.e., when time- dependent potentials are applied to the QD -. A version of the nonequilibrium Kondo problem has been solved exactly to all orders in the applied bias V in .
In this article, we address the nonequilibrium Kondo problem for a QD in the limit of infinite Coulomb repulsion U, using an equation of motion (EOM) method based on a slave-boson Hamiltonian. By means of a strict perturbative solution in the leads-dot coupling, $`T`$, to orders $`T^4`$ and $`T^6`$, we have been able to identify the decoherent processes that broaden the Kondo-level, as a function of the applied bias, and break the strong-coupling regime. The EOM method has been applied before to the Kondo problem but it was mainly used with approximations valid only for high temperatures ,,. We show that this method, when applied consistently to high orders, can be conveniently used in the nonequilibrium case, describing well the inelastic processes contributing to the Kondo resonance decoherence. We first check the method in the equilibrium regime, comparing our results with exact results using the numerical renormalization group (NRG) and also with the NCA. When used in the nonequilibrium case, we find that, for an applied bias larger than $`2T_K`$, the system develops two-peaked resonant structures at $`\pm V/2`$; we show, however, that decoherent processes take the system from the strong-coupling regime, where the spin of the impurity is quenched by the lead electrons, to the weak-coupling regime, if $`V/T_K`$ is larger than a few units. We conclude that our method provides an accurate and powerful tool to analyze nonequilibrium systems and time-dependent problems, situations where very few theoretical approaches are valid.
## 2 Theory
Our starting point for a QD level with infinite-U, coupled to two leads, is the slave-boson Hamiltonian for spin up and down electrons :
$`\widehat{H}={\displaystyle \underset{k\sigma }{}}ϵ_{k\sigma }\widehat{c}_{k\sigma }^{}\widehat{c}_{k\sigma }+{\displaystyle \underset{\sigma }{}}ϵ_0\widehat{d}_\sigma ^{}\widehat{d}_\sigma +{\displaystyle \underset{k\sigma }{}}T_k(\widehat{c}_{k\sigma }^{}\widehat{\chi }_\sigma +\widehat{\chi }_\sigma ^{}\widehat{c}_{k\sigma })`$ (1)
where the operators $`\widehat{\chi }_\sigma `$ and $`\widehat{\chi }_\sigma ^{}`$ are introduced in terms of an auxiliary boson field, $`\widehat{b}`$ and $`\widehat{b}^{}`$, as $`\widehat{\chi }_\sigma =\widehat{b}^{}\widehat{d}_\sigma `$, $`\widehat{\chi }_\sigma ^{}=\widehat{d}_\sigma ^{}\widehat{b}`$ and $`k=k_L,k_R`$ for left and right leads. In eq.(1), the operator $`\widehat{c}_{k\sigma }^{}`$ ($`\widehat{c}_{k\sigma }`$) creates (anhilates) an electron of momentum $`k`$ and spin $`\sigma `$ in one of the leads, while $`\widehat{d}_\sigma ^{}`$ ($`\widehat{d}_\sigma `$) creates (anhilates) an electron of spin $`\sigma `$ in the dot. The constraint relation
$$\widehat{b}^{}\widehat{b}+\underset{\sigma }{}\widehat{d}_\sigma ^{}\widehat{d}_\sigma =1$$
(2)
is imposed to guarantee that no more than one electron occupies the QD since we are in the infinite-U limit.
The physical (retarded) Green function for this problem is defined as ,
$$<<\widehat{\chi }_\sigma (t);\widehat{\chi }_\sigma ^{}(t^{})>>=i<\widehat{\chi }_\sigma ^{}(t^{})\widehat{\chi }_\sigma (t)+\widehat{\chi }_\sigma (t)\widehat{\chi }_\sigma ^{}(t^{})>\theta (tt^{})$$
(3)
and its Fourier transform is written as $`<<\widehat{\chi }_\sigma ;\widehat{\chi }_\sigma ^{}>>(\omega )`$. The EOM for this Green function reads (omitting the argument $`\omega `$ in all the Green functions):
$`(wϵ_0)<<\widehat{\chi }_\sigma ;\widehat{\chi }_\sigma ^{}>>=1<\widehat{n}_\sigma >+{\displaystyle \underset{k}{}}T_k<<(1\widehat{n}_\sigma )\widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>+`$
$`{\displaystyle \underset{k}{}}T_k<<\widehat{d}_\sigma ^{}\widehat{d}_\sigma \widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>`$ (4)
where the new Green functions appearing on the right-hand side of eq.(4) are defined as in eq.(3). $`\widehat{n}_\sigma `$ is the number operator for electrons of spin $`\sigma `$ in the QD and $`<\widehat{A}>`$ will denote the expectation value of $`\widehat{A}`$ from now on. In a second step, the EOM for these Green functions lead to a second generation of Green functions, viz.,
$$(wϵ_k)<<\widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>=T_k<<\widehat{\chi }_\sigma ;\widehat{\chi }_\sigma ^{}>>$$
(5)
$`(wϵ_k)<<\widehat{d}_\sigma ^{}\widehat{d}_\sigma \widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>=<\widehat{\chi }_\sigma ^{}\widehat{c}_{k\sigma }>+`$
$`{\displaystyle \underset{k^{}}{}}T_k^{}<<\widehat{\chi }_\sigma \widehat{c}_{k^{}\sigma }^{}\widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>+{\displaystyle \underset{k^{}}{}}T_k^{}<<\widehat{\chi }_\sigma ^{}\widehat{c}_{k\sigma }\widehat{c}_{k^{}\sigma };\widehat{\chi }_\sigma ^{}>>`$ (6)
$`(wϵ_k)<<\widehat{n}_\sigma \widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>={\displaystyle \underset{k^{}}{}}T_k^{}<<\widehat{\chi }_\sigma \widehat{c}_{k^{}\sigma }^{}\widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>+`$
$`{\displaystyle \underset{k^{}}{}}T_k^{}<<\widehat{\chi }_\sigma ^{}\widehat{c}_{k^{}\sigma }\widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>`$ (7)
Note that the second generation of Green functions involve the excitation of two electrons in the leads. In our approximation to $`O(T^2)`$, we contract operators $`\widehat{c}_{k\sigma }`$, $`\widehat{c}_{k^{}\sigma }^{}`$, where possible and obtain a closed equation for $`<<\widehat{\chi }_\sigma ;\widehat{\chi }_\sigma ^{}>>`$ . In particular we take $`<<\widehat{\chi }_\sigma \widehat{c}_{k^{}\sigma }^{}\widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>=\delta _{k,k^{}}<\widehat{n}_{k\sigma }><<\widehat{\chi }_\sigma ;\widehat{\chi }_\sigma ^{}>>`$ and $`<<\widehat{\chi }_\sigma ^{}\widehat{c}_{k\sigma }\widehat{c}_{k^{}\sigma };\widehat{\chi }_\sigma ^{}>>=<<\widehat{\chi }_\sigma \widehat{c}_{k^{}\sigma }^{}\widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>=<<\widehat{\chi }_\sigma ^{}\widehat{c}_{k^{}\sigma }\widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>=0`$. This yields the following approximation to $`<<\widehat{\chi }_\sigma ;\widehat{\chi }_\sigma ^{}>>`$, which is exact to $`O(T^2)`$ (see Ref. ):
$$<<\widehat{\chi }_\sigma ;\widehat{\chi }_\sigma ^{}>>=\frac{1<\widehat{n}_\sigma >\underset{k}{}\frac{T_k<\widehat{\chi }_\sigma ^{}\widehat{c}_{k\sigma }>}{\omega ϵ_k}}{\omega ϵ_0\mathrm{\Sigma }_0_k\frac{T_k^2<\widehat{n}_k>}{\omega ϵ_k}}$$
(8)
where
$$\mathrm{\Sigma }_0=\underset{k}{}\frac{T_k^2}{\omega ϵ_k}$$
(9)
We can obtain $`<<\widehat{\chi }_\sigma ;\widehat{\chi }_\sigma ^{}>>`$ to $`O(T^4)`$ by making all possible contractions of operators $`\widehat{c}_{k\sigma }`$, $`\widehat{c}_{k^{}\sigma }^{}`$ in a third generation of Green functions. For instance, the EOM for $`<<\widehat{\chi }_\sigma ^{}\widehat{c}_{k^{}\sigma }\widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>`$ in eq.(7) is:
$`(\omega +ϵ_0ϵ_kϵ_k^{})<<\widehat{\chi }_\sigma ^{}\widehat{c}_{k^{}\sigma }\widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>=T_k^{}<<\widehat{n}_\sigma \widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>`$
$`T_k^{}<<\widehat{d}_\sigma ^{}\widehat{d}_\sigma \widehat{c}_{k^{}\sigma };\widehat{\chi }_\sigma ^{}>>+{\displaystyle \underset{k^{\prime \prime }}{}}T_{k^{\prime \prime }}<<\widehat{d}_\sigma ^{}\widehat{d}_\sigma \widehat{c}_{k^{}\sigma }\widehat{c}_{k^{\prime \prime }\sigma }^{}\widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>`$
$`{\displaystyle \underset{k^{\prime \prime }}{}}T_{k^{\prime \prime }}<<(1\widehat{n}_\sigma )\widehat{c}_{k^{\prime \prime }\sigma }^{}\widehat{c}_{k^{}\sigma }\widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>`$ (10)
which is approximated by
$`(\omega +ϵ_0ϵ_kϵ_k^{})<<\widehat{\chi }_\sigma ^{}\widehat{c}_{k^{}\sigma }\widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>=T_k^{}<<\widehat{n}_\sigma \widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>`$
$`T_k^{}<<\widehat{d}_\sigma ^{}\widehat{d}_\sigma \widehat{c}_{k^{}\sigma };\widehat{\chi }_\sigma ^{}>>+T_k<\widehat{n}_k><<\widehat{d}_\sigma ^{}\widehat{d}_\sigma \widehat{c}_{k^{}\sigma };\widehat{\chi }_\sigma ^{}>>`$
$`T_k^{}<\widehat{n}_k^{}><<(1\widehat{n}_\sigma )\widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>`$ (11)
This procedure yields the Green functions already obtained in eqs.(5-7) and a new one, $`<<\widehat{n}_\sigma \widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>`$, which has to be calculated as in eq.(7). The first generation of Green functions are thus obtained as:
$`(\stackrel{~}{\omega }_k\mathrm{\Sigma }_1(\stackrel{~}{\omega }_k))<<\widehat{d}_\sigma ^{}\widehat{d}_\sigma \widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>=<\widehat{\chi }_\sigma ^{}\widehat{c}_{k\sigma }>+`$
$`T_k<n_k><<\widehat{\chi }_\sigma ;\widehat{\chi }_\sigma ^{}>>T_k<\widehat{n}_k>{\displaystyle \underset{k^{}}{}}{\displaystyle \frac{T_k^{}<<\widehat{n}_\sigma \widehat{c}_{k^{}\sigma };\widehat{\chi }_\sigma ^{}>>}{\stackrel{~}{\omega }_kϵ_k^{}+ϵ_0}}`$ (12)
$`(\stackrel{~}{\omega }_k\mathrm{\Sigma }_2(\stackrel{~}{\omega }_k))<<\widehat{n}_\sigma \widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>=\mathrm{\Sigma }_3(\stackrel{~}{\omega }_k)<<\widehat{n}_\sigma \widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>+`$
$`T_k<\widehat{n}_k>{\displaystyle \underset{k^{}}{}}{\displaystyle \frac{T_k^{}<<\widehat{d}_\sigma ^{}\widehat{d}_\sigma \widehat{c}_{k^{}\sigma };\widehat{\chi }_\sigma ^{}>>}{\stackrel{~}{\omega }_kϵ_k^{}+ϵ_0}}`$ (13)
$`(\stackrel{~}{\omega }_k\mathrm{\Sigma }_2(\stackrel{~}{\omega }_k))<<\widehat{n}_\sigma \widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>=<\widehat{\chi }_\sigma ^{}\widehat{c}_{k\sigma }>+T_k<\widehat{n}_k><<\widehat{\chi }_\sigma ;\widehat{\chi }_\sigma ^{}>>+`$
$`\mathrm{\Sigma }_3(\stackrel{~}{\omega }_k)<<\widehat{n}_\sigma \widehat{c}_{k\sigma };\widehat{\chi }_\sigma ^{}>>T_k<\widehat{n}_k>{\displaystyle \underset{k^{}}{}}{\displaystyle \frac{T_k^{}<<\widehat{n}_\sigma \widehat{c}_{k^{}\sigma };\widehat{\chi }_\sigma ^{}>>}{\stackrel{~}{\omega }_kϵ_k^{}+ϵ_0}}`$ (14)
with $`\stackrel{~}{\omega }_k=\omega ϵ_k`$ from now on. Eqs.(12-14) define a set of coupled integral equations in the energy levels of the leads, $`ϵ_k`$, whose solution yields the Green functions to be introduced in eq.(4). In these equations we have neglected expectation values of operators such as $`<\widehat{\chi }_\sigma ^{}\widehat{\chi }_\sigma \widehat{c}_{k\sigma }^{}\widehat{c}_{k\sigma }>`$, which give small contributions to $`<<\widehat{\chi }_\sigma ;\widehat{\chi }_\sigma ^{}>>`$ to $`O(T^4)`$ and have defined several self-energies given by:
$$\mathrm{\Sigma }_i(\stackrel{~}{\omega }_k)=\underset{k^{}}{}T_k^{}^2\nu _i[\frac{1}{\stackrel{~}{\omega }_kϵ_k^{}+ϵ_0}+\frac{1}{\stackrel{~}{\omega }_k+ϵ_k^{}ϵ_0}]$$
(15)
where $`\nu _1=1<\widehat{n}_k^{}>`$, $`\nu _2=1`$ and $`\nu _3=<\widehat{n}_k^{}>`$, $`(\mathrm{\Sigma }_3(\stackrel{~}{\omega }_k)=\mathrm{\Sigma }_2(\stackrel{~}{\omega }_k)\mathrm{\Sigma }_1(\stackrel{~}{\omega }_k)`$). These self-energies are associated with different Green functions and play a crutial role in the Kondo-resonance decoherence, as we will see below. Eqs.(12-14) reproduce the conventional $`T^2`$-approximation of by neglecting all the terms proportional to $`T^3`$. One can get, however, a new (and effective) $`T^2`$-order near the Kondo resonance by neglecting the terms that couple these equations in an integral way. This approximation yields the following Green function:
$`<<\widehat{\chi }_\sigma ;\widehat{\chi }_\sigma ^{}>>=`$
$`[1<\widehat{n}_\sigma >{\displaystyle \underset{k}{}}{\displaystyle \frac{T_k<\widehat{\chi }_\sigma ^{}\widehat{c}_{k\sigma }>}{\stackrel{~}{\omega }_k\mathrm{\Sigma }_1(\stackrel{~}{\omega }_k)}}+{\displaystyle \underset{k}{}}T_k{\displaystyle \frac{\mathrm{\Sigma }_3(\stackrel{~}{\omega }_k)<\widehat{\chi }_\sigma ^{}\widehat{c}_{k\sigma }>}{(\stackrel{~}{\omega }_k\mathrm{\Sigma }_2(\stackrel{~}{\omega }_k))(\stackrel{~}{\omega }_k\mathrm{\Sigma }_4(\stackrel{~}{\omega }_k))}}]`$
$`[\omega ϵ_0\mathrm{\Sigma }_0{\displaystyle \underset{k}{}}{\displaystyle \frac{T_k^2<\widehat{n}_k>}{\stackrel{~}{\omega }_k\mathrm{\Sigma }_1(\stackrel{~}{\omega }_k)}}+{\displaystyle \underset{k}{}}{\displaystyle \frac{T_k^2<\widehat{n}_k>\mathrm{\Sigma }_3(\stackrel{~}{\omega }_k)}{(\stackrel{~}{\omega }_k\mathrm{\Sigma }_2(\stackrel{~}{\omega }_k))(\stackrel{~}{\omega }_k\mathrm{\Sigma }_4(\stackrel{~}{\omega }_k))}}]^1`$ (16)
with
$$\mathrm{\Sigma }_4(\stackrel{~}{\omega }_k)=\mathrm{\Sigma }_2(\stackrel{~}{\omega }_k)+\frac{\mathrm{\Sigma }_3(\stackrel{~}{\omega }_k)^2}{\stackrel{~}{\omega }_k\mathrm{\Sigma }_2(\stackrel{~}{\omega }_k)}.$$
(17)
In eq.(16), $`<n_\sigma >`$ and $`<\widehat{\chi }_\sigma ^{}\widehat{c}_{k\sigma }>`$ have to be calculated self-consistently ,. At zero temperature and in a wide-band approximation, they are given by:
$$<\widehat{n}_\sigma >=_{\mathrm{}}^0\frac{d\omega }{\pi }Im<<\widehat{\chi }_\sigma ;\widehat{\chi }_\sigma ^{}>>(\omega )$$
(18)
and
$$<\widehat{\chi }_\sigma ^{}\widehat{c}_{k\sigma }>=_{\mathrm{}}^0\frac{d\omega }{\pi }Im\frac{T_k<<\widehat{\chi }_\sigma ;\widehat{\chi }_\sigma ^{}>>(\omega )}{\stackrel{~}{\omega }_k}$$
(19)
where $`\omega =0`$ is the Fermi level.
## 3 Results
In this section we first discuss some consequences that our approximation bears in the description of the Kondo peak. Notice that the $`T^2`$-order of eq.(8) is given by eq.(16) neglecting the last terms of the numerator and denominator and taking $`\mathrm{\Sigma }_1(\stackrel{~}{\omega }_k)=0`$. These new terms yield, however, contributions that are proportional to $`T^2`$ around the Kondo peak. This can be seen by considering the self-energies $`\mathrm{\Sigma }_1`$, $`\mathrm{\Sigma }_2`$ and $`\mathrm{\Sigma }_4`$, whose imaginary parts tend to broaden the Kondo-like peak associated with the factors $`_k\frac{T_k^2<\widehat{n}_k>}{\stackrel{~}{\omega }_k\mathrm{\Sigma }_i(\stackrel{~}{\omega }_k)}`$. In the $`T^4`$-order, eq.(15) yields, for $`V=0`$, $`ϵ_0<0`$, $`|ϵ_0|>>\mathrm{\Delta }`$ ($`2\mathrm{\Delta }`$ being the one-electron level-width) and a wide-band approximation, $`\mathrm{\Sigma }_2=2\mathrm{\Sigma }_02i\mathrm{\Delta }`$, $`\mathrm{\Sigma }_3\mathrm{\Sigma }_2`$ and $`\mathrm{\Sigma }_1=0`$, if $`|\omega ϵ_k|<<|ϵ_0|`$. Then, near the Fermi level, we obtain:
$$\underset{k}{}\frac{T_k^2<\widehat{n}_k>}{\stackrel{~}{\omega }_k\mathrm{\Sigma }_1(\stackrel{~}{\omega }_k)}\frac{\mathrm{\Delta }}{\pi }Ln(\frac{D}{\omega })$$
(20)
D being the half-bandwidth, and
$$\underset{k}{}\frac{T_k^2<\widehat{n}_k>}{\stackrel{~}{\omega }_k\mathrm{\Sigma }_2(\stackrel{~}{\omega }_k)}\frac{\mathrm{\Sigma }_3(\stackrel{~}{\omega }_k)}{\stackrel{~}{\omega }_k\mathrm{\Sigma }_4(\stackrel{~}{\omega }_k)}\frac{\mathrm{\Delta }}{2\pi }Ln(\frac{\omega +4i\mathrm{\Delta }}{\omega }).$$
(21)
The important point is that the term given by eq.(21), which is neglected in the $`T^2`$-order approximation, also gives a $`T^2`$-contribution to $`<<\widehat{\chi }_\sigma ;\widehat{\chi }_\sigma ^{}>>`$ going like $`Ln(\omega )`$ at the Fermi energy. In particular, eq.(16) yields the Kondo-energy
$$T_KD^{2/3}(4\mathrm{\Delta })^{1/3}exp[\frac{2\pi |ϵ_0|}{3\mathrm{\Delta }}]$$
(22)
intermediate between $`Dexp[\pi |ϵ_0|/\mathrm{\Delta }]`$ (the $`T^2`$-order) and the exact result of Haldane $`(D\mathrm{\Delta })^{1/2}exp[\pi |ϵ_0|/2\mathrm{\Delta }]`$ .
A check of our solution, eq.(16), can be obtained by considering the case $`V=0`$; this equilibrium case can be compared with the NCA and with the numerically exact calculations of Ref. using the NRG. Fig.1 shows our calculated spectral density for the Kondo regime defined by $`ϵ_0=4\mathrm{\Delta }`$, $`\mathrm{\Delta }=0.01`$ and $`D=1`$ at zero temperature. Two cases, $`T^2`$-order and eq.(16) are shown as a function of $`\omega /T_K`$. Although for this particular case our value of $`T_K`$, eq.(22), is very similar to Haldane’s, we use this last one from now on. Our results indicate that eq.(16) yields a much better spectral density for the Kondo-resonance than the $`T^2`$-order aproximation, with values very similar to the ones given by the NRG above the Fermi energy, although our results yield a too small value of the spectral density below the Fermi energy; this deviation extends to energies of around $`2T_K`$ and suggests that a much more accurate solution is going to be obtained for voltages larger than this value (see below). Fig.2 shows the same results for a mixed-valence regime ($`ϵ_0=0`$, $`\mathrm{\Delta }=0.01`$ and $`D=1`$) at zero temperature. In this case, eq.(16) yields an excellent approximation to the exact NRG-calculation and the Friedel-Langreth sum rule is satisfied with an accuracy better than 10$`\%`$. The agreement is even better in the empty orbital regime, Fig.3, where the sum rule is fulfilled to within 5$`\%`$ accuracy. We should mention that our present results are in much better agreement with the exact results than the ones obtained using Lacroix’s approximation and shown in .
Now we turn our attention to the nonequilibrium case $`V0`$ in the symmetric situation, where the chemical potentials of the left and right leads are $`\pm V/2`$. A fully consistent solution of the case $`V0`$ requires the application of the Keldysh formalism, as explained in . The symmetric case can be analyzed, however, at the level of the above equations because, to good accuracy, $`<n_\sigma >`$ and $`<\widehat{\chi }_\sigma ^{}\widehat{c}_{k\sigma }>`$ do not vary with V for the values used here. We start by analyzing the self-energies $`\mathrm{\Sigma }_1`$ and $`\mathrm{\Sigma }_4`$ in the regime $`\stackrel{~}{\omega }_k0`$, which is the relevant regime for analyzing the Kondo-resonance decoherence. It is easy to see that eq.(15) produces $`\mathrm{\Sigma }_1(\stackrel{~}{\omega }_k0)0`$ , as in the equilibrium case and therefore, we find no decoherence for $`V0`$. It is interesting to proceed further with the EOM method and see how these self-energies are modified. By going to $`O(T^6)`$, we find for $`\mathrm{\Sigma }_1`$ and $`\mathrm{\Sigma }_2`$:
$`\mathrm{\Sigma }_1(\stackrel{~}{\omega }_k)={\displaystyle \underset{k^{}}{}}T_k^{}^2[{\displaystyle \frac{1<\widehat{n}_k^{}>}{\stackrel{~}{\omega }_kϵ_k^{}+ϵ_{}(k^{})+i\mathrm{\Delta }}}+{\displaystyle \frac{1<\widehat{n}_k^{}>}{\stackrel{~}{\omega }_k+ϵ_k^{}ϵ_+(k^{})+i\mathrm{\Delta }}}]+`$
$`2{\displaystyle \underset{k^{},k^{\prime \prime }}{}}{\displaystyle \frac{T_k^{}^2(1<\widehat{n}_k^{}>)}{\stackrel{~}{\omega }_kϵ_k^{}+ϵ_{}(k^{})+i\mathrm{\Delta }}}{\displaystyle \frac{1}{\stackrel{~}{\omega }_k+ϵ_{k^{\prime \prime }}ϵ_k^{}}}{\displaystyle \frac{T_{k^{\prime \prime }}^2(1<\widehat{n}_{k^{\prime \prime }}>)}{\stackrel{~}{\omega }_k+ϵ_{k^{\prime \prime }}ϵ_+(k^{\prime \prime })+i\mathrm{\Delta }}}+`$
$`{\displaystyle \underset{k^{},k^{\prime \prime },q}{}}{\displaystyle \frac{T_k^{}^2(1<\widehat{n}_k^{}>)}{\stackrel{~}{\omega }_kϵ_k^{}+ϵ_{}(k^{})+i\mathrm{\Delta }}}{\displaystyle \frac{1}{\stackrel{~}{\omega }_k+ϵ_{k^{\prime \prime }}ϵ_k^{}}}{\displaystyle \frac{T_{k^{\prime \prime }}^2(1<\widehat{n}_{k^{\prime \prime }}>)}{\stackrel{~}{\omega }_k+ϵ_{k^{\prime \prime }}ϵ_+(k^{\prime \prime })+i\mathrm{\Delta }}}\times `$
$`{\displaystyle \frac{1}{\stackrel{~}{\omega }_k+ϵ_{k^{\prime \prime }}ϵ_q}}{\displaystyle \frac{T_q^2(1<\widehat{n}_q>)}{\stackrel{~}{\omega }_kϵ_q+ϵ_{}(q)+i\mathrm{\Delta }}}+O(T^8)`$ (23)
and
$`\mathrm{\Sigma }_2(\stackrel{~}{\omega }_k)={\displaystyle \underset{k^{}}{}}T_k^{}^2[{\displaystyle \frac{1}{\stackrel{~}{\omega }_kϵ_k^{}+ϵ_{}(k^{})+i\mathrm{\Delta }}}+{\displaystyle \frac{1}{\stackrel{~}{\omega }_k+ϵ_k^{}ϵ_+(k^{})+i\mathrm{\Delta }}}]+`$
$`2{\displaystyle \underset{k^{},k^{\prime \prime }}{}}{\displaystyle \frac{T_k^{}^2<\widehat{n}_k^{}>}{\stackrel{~}{\omega }_kϵ_k^{}+ϵ_{}(k^{})+i\mathrm{\Delta }}}\times {\displaystyle \frac{1}{\stackrel{~}{\omega }_k+ϵ_{k^{\prime \prime }}ϵ_k^{}}}\times {\displaystyle \frac{T_{k^{\prime \prime }}^2<\widehat{n}_{k^{\prime \prime }}>}{\stackrel{~}{\omega }_k+ϵ_{k^{\prime \prime }}ϵ_+(k^{\prime \prime })+i\mathrm{\Delta }}}+O(T^8)`$ (24)
where $`ϵ_\pm (ϵ_k^{})=ϵ_0\frac{\mathrm{\Delta }}{\pi }Ln(\frac{ϵ_k^{}\pm \stackrel{~}{\omega }_k}{D})`$. We also get $`\mathrm{\Sigma }_3(\stackrel{~}{\omega }_k)=\mathrm{\Sigma }_2(\stackrel{~}{\omega }_k)\mathrm{\Sigma }_1(\stackrel{~}{\omega }_k)`$, $`\mathrm{\Sigma }_4(\stackrel{~}{\omega }_k)=\mathrm{\Sigma }_2(\stackrel{~}{\omega }_k)+\frac{\mathrm{\Sigma }_3(\stackrel{~}{\omega }_k)^2}{\stackrel{~}{\omega }_k\mathrm{\Sigma }_2(\stackrel{~}{\omega }_k)}`$, as before. These higher order self-energies give rise to the first non-zero contribution to the Kondo resonance decoherence, due to the cross terms between the right and left leads, as has already been pointed out in . Eqs.(23) and (24) represents the first terms of an expansion for $`\mathrm{\Sigma }_1`$ or $`\mathrm{\Sigma }_2`$, using $`(T_k/ϵ_0)^2`$ as the expansion parameter, and allows us to analyze the Kondo resonance decoherence as a function of an external bias. For $`V=0`$, we again find $`\mathrm{\Sigma }_1(\stackrel{~}{\omega }_k0)0`$ and $`\mathrm{\Sigma }_4(\stackrel{~}{\omega }_k0)0`$ but for $`V0`$ these self-energies develop non-negligible imaginary parts for $`\stackrel{~}{\omega }_k0`$. In the bias regime, defined by $`V<<|\stackrel{~}{ϵ}_0|<\mathrm{\Delta }`$, we find
$$\mathrm{\Sigma }_1(\stackrel{~}{\omega }_k0;V)i(\frac{\mathrm{\Delta }}{\stackrel{~}{ϵ}_0})^2\frac{V}{2\pi }(12\frac{\mathrm{\Delta }}{\pi |\stackrel{~}{ϵ}_0|}Ln\frac{V}{|\stackrel{~}{ϵ}_0|})i\gamma $$
(25)
and $`\mathrm{\Sigma }_2i2\mathrm{\Delta }`$. In eq.(25), $`\stackrel{~}{ϵ}_0`$ is a renormalized level given approximately by $`\stackrel{~}{ϵ}_0ϵ_0+\frac{\mathrm{\Delta }}{\pi }Ln\frac{\pi D}{\mathrm{\Delta }}`$ . Fig.4 shows the evolution of the Kondo-like resonance spectral density for the same parameters used in Fig.1, but introducing different biases on the QD. Here we have used eq.(16) and the self-energies to $`T^6`$-order. Our results are compared with the NCA-calculations of Ref. , where available. The agreement between both approaches, for $`\omega >0`$, is remarkable. The difference between both cases can be attributed to the bandwidth used in , which is ten times smaller than ours. For $`\omega <0`$, however, the agreement is only reasonable at low bias (like in the equilibrium case shown in Fig.1), which improves as $`V`$ increases.
We have also calculated the QD-conductance which, in the symmetric configuration we are considering, is basically proportional to the mean value of the spectral density in the interval $`(V/2,V/2)`$. Fig.5 shows the conductance $`G`$ in units of the conductance quantum $`G_0=2e^2/h`$, as a function of $`V/T_K`$; a good fit is obtained by the expression:
$$G/G_0=1/(1+0.75Ln(1+(\frac{V}{2T_K})^2)),$$
(26)
for large values of $`V/T_K`$. These values also show good agreement with the NCA-calculations of for $`V>2T_K`$, where our solution can be taken with great confidence, although in this case, $`G/G_0`$ behaves like $`1/Ln^2(V/T_K)`$ for very high voltages, $`Ln(V/T_K)>>1`$.
In addition, we have also addressed the important issue of strong versus weak coupling regimes, as a function of the bias. It has been argued that the strong-coupling regime might subsist, even for high bias ($`V>>T_K`$), if the decoherence rate $`\gamma `$ is smaller than $`T^{}=V/2+\sqrt{T_K^2+(V/2)^2}`$. With $`\gamma `$ given by eq.(25), the strong-coupling regime is obtained for $`V\pi ^{1/2}\frac{|\stackrel{~}{ϵ}_0|}{\mathrm{\Delta }}T_K`$. This is a very stringent condition; for instance, for the values of the parameters in Figs.1a and 2, we find strong coupling for $`V3.5T_K`$, for which the conductance is larger than $`0.4G_0`$. We conclude that, when subjected to an external bias $`V>T_K`$, the QD-Kondo resonance develops a two-peaked structure at $`\pm V/2`$ for $`V>2T_K`$; the system, however, evolves quickly as a function of the voltage, moving from a strong-coupling regime with the spin impurity screened by the electrons of the leads, to a weak-coupling regime with that spin unscreened.
## 4 Conclusions
The EOM method provides an accurate approximation to the Kondo problem in and out of equilibrium situations, the only requirement being to include all the Green functions in a strict expansion in $`T^2`$. Our results to $`T^4`$-order in equilibrium provide a good description of the Kondo peak in the Kondo regime and an excellent aproximation to (numerically) exact results in the mixed-valence and empty-orbital regimes. In the nonequilibrium regime, it is necessary to carry the method to $`O(T^6)`$ to find the Kondo resonance decoherence, since it involves the coupling between the right and left electrodes through the QD. In particular, in our approach, the decoherence rate is found analytically to be nearly proportional to $`V`$, which indicates that the system is driven readily to the weak couping regime if $`V\pi ^{1/2}\frac{|\stackrel{~}{ϵ}_0|}{\mathrm{\Delta }}T_K`$, with the spin impurity unquenched by the electrons of the leads.
## 5 Acknowledgements
We are indebted to E.A. Anda and E.C. Goldberg for many fruitful discussions, to S. Davison for a critical reading of the manuscript and to T. Costi and P. Nordlander for providing us with their data. This work has been funded by the Spanish Comisión Interministerial de Ciencia y Tecnología under contracts MAT-2001-0665 and BFM-2001-0150. |
warning/0507/math0507562.html | ar5iv | text | # Elementary elliptic (𝑅,𝑞)-polycycles
## 1 Introduction
Given $`q`$ and $`R`$, a $`(R,q)`$-polycycle is a non-empty $`2`$-connected plane, locally finite (i.e. any circle contain only finite number of its vertices) graph $`G`$ with faces partitioned in two non-empty sets $`F_1`$ and $`F_2`$, so that:
(i) all elements of $`F_1`$ (called proper faces) are combinatorial $`i`$-gons with $`iR`$;
(ii) all elements of $`F_2`$ (called holes, the exterior face(s) <sup>1</sup><sup>1</sup>1Any finite plane graph has an unique exterior face; any infinite plane graph can have any number of exterior faces, including $`0`$ and infinity ($`2`$-gonal faces are permitted). are amongst them) are pair-wisely disjoint, i.e. have no common vertices;
(iii) all vertices have degree within $`\{2,\mathrm{},q\}`$ and all interior (i.e. not on the boundary of a hole) vertices are $`q`$-valent.
The plane graph $`G`$ can be finite or infinite and some of the faces of the set $`F_2`$ can be $`i`$-gons with $`iR`$. Two $`(R,q)`$-polycycles, which are isomorphic as plane graphs, but have different pairs $`(F_1,F_2)`$, will be considered non-isomorphic in our context. The symmetry group $`Aut(P)`$ of a polycycle $`P`$, considered below, consists of all automorphisms of plane graph $`G`$ preserving the pair $`(F_1,F_2)`$. Note that it is different from $`Aut(G)`$, the full automorphism group of the plane graph $`G`$, i.e. the group of transformations, which preserves the edge-set and the face-set of the plane graph $`G`$. In fact, $`Aut(P)`$ is the stabilizer of the pair $`(F_1,F_2)`$ in $`Aut(G)`$.
The notion of $`(R,q)`$-polycycle is a large generalization of the case $`|R|=|F_2|=1`$, i.e. $`(r,q)`$-polycycle introduced by Deza and Shtogrin in \[DS98\] and studied in their papers \[DS98\], \[DS99\], \[DS00a\], \[DS00b\], \[DS00c\], \[DS01\], \[DS02a\], \[DS02b\], \[DS04\], \[DS05a\], \[DS05b\], \[Sh99\], \[Sh00\]. The case $`|R|=1`$, i.e. $`(r,q)`$-polycycles with holes, was considered in \[DDS05b\].
A boundary of a $`(R,q)`$-polycycle $`P`$ is the boundary of any of its holes.
A bridge of a $`(R,q)`$-polycycle is an edge, which is not on a boundary and goes from a hole to a hole (possibly, the same).
An $`(R,q)`$-polycycle is called elementary if it has no bridges.
An open edge of a $`(R,q)`$-polycycle is an edge on a boundary, such that each of its end-vertices have degree less than $`q`$.
###### Theorem 1
Every $`(R,q)`$-polycycle is uniquely formed by the agglomeration of elementary ones along open edges or, in other words, it can be uniquely cut, along the bridges, into the elementary ones.
Hence, the interesting question is to enumerate those elementary $`(R,q)`$-polycycles. Call a $`(R,q)`$-polycycle elliptic, parabolic or hyperbolic if the number $`\frac{1}{q}+\frac{1}{r}\frac{1}{2}`$ (where $`r=max_{iR}i`$) is positive, zero or negative, respectively. The number of elementary $`(r,q)`$-polycycles is uncountable for any parabolic or hyperbolic pairs $`(r,q)`$. But in \[DS01\] and \[DS02b\] all elliptic elementary $`(r,q)`$-polycycles were determined. Namely, the countable set of all elementary $`(5,3)`$\- and $`(3,5)`$-polycycles was described; the cases of $`(3,3)`$-, $`(4,3)`$\- and $`(3,4)`$-polycycles are easy. In \[DDS05b\] all elliptic elementary $`(\{r\},q)`$-polycycles (i.e. $`(r,q)`$-polycycles with holes) were determined as a part of a more elaborate classification.
The purpose of this paper is to generalize it for elliptic $`(R,q)`$-polycycles. In fact, we will consider the main case $`R=\{i\text{ : }2ir\}`$ and so, all such elliptic possibilities are $`(\{2,3,4,5\},3)`$-, $`(\{2,3\},4)`$\- and $`(\{2,3\},5)`$-polycycles.
Given a $`(R,q)`$-polycycle $`P`$, one can define another $`(R,q)`$-polycycle $`P^{}`$ by removing a face $`f`$ from $`F_1`$, i.e. by considering it as a hole. If $`f`$ has no common vertices with other faces from $`F_1`$, then removing of it leaves unchanged the plane graph $`G`$ and only changes the pair $`(F_1,F_2)`$. If $`f`$ has some edges in common with a hole, then we remove them. If $`f`$ has a common vertex with a hole, then we split it in two. See below this operation for a $`2`$-gon, which is incident to a vertex on the boundary.
A $`(R,q)`$-polycycle $`P`$ is called extensible if there exists another $`(R,q)`$-polycycle $`P^{}`$, such that the elimination of a face of $`P^{}`$ yields $`P`$.
The same plane graph $`G`$ can admit several polycycle realizations. For example, $`Prism_3`$ admits two realizations as a $`(\{3,4\},3)`$-polycycle (see $`(4,2022,D_{3h})`$ below). For elementary $`(\{2,3,4,5\},3)`$-polycycles, we give the list in the form $`(n_f,nb_1\mathrm{}nb_p,Aut(G))`$ with $`n_f`$ being the number of proper faces (from $`F_1`$), and $`nb_1,\mathrm{},nb_p`$ being the position in the corresponding lists of elementary polycycles with $`n_f`$ faces and $`Aut(G)`$ being the automorphism group of plane graph $`G`$: $`(4,1213,C_s)`$, $`(4,2022,D_{3h})`$, $`(4,235^{th}ofLemma\mathrm{\hspace{0.17em}1},C_{2\nu })`$, $`(5,262830,C_{2\nu })`$, $`(5,2729,C_s)`$, $`(6,212231,C_{3\nu })`$, $`(6,2834,D_{5h})`$, $`(7,2436,D_{3d})`$, $`(7,2931,C_{2\nu })`$, $`(8,2529,D_{3h})`$, $`(9,1116,D_{4d})`$. For elementary $`(\{2,3\},5)`$-polycycles, this concerns the 10th and 20th of Figure 3 ($`Aut(G)=C_s`$), the 21st of Figure 3 and 1st of Figure 4 ($`Aut(G)=D_3`$), the 14th and 22nd of Figure 4 ($`Aut(G)=D_{2d}`$). If $`Aut(P)Aut(G)`$ and no other polycycle with the same graph $`G`$ exists, we indicate it by putting the second group in parenthesis.
Theorem 1, together with the determination of the elementary $`(R,q)`$-polycycles, is especially useful in extremal problems, related to the number of interior/exterior vertices, see \[DS01\] and \[DS02b\], and in classification of face-regular two-faced maps, see \[DD05\].
We thank Gil Kalai for putting a question, which lead us to this study.
## 2 Classification of elementary $`(\{2,3,4,5\},3)`$-polycycles
A $`(R,3)`$-polycycle is called totally elementary if it is elementary and if, after removing any face adjacent to a hole, one obtains a non-elementary $`(R,3)`$-polycycle. So, an elementary $`(R,3)`$-polycycle is totally elementary if and only if it is not the result of an extension of some elementary $`(R,3)`$-polycycle.
We will classify those polycycles in a number of steps. At first, we find all totally elementary $`(\{2,3,4,5\},3)`$-polycycles. Then, all other elementary ones are obtained by adding faces to the existing elementary or totally elementary polycycles.
###### Theorem 2
The list of totally elementary $`(\{3,4,5\},3)`$-polycycles consists of:
1. three isolated $`i`$-gons, $`i\{3,4,5\}`$<sup>2</sup><sup>2</sup>2In fact, they belong to the infinite series of $`(r,2)`$-polycycles (monocycles) consisting of isolated $`r`$-gons, $`r2`$.:
$`C_{3\nu }`$ $`(D_{3h})`$
$`C_{4\nu }`$ $`(D_{4h})`$
$`C_{5\nu }`$ $`(D_{5h})`$
2. all ten triples of $`i`$-gons, $`i\{3,4,5\}`$:
$`C_{3\nu }`$, nonext. $`(T_d)`$
$`C_s`$, nonext. $`(C_{2\nu })`$
$`C_s`$ $`(C_{2\nu })`$
$`C_1`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_{3\nu }`$
$`C_{3\nu }`$
3. the following doubly infinite $`(\{5\},3)`$-polycycle, denoted by $`Barrel_{\mathrm{}}`$:
pma2, nonext.
4. the infinite series<sup>3</sup><sup>3</sup>3The $`Barrel_m`$ is a $`3`$-valent plane graph, consisting of two $`m`$-gons separated by two $`m`$-rings of $`5`$-gons. of $`Barrel_m`$, $`m2`$ (two $`m`$-gonal holes, non-extensible for $`m3,4,5`$, with symmetry $`D_{md}`$), represented below for $`m=2,3,4,5`$:
$`D_{2d}`$
$`D_{3d}`$
$`D_{4d}`$
$`D_{5d}`$ $`(I_h)`$
Proof. Take a totally elementary $`(\{2,3,4,5\},3)`$-polycycle $`P`$. If $`|F_1|=1`$, then $`P`$ is, clearly, totally elementary; so, let us assume that $`|F_1|2`$. If $`|F_1|=2`$, then it is, clearly, not elementary; so, assume $`|F_1|3`$. Of course, $`P`$ has at least one interior vertex; let $`v`$ be such a vertex. Furthermore, one can assume that $`v`$ is adjacent to a vertex $`v^{}`$, which is incident to a hole.
The vertex $`v`$ is incident to three faces $`f_1`$, $`f_2`$, $`f_3`$. Let us denote by $`v_{ij}`$ unique vertex incident to $`f_i`$, $`f_j`$ and adjacent to $`v`$. Without loss of generality, one can suppose that $`v^{}`$ is incident to the faces $`f_1`$ and $`f_2`$, i.e. that $`v^{}=v_{12}`$.
The removal of the face $`f_1`$ yields a non-elementary polycycle; so, there is at least one bridge separating $`Pf_1`$ in two parts. Such bridge should have an end-vertex incident to $`f_1`$. The same holds for $`f_2`$. The proof consists of a number of cases.
$`1^{st}`$ case: If $`e_1=\{v,v_{23}\}`$ and $`e_2=\{v,v_{13}\}`$ are bridges for $`Pf_1`$, $`Pf_2`$, respectively, then from the constraint that faces $`f_i`$ are $`p`$-gons with $`p5`$, one sees that each face $`f_i`$ is adjacent in $`P`$ to at most one other face. Furthermore, if $`f_i`$ is adjacent to another face, then this adjacence is along a bridge, which is forbidden. Hence, $`F_1=\{f_1,f_2,f_3\}`$.
$`2^{nd}`$ case: Let us assume now that $`e_1=\{v,v_{23}\}`$ is a bridge for $`Pf_1`$, but $`e_2=\{v,v_{13}\}`$ is not a bridge for $`Pf_2`$. Then, since $`f_2`$ is a $`p`$-gon with $`p5`$, it is adjacent to at most one other face and, if so, then along a bridge, which is impossible. So, $`f_2`$ is adjacent to only $`f_1`$ and $`f_3`$ and, since $`e_2`$ is not a bridge for $`Pf_2`$, one obtains that $`Pf_2`$ is elementary, which contradicts the hypothesis.
$`3^{rd}`$ case: Let us assume that neither $`e_1`$, nor $`e_2`$ are bridges for $`Pf_1`$ and $`Pf_2`$. From the consideration of previous two cases, one has that every vertex $`v`$, adjacent to a vertex on the boundary, is in this $`3^{rd}`$ case.
The first subcase, which can happen only if $`f_1`$ is $`5`$-gon, happens, when the over-lined edge $`e^{}`$, in the drawing below, is a bridge.
The face $`g`$ is adjacent to the faces $`h`$ and $`f_1`$ and, possibly, to another face $`g^{}`$. But if $`g`$ is adjacent to such a face $`g^{}`$, it is along a bridge of $`P`$; hence, $`g`$ is adjacent only to $`h`$ and $`f_1`$. So, $`Pg`$ is elementary, which is impossible.
So, the edge $`e^{}`$ is not a bridge and this forces the face $`h`$ to be $`5`$-gonal. Hence, the vertex $`v_{13}`$ is in the same situation as the vertex $`v`$, described in the diagram below:
So, one can repeat the construction. If, at some point, $`e_1`$ is a bridge, then the construction stops; otherwise, one can continue indefinitely. If one obtains only different vertices, then it means that we have the $`(\{5\},3)`$-polycycle $`Barrel_{\mathrm{}}`$; otherwise, we obtain a loop of vertices, i.e. a circuit of vertices, which appear over and over, i.e. $`Barrel_m`$ for some $`m2`$. $`\mathrm{}`$
###### Lemma 1
All elementary $`(\{2,3,4,5\},3)`$-polycycles, containing a $`2`$-gon, are following eight ones:
$`C_{2\nu }`$ $`(D_{2h})`$
$`C_{2\nu }`$, nonext. $`(D_{2h})`$
$`C_{2\nu }`$
$`C_s`$, nonext.
$`C_s`$, nonext. $`(C_{2\nu })`$
$`C_s`$, nonext. $`(C_{2\nu })`$
$`C_s`$ $`(C_{2\nu })`$
$`C_{2\nu }`$, nonext. $`(D_{2d})`$
Proof. Let $`P`$ be such polycycle. Clearly, the $`2`$-gon is the only possibility if $`|F_1|=1`$. If $`|F_1|=2`$, then it is not elementary. If $`|F_1|3`$, then the $`2`$-gon should be inside of the structure. So, $`P`$ contain, as a subgraph, one of three following graphs:
So, the only possibilities for $`P`$ are those given in above Lemma. $`\mathrm{}`$
We will enumerate now all elementary $`(\{3,4,5\},3)`$-polycycles. If such a polycycle is not totally elementary, then it is obtained from another elementary $`(\{3,4,5\},3)`$-polycycle (totally elementary or not) by addition of another face.
###### Theorem 3
The list of elementary $`(\{2,3,4,5\},3)`$-polycycles consists of:
1. eight $`(\{2,3,4,5\},3)`$-polycycles, given in the list of Lemma 1.
2. the sporadic $`(\{3,4,5\},3)`$-polycycles, given for $`1`$ and $`3`$ faces in Theorem 2 (i),(ii) and the remainder, after proof of this Theorem.
3. the infinite series of $`(\{5\},3)`$-polycycles $`Barrel_m`$ with $`2m\mathrm{}`$, given in Theorem 2 (iii), (iv).
4. six $`(\{3,4,5\},3)`$-polycycles, infinite in one direction:
$`\alpha `$: $`C_1`$
$`\delta `$: $`C_1`$
$`\beta `$: $`C_1`$
$`\epsilon `$: $`C_1`$, nonext.
$`\gamma `$: $`C_1`$, nonext.
$`\mu `$: $`C_1`$, nonext.
5. $`21=\left(\genfrac{}{}{0pt}{}{6+1}{2}\right)`$ infinite series obtained by taking two endings of the infinite polycycles of (iv) above and concatenating them.
For example, merging of $`\alpha `$ with itself produces the infinite series of elementary $`(\{5\},3)`$-polycycles, denoted by $`E_n`$ in \[DS01\], $`0n\mathrm{}`$. For $`n=0,1,2,3,\mathrm{}`$, it is the $`5`$-gon (Theorem 2 (i)), the triple of $`5`$-gons (Theorem 2 (ii)), 16th in the list for $`4`$ faces below, $`22th`$ in the list for $`5`$ faces below, $`Barrel_{\mathrm{}}`$. See Figure 1 for the first $`3`$ members (starting with 6 faces) of two such series: $`\alpha \alpha `$ and $`\beta \epsilon `$.
Proof. The proof consists of taking the totally elementary polycycles, given in Theorem 2, adding a face with right number of sides, which preserves the elementarity in all possible ways. Then we reduce, by isomorphism of $`(R,q)`$-polycycles, and obtain the list of finite elementary $`(\{2,3,4,5\},3)`$-polycycles with one hole. If a $`(\{3,4,5\},3)`$-polycycle has two holes and is not a $`Barrel_m`$, then it is not elementary. So, it can be obtained from another elementary polycycle with one hole less, by the addition of one face. It is easy to see that this cannot happen. So, we have the complete list of finite $`(\{2,3,4,5\},3)`$-polycycles.
Take now an elementary $`(\{3,4,5\},3)`$-polycycle $`P`$, which is infinite. Remove all $`3`$\- or $`4`$-gonal faces of it. The result is a $`(\{5\},3)`$-polycycle $`P^{}`$, which is not necessarily elementary. We will now use the classification of elementary $`(\{5\},3)`$-polycycles (possibly, infinite) done in \[DDS05b\]. If the infinite $`(\{5\},3)`$-polycycle $`Barrel_{\mathrm{}}`$ appears in the decomposition, then, clearly, $`P`$ is reduced to it. If the infinite polycycle $`\alpha `$ appear in the decomposition, then there are two possibilities for extending it, indicated below.
If a $`3`$\- or $`4`$-gonal face is adjacent on the dotted line, then there should be another face on the over-line edges. So, in any case, there is a face, adjacent on the over-line edges, and we can assume that it is a $`3`$\- or $`4`$-gonal face. Then, consideration of all possibilities to extend it, yields $`\beta `$, …, $`\mu `$. Suppose now, that $`P`$ does not contain any infinite $`(\{5\},3)`$-polycycles. Then we can find an infinite path $`f_0,\mathrm{},f_i,\mathrm{}`$ of distinct faces of $`P`$ in $`F_1`$, such that $`f_i`$ is adjacent to $`f_{i+1}`$ and $`f_{i1}`$ is not adjacent to $`f_{i+1}`$. The condition on $`P`$ implies that an infinite number of faces are $`3`$\- or $`4`$-gons, but the condition of non-adjacency of $`f_{i1}`$ with $`f_{i+1}`$ forbids $`3`$-gons. Take now a $`4`$-gon $`f_i`$ and assume that $`f_{i1}`$ and $`f_{i+1}`$ are $`5`$-gons. The consideration of all possibilities of extension around that face, lead us to an impossibility. If some of $`f_{i1}`$ or $`f_{i+1}`$ are $`4`$-gons, then we have a path of $`4`$-gons and the case is even simpler. $`\mathrm{}`$
List of elementary $`(\{3,4,5\},3)`$-polycycles with $`4`$ faces:
$`C_1`$
$`C_1`$
$`C_1`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_s`$, nonext. $`(C_{2\nu })`$
$`C_{2\nu }`$
$`C_{2\nu }`$
$`C_{2\nu }`$
$`C_{2\nu }`$
$`C_{2\nu }`$
$`C_{2\nu }`$, nonext.
$`C_{3\nu }`$
$`C_{3\nu }`$, nonext.
$`C_{2\nu }`$, nonext.
List of elementary $`(\{3,4,5\},3)`$-polycycles with $`5`$ faces:
$`C_1`$
$`C_1`$
$`C_1`$
$`C_1`$
$`C_1`$
$`C_1`$
$`C_1`$
$`C_1`$
$`C_1`$
$`C_1`$
$`C_1`$
$`C_1`$
$`C_1`$
$`C_1`$
$`C_1`$
$`C_1`$, nonext.
$`C_1`$, nonext.
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_s`$, nonext. $`(C_{2\nu })`$
$`C_s`$, nonext. $`(C_{2\nu })`$
$`C_{2\nu }`$
$`C_{4\nu }`$
$`C_{4\nu }`$, nonext. $`(O_h)`$
List of sporadic elementary $`(\{3,4,5\},3)`$-polycycles with $`6`$ faces:
$`C_1`$
$`C_1`$
$`C_1`$
$`C_1`$, nonext.
$`C_1`$, nonext.
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_s`$, nonext. $`(C_{2\nu })`$
$`C_{2\nu }`$
$`C_{2\nu }`$, nonext.
$`C_{3\nu }`$
$`C_{3\nu }`$
$`C_{3\nu }`$, nonext.
$`C_{3\nu }`$, nonext.
$`C_{5\nu }`$
$`C_{5\nu }`$, nonext.
List of sporadic elementary $`(\{3,4,5\},3)`$-polycycles with $`7`$ faces:
$`C_1`$
$`C_1`$
$`C_1`$
$`C_1`$
$`C_1`$
$`C_1`$
$`C_1`$, nonext.
$`C_1`$, nonext.
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_s`$, nonext. $`(C_{2\nu })`$
$`C_{2\nu }`$
$`C_{2\nu }`$
$`C_{2\nu }`$, nonext.
$`C_{3\nu }`$, nonext.
List of sporadic elementary $`(\{3,4,5\},3)`$-polycycles with $`8`$ faces:
$`C_1`$
$`C_1`$
$`C_1`$
$`C_1`$
$`C_1`$
$`C_1`$
$`C_1`$, nonext.
$`C_1`$, nonext.
$`C_1`$, nonext.
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_2`$
$`C_{2\nu }`$
$`C_{2\nu }`$, nonext.
$`C_{2\nu }`$, nonext.
List of sporadic elementary $`(\{3,4,5\},3)`$-polycycles with $`9`$ faces:
$`C_1`$
$`C_1`$
$`C_1`$, nonext.
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_{3\nu }`$
$`C_{3\nu }`$, nonext.
$`C_{4\nu }`$, nonext.
List of sporadic elementary $`(\{3,4,5\},3)`$-polycycles with $`10`$ faces:
$`C_s`$
$`C_s`$
$`C_s`$, nonext.
$`C_s`$, nonext.
$`C_{2\nu }`$
$`C_{2\nu }`$
$`C_{2\nu }`$, nonext.
$`C_{2\nu }`$, nonext.
$`C_{2\nu }`$, nonext.
List of sporadic elementary $`(\{3,4,5\},3)`$-polycycles with at least $`11`$ faces:
$`C_{5\nu }`$, nonext. $`(I_h)`$
## 3 Classification of elementary $`(\{2,3\},4)`$-polycycles
###### Theorem 4
Any elementary $`(\{2,3\},4)`$-polycycle is one of the following eight:
$`C_{3\nu }`$ $`(D_{3h})`$
$`C_{4\nu }`$
$`C_{2\nu }`$
$`C_{3\nu }`$, nonext. $`(O_h)`$
$`C_{2\nu }`$ $`(D_{2h})`$
$`C_s`$
$`C_{2\nu }`$, nonext. $`(D_{2d})`$
$`C_{2\nu }`$, nonext. $`(D_{3h})`$
Proof. The list of elementary $`(\{3\},4)`$-polycycles is determined in \[DDS05b\] and consists of the first four graphs of this theorem. Let $`P`$ be a $`(\{2,3\},4)`$-polycycle, containing a $`2`$-gon. If $`|F_1|=1`$, then it is the $`2`$-gon. Clearly, the case, in which two $`2`$-gons share one edge, is impossible. Assume that $`P`$ contains two $`2`$-gons, which share a vertex. Then we should add triangle on both sides and so, obtain the second above polycycle. If there is a $`2`$-gon, which does not share a vertex with a $`2`$-gon, then $`P`$ contains the following pattern:
So, clearly, $`P`$ is one of the last two possibilities above. $`\mathrm{}`$
Note that seventh and fourth polycycles in Theorem 4 are, respectively, $`2`$\- and $`3`$-antiprisms; here the exterior face is the unique hole. The $`m`$-antiprism for any $`m2`$ can also be seen as $`(\{2,3\},4)`$-polycycle with $`F_2`$ consisting of the exterior and interior $`m`$-gons; this polycycle is not elementary.
## 4 Classification of elementary $`(\{2,3\},5)`$-polycycles
Let us consider an elementary $`(\{2,3\},5)`$-polycycle $`P`$. Assume that $`P`$ is not an $`i`$-gon and has a $`2`$-gonal face $`f`$. If $`f`$ is adjacent to a hole, then the polycycle is not elementary. So, holes are adjacent only to $`3`$-gons. If one remove such a $`3`$-gon $`t`$, then the third vertex $`v`$ of $`t`$, which is necessarily interior in $`P`$, becomes non-interior in $`Pt`$. The polycycle $`Pt`$ is not necessarily elementary. Let us denote by $`e_1,\mathrm{},e_5`$ the edges incident to $`v`$ and assume that $`e_1`$, $`e_2`$ are edges of $`t`$. The boundary is adjacent only to $`3`$-gons. The potential bridges in $`Pt`$ are $`e_3`$, $`e_4`$ and $`e_5`$. Let us check all five cases:
* If no edge $`e_k`$ is a bridge, then $`Pt`$ is elementary.
* If only $`e_4`$ is a bridge, then two cases can happen:
+ This bridge goes from a hole to the same hole. This means that $`P`$ is formed by the merging of two elementary $`(\{2,3\},5)`$-polycycles.
+ This bridge goes from a hole to another hole. This means that $`P`$ has at least two holes and that $`Pt`$ is formed by the merging of two open edges of an elementary $`(\{2,3\},5)`$-polycycle, which has one hole less.
* If $`e_3`$ or $`e_5`$ is a bridge, then $`Pt`$ is formed by the agglomeration of an elementary $`(\{2,3\},5)`$-polycycle and a $`i`$-gon with $`i=2`$ or $`3`$.
* If $`e_4`$ is a bridge and $`e_3`$ or $`e_5`$ is a bridge, then $`Pt`$ is formed by the agglomeration of an elementary $`(\{2,3\},5)`$-polycycle and two $`i`$-gons with $`i=2`$ or $`3`$.
* If all $`e_k`$ are bridges, then $`P`$ has only one interior vertex.
Given a hole of a $`(R,q)`$-polycycle, its boundary sequence is the sequence of degrees of all consecutive vertices of the boundary of this hole.
###### Theorem 5
The list of elementary $`(\{2,3\},5)`$-polycycles consists of:
1. $`57`$ sporadic $`(\{2,3\},5)`$-polycycles given on Figures 3 and 4,
2. three following infinite $`(\{2,3\},5)`$-polycycles:
$`\alpha `$: $`C_1`$
$`\beta `$: $`C_1`$
$`\gamma `$: $`C_1`$, nonext.
3. six infinite series of $`(\{2,3\},5)`$-polycycles with one hole (they are obtained by concatenating endings of a pair of polycycles, given in (ii); see Figure 2 for the first $`5`$ graphs),
4. the following $`5`$-valent doubly infinite $`(\{2,3\},5)`$-polycycle, called snub $`\mathrm{}`$-antiprism:
pma2, nonext.
5. the infinite series of snub $`m`$-antiprisms<sup>4</sup><sup>4</sup>4Snub $`m`$-antiprism is a $`5`$-valent plane graph defined on page 119 of \[DGS04\]. For $`m=4`$ it is one of well-known regular-faced polyhedra (Johnson polyhedra). For $`m=3`$ it has the skeleton of Icosahedron, but it is distinct from the corresponding sporadic one, having only one hole (exterior $`3`$-gonal face), $`m2`$ (two $`m`$-gonal holes, non-extensible for $`m4`$, with symmetry $`D_{md}`$), represented below for $`m=2,3,4,5`$:
$`D_{2d}`$
$`D_{3d}`$ $`(I_h)`$
$`D_{4d}`$
$`D_{5d}`$
Proof. Let us take an elementary $`(\{2,3\},5)`$-polycycle, which is finite. Then, by removing a triangle, which is adjacent to a boundary, one is led to the situation described above. Hence, the algorithm for enumerating finite elementary $`(\{2,3\},5)`$ polycycles is the following.
1. Begin with isolated $`i`$-gons with $`i=2`$ or $`3`$.
2. For every vertex $`v`$ of an elementary polycycle with $`n`$ interior vertices, consider all possibilities of adding $`2`$\- and $`3`$-gons incident to $`v`$, such that the obtained polycycle is elementary and $`v`$ has become an interior vertex.
3. Reduce by isomorphism.
The above algorithm first finds some sporadic $`(\{2,3\},5)`$-polycycles and the first elements of the infinite series and then find only the elements of the infinite series. In order to prove that this is the complete list of all finite $`(\{2,3\},5)`$-polycycles with only one hole, one needs to consider the case, in which only $`e_4`$ is a bridge going from a hole to the same hole. So, we need to consider all possibilities, where the addition of two elementary $`(\{2,3\},5)`$-polycycles and one $`3`$-gon make a larger elementary $`(\{2,3\},5)`$-polycycle. Given a sequence $`a_1,\mathrm{},a_n`$, we say that a sequence $`b_1,\mathrm{},b_p`$ with $`p<n`$ is a pattern of that sequence if, for some $`n_0`$, one has $`a_{n_0+j1}=b_j`$ or $`a_{n_0+1j}=b_j`$ with the addition being modulo $`n`$. The $`(\{2,3\},5)`$-polycycles, used in that construction, should have the pattern $`3,3,x`$ with $`x4`$ in their boundary sequence. Only the polycycles, which belong to the six infinite series, satisfy this and it is easy to see, that the result of the operation is still one of the six infinite series. So, the list of finite elementary $`(\{2,3\},5)`$-polycycles with one hole is the announced one.
If a polycycle has more than one hole, then it is obtained by the addition of $`2`$\- and $`3`$-gonal faces to a vertex, incident to the boundary (or from another $`(\{2,3\},5)`$-polycycle with a smaller number of holes) by the merging of two open edges in a bridge and addition of a $`3`$-gon. This suggest an iterative procedure, where we begin with a $`(\{2,3\},5)`$-polycycle with one hole and consider all possibilities of extension to a polycycle with two holes and then continue for more complicated polycycles.
The polycycles, which are suitable for such a construction, should contain the pattern $`3,3,x`$ and $`y,3,3`$ in their boundary sequence with $`x4`$ and $`y3`$. Amongst all finite $`(\{2,3\},5)`$-polycycles with one hole, only members of the infinite series $`\alpha \alpha `$ of $`(\{2,3\},5)`$-polycycles with one hole, are suitable ones and one obtains the infinite series of the above theorem. Furthermore, the members of this infinite series admit no extension. This proves that there are no other $`(\{2,3\},5)`$-polycycles with at least two holes.
Consider now an elementary infinite $`(\{2,3\},5)`$-polycycle $`P`$. Eliminate all $`2`$-gonal faces of $`P`$ and obtain another $`(\{3\},5)`$-polycycle $`P^{}`$, which is not necessarily elementary. We do a decomposition of $`P^{}`$ along its elementary components, which are enumerated in \[DDS05b\]. If snub $`\mathrm{}`$-antiprism is one of the components, then we are finished and $`P=P^{}`$ is the snub $`\mathrm{}`$-antiprism. If $`\alpha `$ is one of the components, then one has two edges, along which to extend the polycycle; they are depicted below:
Clearly, if we extend the polycycle along only one of those edges, then the result is not an elementary polycycle. The consideration of all possibilities yields $`\beta `$ and $`\gamma `$. Suppose now that $`P^{}`$ has no infinite components. Then $`P`$ has at least one infinite path $`f_0,\mathrm{},f_i,\mathrm{}`$, such that $`f_i`$ is adjacent to $`f_{i+1}`$, but $`f_{i1}`$ is not adjacent to $`f_{i+1}`$. The considerations, analogous to the $`3`$-valent case, yield the result. $`\mathrm{}`$
## 5 Possible extensions of the setting
The classification, developped in this paper, is an extension of previous work of Deza and Shtogrin, by allowing various possible number of sides of interior faces and several possible holes. A natural question is, if one can further enlarge the class of polycycles.
We required $`2`$-connectivity and that any two holes do not share a vertex. If one removes those two hypothesis, then many other graphs do appear.
Consider infinite series of $`(\{2\},6)`$-polycycles, m-bracelets, $`m2`$ (i.e. $`m`$-circle, but each edge is tripled). The central edge is a bridge for those polycycles, for both $`2`$-gons of the triple of edges. But if one removes those two digons, then the resulting plane graph has two holes sharing a face, i.e. violates the second of the crucial points (i)–(iii) of the definition of $`(R,q)`$-polycycle. This shows that our hypothesis were necessary.
For even $`m`$, each even edge (for some order $`1,\mathrm{},m`$ of them) can be duplicated $`t`$ times (for fixed $`t`$, $`1t5`$), and each odd edge duplicated $`6t`$ times; so, all vertex-valencies will be still $`6`$. On the other hand, two holes ($`m`$-gons inside and outside of the $`m`$-bracelet), have common vertices; so, it is again not our polycycle. |
warning/0507/math0507404.html | ar5iv | text | # A short course on geometric motivic integration
## 1. The invention of motivic integration.
Motivic integration was introduced by Kontsevich to prove the following result conjectured by Batyrev: Let
be two crepant resolutions of the singularities of $`X`$, which itself is a complex projective Calabi-Yau<sup>1</sup><sup>1</sup>1Usually, a normal projective variety $`X`$ of dimension $`n`$ is called Calabi-Yau if the canonical divisor $`K_X`$ is trivial and $`H^i(X,𝒪_𝒳)=\mathcal{0}`$ for $`0<i<n`$. This last condition on the cohomology vanishing is not necessary for the statements below. In the context of mirror symmetry it seems customary to drop this last condition and call $`X`$ Calabi-Yau as soon as $`K_X=0`$ (and the singularities are mild), see . variety with at worst canonical Gorenstein singularities. Crepant (as in *non discrepant*) means that the pullback of the canonical divisor class on $`X`$ is the canonical divisor class on $`X_i`$, i.e. the discrepancy divisor $`E_i=K_{X_i}\pi _i^{}K_X`$ is numerically equivalent to zero. In this situation Batyrev showed, using $`p`$-adic integration, that $`X_1`$ and $`X_2`$ have the same betti numbers $`h_i=dimH^i(\underset{¯}{},)`$. This lead Kontsevich to invent *motivic integration* to show that $`X_1`$ and $`X_2`$ even have the same Hodge numbers $`h^{i,j}=dimH^i(\underset{¯}{},\mathrm{\Omega }^j)`$.
This problem was motivated by the *topological mirror symmetry test* of string theory which asserts that if $`X`$ and $`X^{}`$ are a mirror pair<sup>2</sup><sup>2</sup>2To explain what a mirror pair is in a useful manner lies beyond my abilities. For our purpose one can think of a mirror pair (somewhat tautologically) as being a pair that passes the topological mirror symmetry test. Another achievement of Batyrev was to explicitly construct the mirror to a mildly singular (toric) Calabi-Yau variety. of smooth Calabi-Yau varieties then they have mirrored Hodge numbers
$$h^{i,j}(X)=h^{ni,j}(X^{}).$$
As the mirror of a smooth Calabi-Yau might be singular, one cannot restrict to the smooth case and the equality of Hodge numbers actually fails in this case. Therefore Batyrev suggested, inspired by string theory, that one should look instead at the Hodge numbers of a crepant resolution, if such exists<sup>3</sup><sup>3</sup>3 Calabi-Yau varieties do not always have crepant resolutions. I think one of Batyrev’s papers discusses this.. The independence of these numbers from the chosen crepant resolution is Kontsevich’s result. This makes the *stringy Hodge numbers* $`h_{st}^{i,j}(X)`$ of $`X`$, defined as $`h^{i,j}(X^{})`$ for a crepant resolution $`X^{}`$ of $`X`$, well defined. This leads to a modified mirror symmetry conjecture, asserting that the stringy Hodge numbers of a mirror pair are equal .
Batyrev’s conjecture is now Kontsevich’s theorem and the simplest form to phrase it might be:
###### Theorem 1.1 (Kontsevich).
Birationally equivalent smooth Calabi-Yau varieties have the same Hodge numbers.<sup>4</sup><sup>4</sup>4There is now a proof by Ito of this result using $`p`$-adic integration, thus continuing the ideas of Batyrev who proved the result for Betti numbers using this technique. Furthermore the recent weak factorization theorem of Włodarczyk allows for a proof avoiding integration of any sort.
###### Proof.
The idea now is to assign to any variety a *volume* in a suitable ring $`\widehat{}_k`$ such that the information about the Hodge numbers is retained. The following diagram illustrates the construction of $`\widehat{}_k`$:
The diagonal map is the (compactly supported) Hodge characteristic, which on a smooth projective variety $`X`$ is given by $`E(X)=(1)^idimH^i(X,\mathrm{\Omega }_X^j)u^iv^j`$. In general it is defined via mixed Hodge structures<sup>5</sup><sup>5</sup>5Recently, Bittner gave an alternative construction of the compactly supported Hodge characteristic. She uses the weak factorization theorem of Włodarczyk to reduce the definition of $`E`$ to the case of $`X`$ smooth and projective, where it is as given above. , satisfies $`E(X\times Y)=E(X)E(Y)`$ for all varieties $`X,Y`$ and has the property that for $`YX`$ a closed $`k`$-subvariety one has $`E(X)=E(Y)+E(XY)`$. Therefore the Hodge characteristic factors through the *naive Grothendieck ring* $`K_0(Var_k)`$ which is the universal object with the latter property.<sup>6</sup><sup>6</sup>6$`K_0(\mathrm{Var}_k)`$ is the free abelian group on the isomorphism classes $`[X]`$ of $`k`$-varieties subject to the relations $`[X]=[XY]+[Y]`$ for $`Y`$ a closed subvariety of $`X`$. The product is given by $`[X][Y]=[X\times _kY]`$. The symbol $`𝕃`$ denotes the class of the affine line $`[𝔸_k^1]`$. This explains the left triangle of the diagram.
The bottom row of the diagram is the composition of a localization (inverting $`uv`$) and a completion with respect to negative degree. $`_𝓀`$ is constructed analogously, by first inverting $`𝕃^1=[𝔸_k^1]`$ (a pre-image of $`uv`$) and then completing appropriately (negative dimension). Whereas the bottom maps are injective (easy exercise), the map $`K_0(\mathrm{Var}_k)\stackrel{}{}\widehat{}_k`$ is most likely not injective. The need to work with $`\widehat{}_k`$ instead of $`K_0(\mathrm{Var}_k)`$ arises in the setup of the integration theory an will become clear later.<sup>7</sup><sup>7</sup>7In fact, recent results of F. Loeser and R. Cluckers , and J. Sebag indicate that the full completion may not be necessary, and all the volumes of measurable sets are contained in a subring of $`\widehat{}_k`$ that can be constructed explicitly.
Clearly, by construction it is now enough to show that birationally equivalent Calabi-Yau varieties have the same *volume*, i.e. the same class in $`\widehat{}_k`$. This is achieved via the all important *birational transformation rule* of motivic integration. Roughly it asserts that for a proper birational map $`\pi :Y\stackrel{}{}X`$ the class $`[X]\widehat{}_k`$ is an *expression* in $`Y`$ and $`K_{Y/X}`$ only:
$$[X]=_Y𝕃^{\mathrm{ord}_{K_{Y/X}}}$$
To finish off the proof let $`X_1`$ and $`X_2`$ be birationally equivalent Calabi-Yau varieties. We resolve the birational map to a Hironaka hut:
By the Calabi-Yau assumption we have $`K_{X_i}0`$ and therefore $`K_{Y/X_i}K_Y\pi _i^{}K_{X_i}K_Y`$. Hence the divisors $`K_{Y/X_1}`$ and $`K_{Y/X_2}`$ are numerically equivalent. This numerical equivalence implies in fact an equality of divisors $`K_{X/X_1}=K_{X/X_2}`$ since, again by the Calabi-Yau assumption, $`dimH^0(X,K_Y)=dimH^0(X_i,𝒪_{𝒳_𝒾})=\mathcal{1}`$.<sup>8</sup><sup>8</sup>8In general, the condition that $`X_1`$ and $`X_2`$ have a common resolution $`Y`$ such that $`K_{Y/X_1}`$ is numerically equivalent $`K_{Y/X_1}`$ is called $`K`$–equivalence. We showed above that two birational Calabi–Yau varieties are $`K`$–equivalent. For mildly singular $`X_i`$ (say canonical) it can be derived from the negativity lemma \[30, Lemma 3.39\] that $`K`$–equivalence implies actual equality of divisors $`K_{Y/X_1}`$ and $`K_{Y/X_1}`$. Hence the Calabi–Yau assumption was not essential to conclude this (but provides a simple argument). By the transformation rule, $`[X_1]`$ is an expression depending only on $`Y`$ and $`K_{X/X_1}=K_{X/X_2}`$. The same is true for $`[X_2]`$ and thus we have $`[X_1]=[X_2]`$ as desired. ∎
These notes were started during a working seminar at MSRI during the year of 2003 and took shape in the course of the past 2 years while I was giving introductory lectures on the subject. They have taken me way too much time to finish and would not have been finished at all if it weren’t for the encouragement of many people: Thanks goes to all the participants of the seminar on motivic integration at MSRI (2002/2003), of the Schwerpunkt Junioren Tagung in Bayreuth (2003) and the patient listeners of the mini-courses at KTH, Stockholm (2003), the University of Helsinki (2004) and the Vigre graduate course in Salt Lake City (2005). Special thanks goes to Karen Smith for encouragement to start this project and to Julia Gordon for numerous comments, suggestions and careful reading.
## 2. Geometric motivic integration
We now assume that $`k`$ is algebraically closed and of characteristic zero. In fact, there is one point (see section 4.1) where we will assume that $`k=`$ in order to avoid some technicalities which arise if the field is not uncountable. Thus the reader may choose to replace $`k`$ by $``$ whenever it is comforting. We stress that there are significant (though manageable) obstacles one has to overcome if one wants to (a) work with singular spaces or (b) with varieties defined over fields which are not uncountable or not algebraically closed. Or put differently: The theory develops naturally (for an algebraic geometer), and easily, in the smooth case over $``$, as we hope to demonstrate below. In order to transfer this intuition to any other situation, nontivial results and extra care is necessary.
All the results in these notes appeared in the papers of Denef and Loeser, Batyrev and Looijenga. Our exposition is particularly influenced by Looijenga and Batyrev . Also Craw was very helpful as a first reading. We also recommend the articles of Hales and Veys , both explain the connection to $`p`$-adic integration in detail, which we do not discuss in these notes at all. The above mentioned references are also a great source to learn about the various different applications the theory had to date. We will discuss none of them except for certain applications to birational geometry.
We will now introduce the building blocks of the theory. These are:
1. The value ring of the measure: *$`\widehat{}_k`$, a localized and completed Grothendieck ring.*
2. A domain of integration: *$`𝒥_{\mathrm{}}(𝒳)`$, the space of formal arcs over $`X`$.*
3. An algebra of measurable sets of $`𝒥_{\mathrm{}}(𝒳)`$ and a measure: *cylinders/stable sets and the virtual euler characteristic.*
4. An interesting class of measurable/integrable functions: *Contact order of an arc along a divisor.*
5. A change of variables formula: *Kontsevich’s birational transformation rule.*
These basic ingredients appear with variations in all versions of motivic integration (one could argue: of any theory of integration).
### 2.1. The value ring of the motivic measure
Here we already gravely depart from any previous (classical?) theory of integration since the values of our measure do not lie in $``$. Instead they lie in a huge ring, constructed from the Grothendieck ring of varieties by a process of localization and completion. This ingenious choice is a key feature of the theory.
We start with the naive Grothendieck ring of the category of varieties over $`k`$.<sup>9</sup><sup>9</sup>9Alternatively, the Grothendieck ring of finite type schemes over $`k`$ leads to the same object because $`XX_{\mathrm{red}}=\mathrm{}`$. As Bjorn Poonen points out, the finite type assumption is crucial here. If one would allow non finite type schemes, $`K_0(\mathrm{Var}_k)`$ would be zero. For this let $`Y`$ be any $`k`$–scheme and let $`X`$ be an infinite disjoint union of copies of $`Y`$. Then $`[X]+[Y]=[X]`$ and therefore $`[Y]=0`$. This is the ring $`K_0(\mathrm{Var}_k)`$ generated by the isomorphism classes of all finite type $`k`$–varieties and with relation $`[X]=[Y]+[XY]`$ for a closed $`k`$-subvariety $`YX`$, that is such that the inclusion $`YX`$ is defined over $`k`$. The square brackets denote the image of $`X`$ in $`K_0(\mathrm{Var}/k)`$. The product structure is given by the fiber product, $`[X][Y]=[X\times _kY](=[(X\times _kY)_{\mathrm{red}}])`$. The symbol $`𝕃`$ is reserved for the class of the affine line $`[𝔸_k^1]`$ and $`1=1_k`$ denotes $`\mathrm{Spec}k`$. Thus, for example, $`[^n]=𝕃^n+𝕃^{n1}+\mathrm{}+𝕃+1`$.
Roughly speaking the map $`X[X]`$ is robust with respect to chopping up $`X`$ into a disjoint union of locally closed subvarieties.<sup>10</sup><sup>10</sup>10This is elegantly illustrated in the article of Hales which emphasizes precisely this point of $`K_0(\mathrm{Var}_k)`$ being a *scissor group*. By using a stratification of $`X`$ by smooth subvarieties, this shows that $`K_0(\mathrm{Var}_k)`$ is generated by the classes of smooth varieties.<sup>11</sup><sup>11</sup>11In , Bittner shows that $`K_0(\mathrm{Var}_k)`$ is the abelian group generated by smooth projective varieties subject to a class of relations which arise from blowing up at a smooth center: If $`Z`$ is a smooth subvariety of $`X`$, then the relation is $`[X][Z]=[\mathrm{Bl}_ZX][E]`$, where $`E`$ is the exceptional divisor of the blowup. In a similar fashion one can assign to every constructible subset $`C`$ of $`X`$ a class $`[C]`$ by expressing $`C`$ as a combination of subvarieties.
###### Exercise 2.1.
Verify the claim in the last sentence. That is: show that the map $`Y[Y]`$ for $`Y`$ a closed subvariety of $`X`$ naturally extends to the algebra of constructible subsets of $`X`$.
###### Exercise 2.2.
Let $`Y\stackrel{}{}X`$ be a piecewise trivial fibration with constant fiber $`Z`$. This means one can write $`X=X_i`$ as a finite disjoint union of locally closed subsets $`X_i`$ such that over each $`X_i`$ one has $`f^1X_iX_i\times Z`$ and $`f`$ is given by the projection onto $`X_i`$. Show that in $`K_0(\mathrm{Var}_k)`$ one has $`[Y]=[X][Z]`$.
There is a natural notion of dimension of an element of $`K_0(\mathrm{Var}_k)`$. We say that $`\tau K_0(\mathrm{Var}_k)`$ is *$`d`$–dimensional* if there is an expression in $`K_0(\mathrm{Var}_k)`$
$$\tau =a_i[X_i]$$
with $`a_i`$ and $`k`$-varieties $`X_i`$ of dimension $`d`$, and if there is no expression like this with all $`dimX_id1`$. The dimension of the class of the empty variety is set to be $`\mathrm{}`$. It is easy to verify (exercise!) that the map
$$dim:K_0(\mathrm{Var}_k)\stackrel{}{}\{\mathrm{}\}$$
satisfies $`dim(\tau \tau ^{})dim\tau +dim\tau ^{}`$ and $`dim(\tau +\tau ^{})\mathrm{max}\{dim\tau ,dim\tau ^{}\}`$ with equality in the latter if $`dim\tau dim\tau ^{}`$.
The dimension can be extended to the localization $`_𝓀\stackrel{\mathrm{def}}{=}𝒦_\mathcal{0}(\mathrm{Var}_𝓀)[𝕃^\mathcal{1}]`$ simply by demanding that $`𝕃^1`$ has dimension $`1`$. To obtain the ring $`\widehat{}_k`$ in which the desired measure will take values we further complete $`_𝓀`$ with respect to the filtration induced by the dimension.<sup>12</sup><sup>12</sup>12In Looijenga , this is called the virtual dimension. As described by Batyrev, composing the dimension $`dim:_𝓀\stackrel{}{}\{\mathrm{}\}`$ with the exponential $`\stackrel{\mathrm{exp}(\underset{¯}{})}{}_+`$ and by further defining $`\mathrm{}0`$ we get a map
$$\delta _k:_𝓀\stackrel{}{}_{+,\mathcal{0}}$$
which is a *non-archimedian norm*. That means the following properties hold: $`\delta _k(A)=0`$ iff $`A=0=[\mathrm{}]`$ in $`_𝓀`$. $`\delta _k(A+B)\mathrm{max}\{\delta _k(A),\delta _k(B)\}`$ $`\delta _k(AB)\delta _k(A)\delta _k(B)`$ The ring $`\widehat{}_k`$ is then the completion with respect to this norm, and therefore $`\widehat{}_k`$ is complete in the sense that all Cauchy sequences uniquely converge. The condition (2) is stronger than the one used in the definition of an archimedian norm. This non-archimedian ingredient makes the notion of convergence of sums conveniently simple; a sum converges if and only if the sequence of summands converges to zero.
If there was an equality in condition (3) the norm would be called *multiplicative*. It is unknown whether $`\delta `$ is multiplicative on $`_𝓀`$. However, Poonen shows that $`K_0(\mathrm{Var}_k)`$ contains zero divisors, thus $`\delta `$ restricted to $`K_0(\mathrm{Var}_k)`$ is *not* multiplicative on $`K_0(\mathrm{Var}_k)`$. The $`n`$th filtered subgroup is
$$^𝓃(_𝓀)=\{\tau _𝓀|dim\tau 𝓃\}.$$
This gives us the following maps which will be the basis for constructing the sought after motivic measure:
$$\mathrm{Var}_k\stackrel{}{}K_0(\mathrm{Var}_k)\stackrel{\text{invert }𝕃}{}_𝓀\overline{}\widehat{}_𝓀.$$
We will somewhat ambiguously denote the image of $`X\mathrm{Var}_k`$ in any of the rings to the right by $`[X]`$. It is important to point out here that it is unknown whether the completion map $`\underset{¯}{}^{}`$ is injective, i.e. whether its kernel, $`^𝒹(𝒦_\mathcal{0}(\mathrm{Var}/𝓀)[𝕃^\mathcal{1}])`$, is zero. It is also unknown whether the localization is injective. <sup>13</sup><sup>13</sup>13In Poonen shows that $`K_0(\mathrm{Var}/k)`$ is not a domain in characteristic zero. It is expected though that the localization map is not injective and that $`_𝓀`$ is a domain and that the completion map $`_𝓀\stackrel{}{}\widehat{}_𝓀`$ is injective. But recently Naumann found in his dissertation zero-divisors in $`K_0(\mathrm{Var}_k)`$ for $`k`$ a finite field and these are non-zero even after localizing at $`𝕃`$ – thus for a finite field $`_𝓀`$ is not a domain. For infinite fields (e.g. $`k`$ algebraically closed) the above questions remain open.
###### Exercise 2.3.
Convergence of series in $`\widehat{}_k`$ is rather easy. For this observe that a sequence of elements $`\tau _i_𝓀`$ converges to zero in $`\widehat{}_k`$ if and only if the dimensions $`dim\tau _i`$ tend to $`\mathrm{}`$ as $`i`$ approaches $`\mathrm{}`$. Show that a sum $`_{i=0}^{\mathrm{}}\tau _i`$ converges if and only the sequence of summands converges to zero.
###### Exercise 2.4.
Show that in $`\widehat{}_k`$ the equality $`_{i=0}^{\mathrm{}}𝕃^{ki}=\frac{1}{1𝕃^k}`$ holds.
### 2.2. The arc space $`𝒥_{\mathrm{}}(X)`$
Arc spaces were first studied seriously by Nash who conjectured a tight relationship between the geometry of the arc space and the singularities of $`X`$, see Ishii and Kollar for a recent exposition of Nash’s ideas in modern language. Recent work of Mustaţǎ supports these predictions by showing that the arc spaces contain information about singularities, for example rational singularities of $`X`$ can be detected by the irreducibility of the jet schemes for complete intersections. In subsequent investigations he and his collaborators show that certain invariants of birational geometry, such as the log canonical threshold of a pair, for example, can be read off from the dimensions of certain components of the jet schemes, see and Section 5 where we will discuss some of these results in detail.
Let $`X`$ be a (smooth) scheme of finite type over $`k`$ of dimension $`n`$. An $`m`$-jet of $`X`$ is an order $`m`$ infinitesimal curve in $`X`$, i.e. it is a morphism
$$\vartheta :\mathrm{Spec}k[t]/t^{m+1}\stackrel{}{}X.$$
The set of all $`m`$-jets carries the structure of a scheme $`𝒥_m(X)`$, called the $`m`$th *jet scheme*, or space of truncated arcs. It’s characterizing property is that it is right adjoint to the functor $`\underset{¯}{}\times \mathrm{Spec}k[t]/t^{m+1}`$. In other words,
$$\mathrm{Hom}(Z\times \mathrm{Spec}k[t]/t^{m+1},X)=\mathrm{Hom}(Z,𝒥_m(X))$$
for all $`k`$-schemes $`Z`$, i.e. $`𝒥_m(X)`$ is the scheme which represents the contravariant functor $`\mathrm{Hom}(\underset{¯}{}\times \mathrm{Spec}k[t]/t^{m+1},X)`$.<sup>14</sup><sup>14</sup>14Representability of this functor was proved by Greenberg ; another reference for this fact is . In particular this means that the $`k`$–valued points of $`𝒥_m(X)`$ are precisely the $`k[t]/t^{m+1}`$–valued points of $`X`$. The so called Weil restriction of scalars, i.e. the natural map $`k[t]/t^{m+1}\stackrel{}{}k[t]/t^m`$, induces a map $`\pi _{m1}^m:𝒥_m(X)\stackrel{}{}𝒥_{m1}(X)`$ and composition gives a map $`\pi ^m:𝒥_m(X)\stackrel{}{}𝒥_0(X)=X`$. As upper indices are often cumbersome we define $`\eta _m=\pi ^m`$ and $`\phi _m=\pi _{m1}^m`$.
Taking the inverse limit<sup>15</sup><sup>15</sup>15For this to be defined one crucially uses that the restriction maps are *affine* morphisms. of the resulting system yields the definition of the *infinite jet scheme*, or the *arc space*
$$𝒥_{\mathrm{}}(X)=\underset{}{\mathrm{lim}}𝒥_m(X).$$
Its $`k`$-points are the limit of the $`k`$-valued points $`\mathrm{Hom}(\mathrm{Spec}k[t]/t^{m+1},X)`$ of the jet spaces $`𝒥_m(X)`$. Therefore they correspond to the formal curves (or arcs) in $`X`$, that is to maps $`\mathrm{Spec}kt\stackrel{}{}X`$.<sup>16</sup><sup>16</sup>16For this observe that $`\underset{}{\mathrm{lim}}\mathrm{Hom}(R,k[t]/t^{m+1})\mathrm{Hom}(R,\underset{}{\mathrm{lim}}k[t]/t^{m+1})=\mathrm{Hom}(R,kt)`$. There are also maps $`\pi _m:𝒥_{\mathrm{}}(X)\stackrel{}{}𝒥_m(X)`$ again induced by the truncation map $`kt\stackrel{}{}kt/t^{m+1}`$. If there is danger of confusion we sometimes decorate the projections $`\pi `$ with the space. The following picture should help to remember the notation.
(1)
are the maps induced by the natural surjections
$$\text{}.$$
###### Example 2.1.
Let $`X=\mathrm{Spec}k[x_1,\mathrm{},x_n]=𝔸^n`$. Then, on the level of $`k`$-points one has
$$𝒥_m(X)=\{\vartheta :k[x_1,\mathrm{},x_n]\stackrel{}{}kt/t^{m+1}\}$$
Such a map $`\vartheta `$ is determined by its values on the $`x_i`$’s, i.e. it is determined by the coefficients of $`\vartheta (x_i)=_{j=0}^m\vartheta _i^{(j)}t^j`$. Conversely, any choice of coefficients $`\vartheta _i^{(j)}`$ determines a point in $`𝒥_m(𝔸^n)`$. Choosing coordinates $`x_i^{(j)}`$ of $`𝒥_m(X)`$ with $`x_i^{(j)}(\vartheta )=\vartheta _i^{(j)}`$ we see that
$$𝒥_m(X)\mathrm{Spec}k[x_1^{(0)},\mathrm{},x_n^{(0)},\mathrm{}\mathrm{},x_1^{(m)},\mathrm{},x_n^{(m)}]𝔸^{n(m+1)}.$$
Furthermore observe that, somewhat intuitively, the truncation map $`\pi ^m:𝒥_m(X)\stackrel{}{}X`$ is induced by the inclusion
$$k[x_1,\mathrm{},x_n]k[x_1^{(0)},\mathrm{},x_n^{(0)},\mathrm{}\mathrm{},x_1^{(m)},\mathrm{},x_n^{(m)}]$$
sending $`x_i`$ to $`x_i^{(0)}`$.
###### Exercise 2.5.
Let $`Y𝔸^n`$ be a hypersurface given by the vanishing of one equation $`f=0`$. Show that $`𝒥_m(Y)𝒥_m(𝔸^n)`$ is given by the vanishing of $`m+1`$ equations $`f^{(0)},\mathrm{},f^{(m)}`$ in the coordinates of $`𝒥_m(𝔸^n)`$ described above. (Observe that $`f^{(0)}=f(x^{(0)})`$ and $`f^{(1)}=\frac{}{x_i}f(x^{(0)})x_i^{(1)}`$). Show that
1. $`𝒥_m(Y)`$ is pure dimensional if and only if $`dim𝒥_m(Y)=(m+1)(n1)`$, in which case $`𝒥_m(Y)`$ is a complete intersection.
2. $`𝒥_m(Y)`$ is irreducible if and only if $`dim(\pi _Y^m)^1(Y_{\mathrm{Sing}})<(m+1)(n1)`$.
Similar statements hold if $`Y`$ is locally a complete intersection.
The existence of the jet schemes in general (that is to show the representability of the functor defined above) is proved, for example, in . From the very definition one can easily derive the following étale invariance of jet schemes, which, together with the example of $`𝔸^n`$ above gives us a pretty good understanding of the jet schemes of a smooth variety.
###### Proposition 2.2.
Let $`X\stackrel{}{}Y`$ be étale, then $`𝒥_m(X)𝒥_m(Y)\times _YX`$.
###### Proof.
We show the equality on the level of the corresponding functors of points
$$\mathrm{Hom}(\underset{¯}{},𝒥_m(X))\mathrm{Hom}(\underset{¯}{}\times _k\mathrm{Spec}\frac{kt}{t^{m+1}},X)$$
and
$$\begin{array}{cc}\hfill \mathrm{Hom}(\underset{¯}{},𝒥_m(Y)\times _YX)& =\mathrm{Hom}(\underset{¯}{},𝒥_m(Y))\times \mathrm{Hom}(\underset{¯}{},X)\hfill \\ & =\mathrm{Hom}(\underset{¯}{}\times _k\mathrm{Spec}\frac{kt}{t^{m+1}},Y)\times \mathrm{Hom}(\underset{¯}{},X).\hfill \end{array}$$
For this let $`Z`$ be a $`k`$–scheme and consider the diagram
to see that a $`Z`$-valued $`m`$-jet $`\tau \mathrm{Hom}(Z\times _k\mathrm{Spec}\frac{kt}{t^{m+1}},X)`$ of $`X`$ induces a $`Z`$-valued $`m`$-jet $`\vartheta \mathrm{Hom}(Z\times _k\mathrm{Spec}\frac{kt}{t^{m+1}},Y)`$ and a map $`p\mathrm{Hom}(Z,X)`$. Virtually by definition formally étaleness \[23, Definition (17.1.1)\] for the map from $`X`$ to $`Y`$, the converse holds also, i.e. $`\vartheta `$ and $`p`$ together induce a unique map $`\tau `$ as indicated. ∎
Using this étale invariance of jet schemes the computation carried out for $`𝔸^n`$ above holds locally on any smooth $`X`$. Thus we obtain:
###### Proposition 2.3.
Let $`X`$ be a smooth $`k`$-scheme of dimension $`n`$. Then $`𝒥_m(X)`$ is locally an $`𝔸^{nm}`$–bundle over $`X`$. In particular $`𝒥_m(X)`$ is smooth of dimension $`n(m+1)`$. In the same way, $`𝒥_{m+1}(X)`$ is locally an $`𝔸^n`$–bundle over $`𝒥_m(X)`$.
Note that this is not true for a singular $`X`$ as can be seen already by looking at the tangent bundle $`TX=𝒥_1(X)`$ which is well-known to be a bundle if and only if $`X`$ is smooth. In fact, over a singular $`X`$ the jet schemes need not even be irreducible nor reduced and can also be badly singular.
### 2.3. An algebra of measurable sets
The prototype of a measurable subset of $`𝒥_{\mathrm{}}(X)`$ is a *stable set*. They are defined just right so that they receive a natural volume in $`_𝓀`$.
###### Definition 2.4.
A subset $`A𝒥_{\mathrm{}}(X)`$ is called *stable* if for all $`m0`$, $`A_m\stackrel{\mathrm{def}}{=}\pi _m(A)`$ is a constructible subset<sup>17</sup><sup>17</sup>17The constructible subsets of a scheme $`Y`$ are the smallest algebra of sets containing the closed sets in Zariski topology. of $`𝒥_𝓂(𝒳)`$, $`A=\pi _m^1(A_m)`$ and the map
(2)
$$\pi _m^{m+1}:A_{m+1}\stackrel{}{}A_m\text{is a locally trivial }𝔸^n\text{–bundle.}$$
For any $`m0`$ we define the *volume* of the stable set $`A`$ by
$$\mu _X(A)=[A_m]𝕃^{nm}_𝓀.$$
That this is independent of $`m`$ is ensured by condition (2) which implies that $`[A_{m+1}]=[A_m]𝕃^n`$.<sup>18</sup><sup>18</sup>18Reid , Batyrev and Looijenga use this definition which gives the volume $`\mu _X(𝒥_m(X))_𝓀`$ of $`X`$ virtual dimension $`n`$. Denef, Loeser and Craw use an additional factor $`𝕃^n`$ to give $`\mu _X(𝒥_m(X))`$ virtual dimension $`0`$. It seems to be essentially a matter of taste which definition one uses. Just keep this in mind while browsing through different sources in the literature to avoid unnecessary confusion.
Assuming that $`X`$ is smooth one uses Proposition 2.3 to show that the collection of stable sets forms an algebra of sets, which means that $`𝒥_{\mathrm{}}(X)`$ is stable and with $`A`$ and $`A^{}`$ stable the sets $`𝒥_{\mathrm{}}(X)A`$ and $`AA^{}`$ are also stable. The smoothness of $`X`$ furthermore warrants that so called *cylinder sets* are stable (a cylinder is a set $`A=\pi _m^1B`$ for some constructible $`B𝒥_𝓂(𝒳)`$). Thus in the smooth case condition (2) is superfluous whereas in general it is absolutely crucial. In fact, a main technical point in defining the motivic measure on singular varieties is to show that the class of stable sets can be enlarged to an algebra of *measurable* sets which contains the cylinders. In particular $`𝒥_{\mathrm{}}(X)`$ is then measurable. This is achieved as one would expect by declaring a set measurable if it is approximated in a suitable sense by stable sets. This is essentially carried out in , though there are some inaccuracies; but everything should be fine if one works over $``$ and makes some adjustments following \[2, Appendix\].<sup>19</sup><sup>19</sup>19Of course Denef and Loeser also set up motivic integration over singular spaces but their approach differs from the one discussed here in the sense that they assign a volume to the *formula* defining a constructible set rather than to the set of ($`k`$-rational) points itself. Thus they elegantly avoid any problems which arise if $`k`$ is small. To avoid these technicalities we assume until the end of this section that $`X`$ is smooth over the complex numbers $``$.
### 2.4. The measurable function associated to a subscheme
From an algebra of measurable sets there arises naturally a notion of measurable function. Since we did not carefully define the measurable sets — we merely described the prototypes — we will for now only discuss an important class of measurable functions.
Let $`YX`$ be a subscheme of $`X`$ defined by the sheaf of ideals $`I_Y`$. To $`Y`$ one associates the function
$$\mathrm{ord}_Y:𝒥_{\mathrm{}}(X)\stackrel{}{}\{\mathrm{}\}$$
sending an arc $`\vartheta :𝒪_𝒳\stackrel{}{}𝓀𝓉`$ to the order of vanishing of $`\vartheta `$ along $`Y`$, i.e. to the supremum of all $`e`$ such that ideal $`\vartheta (I_Y)`$ of $`kt`$ is contained in the ideal $`(t^e)`$. Equivalently, $`\mathrm{ord}_Y(\vartheta )`$ is the supremum of all $`e`$ such that the map
$$𝒪_𝒳\stackrel{\mathit{\vartheta }}{}𝓀𝓉\stackrel{}{}𝓀𝓉/𝓉^{}$$
sends $`I_Y`$ to zero. Note that this map is nothing but the truncation $`\pi _{e1}(\vartheta )𝒥_{e1}(X)`$ of $`\vartheta `$.<sup>20</sup><sup>20</sup>20At this point we better set $`𝒥_1(X)=X`$ and $`\pi _1=\pi _0=\pi `$ to avoid dealing with the case $`e=0`$ separately. For a $`(e1)`$-jet $`\gamma 𝒥_{e1}(X)`$ to send $`I_Y`$ to zero means precisely that $`\gamma 𝒥_{e1}(Y)`$.
Thus we can rephrase this by saying that $`\mathrm{ord}_Y(\vartheta )`$ is the supremum of all $`e`$ such that the truncation $`\pi _{e1}(\vartheta )`$ lies in $`𝒥_{e1}(Y)`$. Now it is clear that
$`\mathrm{ord}_Y(\vartheta )0`$ $``$ $`\pi (\vartheta )Y,`$
$`\mathrm{ord}_Y(\vartheta )s`$ $``$ $`\pi _{s1}(\vartheta )𝒥_{s1}(Y)\text{ and}`$
$`\mathrm{ord}_Y(\vartheta )=\mathrm{}`$ $``$ $`\vartheta 𝒥_{\mathrm{}}(Y).`$
The functions $`\mathrm{ord}_Y`$ just introduced are examples of measurable functions which come up in the applications of motivic integration. For a function to be measurable, one requires that the level sets $`\mathrm{ord}_Y^1(s)`$ are measurable sets, i.e. stable sets (or at least suitably approximated by stable sets). For this consider the set $`\mathrm{ord}_Y^1(s)=\{\vartheta 𝒥_{\mathrm{}}(X)|\mathrm{ord}_Y(\vartheta )s\}`$ consisting of all arcs in $`X`$ which vanish of order *at least* $`s`$ along $`Y`$. By what we just observed $`\mathrm{ord}_Y^1(s)=\pi _{s1}^1(𝒥_{s1}(Y))`$ is a cylinder. Therefore, the level set $`\mathrm{ord}_Y^1(s)`$ is also a cylinder equal to
$$\mathrm{ord}_Y^1(s)\mathrm{ord}_Y^1(s+1)=\pi _{s1}^1𝒥_{s1}(Y)\pi _s^1𝒥_s(Y).$$
The exception is $`s=0`$ in which case $`\mathrm{ord}_Y^1(0)=\pi ^1(X)=𝒥_{\mathrm{}}(X)`$ and $`\mathrm{ord}_Y^1(0)=\pi ^1(X)\pi ^1(Y)`$.
Note that the level set at infinity,
$$\mathrm{ord}_Y^1(\mathrm{})=\mathrm{ord}_Y^1(s)=𝒥_{\mathrm{}}(Y)$$
on the other hand is *not* a cylinder.<sup>21</sup><sup>21</sup>21This just says that for an arc $`\vartheta 𝒥_{\mathrm{}}(X)`$ to lie in $`𝒥_{\mathrm{}}(Y)`$ is a condition that cannot be checked on any truncation. Exercise! Still, since it is the decreasing intersection of the cylinders $`\mathrm{ord}_Y^1(s+1)=\pi _s^1𝒥_s(Y)`$ its alleged volume should be obtained as the limit of the volumes of these cylinders. The volume of $`\pi _s^1𝒥_s(Y)`$ is $`[𝒥_s(Y)]𝕃^{ns}`$. The dimension of this element of $`_𝓀`$ is $`dim𝒥_s(Y)ns`$. If $`Y`$ is smooth of dimension $`nc`$, the we have seen that $`dim𝒥_s(Y)=(nc)(s+1)`$. Therefore, $`𝒥_{\mathrm{}}(Y)`$ is the intersection of cylinder sets whose volumes have dimension $`(nc)(s+1)ns=(nc)cs`$. For increasing $`s`$ these dimensions tend to negative infinity. But recall that in the ring $`\widehat{}_k`$ this is exactly the condition of convergence to zero. Thus the only sensible assignment of a volume to $`\mathrm{ord}_Y^1(\mathrm{})`$ is zero. This argument used that $`Y`$ is smooth to show that the dimension of $`𝒥_m(Y)`$ grows significantly slower than the dimension of $`𝒥_m(X)`$. This holds in general for singular $`Y`$ and we phrase it as a proposition whose proof however is postponed until Section 4.
###### Proposition 2.5.
Let $`YX`$ be a nowhere dense subvariety of $`X`$, then $`𝒥_{\mathrm{}}(Y)`$ is measurable and has measure $`\mu _X(𝒥_{\mathrm{}}(Y))`$ equal to zero.
To make this idea into a rigorous theory one has to define a larger class of measurable subsets of $`𝒥_{\mathrm{}}(X)`$ and this turns out to be somewhat subtle. In section 4 we will outline this briefly – but since, no matter what, the measure of $`\mathrm{ord}_Y^1(\mathrm{})`$ will be zero, we will move on at this point and start to integrate.
### 2.5. Definition and computation of the motivic integral
As before let $`X`$ be a smooth $`k`$-scheme and $`Y`$ a subscheme. We define the *motivic integral* of $`𝕃^{\mathrm{ord}_Y}`$ on $`X`$ as
$$_{𝒥_{\mathrm{}}(X)}𝕃^{\mathrm{ord}_Y}𝑑\mu _X=\underset{s=0}{\overset{\mathrm{}}{}}\mu (\mathrm{ord}_Y^1(s))𝕃^s.$$
Observe that the level set at infinity is already left out from this summation as it has measure zero.
Note that the sum on the right does converge since the virtual dimension of the summands approaches negative infinity.<sup>22</sup><sup>22</sup>22The fact that for stable sets $`AB`$ we have $`dim\mu (A)dim\mu (B)`$ applied to $`\mathrm{ord}_Y^1(s)𝒥_{\mathrm{}}(X)`$ gives that the dimension of $`\mu (\mathrm{ord}_Y^1(s))𝕃^s`$ is less or equal to $`ns`$. The notion of convergence in the ring $`\widehat{}_k`$ is such that this alone is enough to ensure the convergence of the sum. Thus it is justified to call $`𝕃^{\mathrm{ord}_Y}`$ integrable with integral as above.
It is useful to calculate at least one example. For $`Y=\mathrm{}`$ one has $`\mathrm{ord}_Y0`$ and thus we get
$$_{𝒥_{\mathrm{}}(X)}𝕃^{\mathrm{ord}_Y}𝑑\mu _X=\mu (\mathrm{ord}_Y^1(0))=[X]$$
where we used that $`X`$ is smooth. A less trivial example is $`Y`$ a smooth divisor in $`X`$. Then $`𝒥_s(Y)`$ is locally a $`𝔸^{(n1)s}`$–bundle over $`Y`$. The level set is $`\mathrm{ord}_Y^1(s)=\pi _{s1}^1𝒥_{s1}(Y)\pi _s^1𝒥_s(Y)`$ and, using that $`[𝒥_s(Y)]=[Y]𝕃^{(n1)s}`$, its measure is
$$[𝒥_{s1}(Y)]𝕃^{n(s1)}[𝒥_s(Y)]𝕃^{ns}=[Y](𝕃1)𝕃^s.$$
The integral of $`\mathrm{ord}_Y^1`$ is therefore
(3)
$$\begin{array}{cc}\hfill _{𝒥_{\mathrm{}}(X)}𝕃^{\mathrm{ord}_Y}𝑑\mu _X& =[XY]+\underset{s=1}{\overset{\mathrm{}}{}}[Y](𝕃1)𝕃^s𝕃^s\hfill \\ & =[XY]+[Y](𝕃1)𝕃^2\underset{s=0}{\overset{\mathrm{}}{}}𝕃^{2s}\hfill \\ & =[XY]+[Y](𝕃1)\frac{1}{𝕃^2(1𝕃^2)}\hfill \\ & =[XY]+[Y](𝕃1)(𝕃^21)^1\hfill \\ & =[XY]+\frac{[Y]}{𝕃+1}=[XY]+\frac{[Y]}{[^1]}\hfill \end{array}$$
Note the appearance of a geometric series in line 3 which is typical for these calculations (cf. Exercise 2.4). In fact, the motivic volumes of a wide class of measurable subsets (namely, the semi-algebraic subsets of Denef and Loeser in ) belong to the ring generated by the image of $`_𝓀`$ under the completion map, and the sums of geometric series with denominators $`𝕃^j`$, $`j>0`$ \[11, Corollary 5.2\].
###### Exercise 2.6.
<sup>23</sup><sup>23</sup>23The answer in the first case is $`[XY]+[Y]\frac{𝕃^c1}{𝕃^{c+1}1}`$. Otherwise it can be read off from the general formula below; e.g. in the second case it is $`[XD]+\frac{[D]}{[^a]}`$.
Compute in a similar fashion the motivic integrals $`𝕃^{\mathrm{ord}_Y}`$ where $`Y`$ is as follows.
1. $`Y`$ is a smooth subscheme of codimension $`c`$ in $`X`$.
2. $`Y=aD`$ where $`D`$ is a smooth divisor and $`a`$.
3. $`Y=D_1+D_2`$ where the $`D_i`$’s are smooth and in normal crossing.
4. $`Y=a_1D_1+a_2D_2`$ with $`D_i`$ as above and $`a_i`$ positive integers.
These computations are a special case of a formula which explicitly computes the motivic integral over $`𝕃^{\mathrm{ord}_Y}`$ where $`Y`$ is an effective divisor with normal crossing support.
###### Proposition 2.6.
Let $`Y=_{i=1}^sr_iD_i`$ ($`r_i>0`$) be an effective divisor on $`X`$ with normal crossing support and such that all $`D_i`$ are smooth. Then
$$_{𝒥_{\mathrm{}}(𝒳)}𝕃^{\mathrm{ord}_Y}𝑑\mu _X=\underset{J\{1,\mathrm{},s\}}{}[D_J^{}](\underset{jJ}{}\frac{𝕃1}{𝕃^{r_j+1}1})=\underset{J\{1,\mathrm{},s\}}{}\frac{[D_J^{}]}{_{jJ}[^{r_j}]}$$
$`D_J=_{jJ}D_j`$ (note that $`D_{\mathrm{}}=X`$) and $`D_J^{}=D_J_{jJ}D_j`$.
The proof of this is a computation entirely similar to (though significantly more complicated than) the one carried out in (3) above; for complete detail see either Batyrev \[2, Theorem 6.28\] or Craw \[7, Theorem 1.17\]. We suggest doing it as an exercise using the following lemma.
###### Lemma 2.7.
For $`J\{1,\mathrm{},s\}`$ redefine $`r_i=0`$ if $`iJ`$. Then
$$\mu (\mathrm{ord}_{D_j}^1(r_j))=[D_J^0](𝕃1)^{|J|}𝕃^{{\scriptscriptstyle r_j}}$$
where $`D_J^0=_{jJ}D_j_{jJ}D_J`$ and $`|J|`$ denotes the cardinality of $`J`$.
###### Exercise 2.7.
<sup>24</sup><sup>24</sup>24Hint/Solution: Clearly, for an arc $`\vartheta `$ membership in $`\mathrm{ord}_{D_i}^1(r_i)`$ only depends on its truncation $`\pi _t(\vartheta )`$ as long as $`t\mathrm{max}\{r_i\}`$. With the notation of Example 2.1 we write
$$\pi _t(\vartheta (x_i))=\underset{j=0}{\overset{t}{}}\vartheta _i^{(j)}t^j$$
such that $`\vartheta `$ is determined by the coefficients $`\vartheta _i^{(j)}`$. Now we determine what the condition $`\vartheta \mathrm{ord}_{D_i}^1(r_i)`$ imposes on the coefficients $`\vartheta _i^{(j)}`$ (here it is convenient to set $`r_i=0`$ for $`i>s`$). For $`j<r_i`$ one has $`\vartheta _i^{(j)}=0`$. For $`j=r_i`$ one has $`\vartheta _i^{(j)}0`$. For $`j>r_i`$ one has no condition $`\vartheta _i^{(j)}`$. Thus the set $`\pi _t(_i\mathrm{ord}_{D_i}^1(r_i))`$ is the product made up form the factors $`D_{\mathrm{supp}r}^0`$ corresponding to all possible $`\vartheta _i^{(0)}`$, a copy of $`k\{0\}`$ for each possible $`\vartheta _i^{(r_i)}`$ and $`r_i>0`$ and a copy of $`k`$ for each $`r_i<jk`$. Putting this together we obtain the above formula.
Show that for $`t\mathrm{max}\{r_i\}`$ one has locally an isomorphism of the $`k`$-points
$$\pi _t(_i\mathrm{ord}_{D_i}^1(r_i))D_{\mathrm{supp}r}^0\times (k\{0\})^{|\mathrm{supp}r|}\times k^{nt{\scriptscriptstyle r_i}}.$$
Show that this implies the lemma. To prove the preceding statement reduce to the case that $`X=U𝔸^n`$ is an open subvariety of $`𝔸^n`$ and $`D_i`$ is given by the vanishing of $`x_1\mathrm{}x_s`$ where $`x_1,\mathrm{},x_n`$ is a local system of coordinates (use Proposition 2.2). Then finish this case using the description of the arc space of $`𝔸^n`$ as given in Example 2.1.
The explicit formulas of Proposition 2.6 and Lemma 2.7 are one cornerstone underlying many applications of motivic integration. The philosophy one employs is to encode information in a motivic integral, then using the transformation rule of the next section the computation of this integral can be reduced to the computation of an integral over $`𝕃^{\mathrm{ord}_Y}`$ for $`Y`$ a normal crossing divisor. In this case the above formula gives the answer. Thus we shall proceed to the all important birational transformation rule for motivic integrals.
## 3. The transformation rule
The power of the theory stems from a formula describing how the motivic integral transforms under birational morphisms:
###### Theorem 3.1.
Let $`X^{}\stackrel{𝑓}{}X`$ be a proper birational morphism of smooth $`k`$-schemes and let $`D`$ be an effective divisor on $`X`$, then
$$_{𝒥_{\mathrm{}}(X)}𝕃^{\mathrm{ord}_D}𝑑\mu _X=_{𝒥_{\mathrm{}}(X^{})}𝕃^{\mathrm{ord}_{f^1D+K_{X^{}/X}}}𝑑\mu _X^{}.$$
As the relative canonical sheaf $`K_{X^{}/X}`$ is defined by the Jacobian ideal of $`f`$, this should be thought of as the change of variables formula for the motivic integral.
As a warmup for the proof, we verify the the transformation rule first in the special case of blowing up a smooth subvariety and $`D=\mathrm{}`$. Let $`X^{}=\mathrm{Bl}_YX`$ be the blowup of $`X`$ along the smooth center $`Y`$ of codimension $`c`$ in $`X`$. Then by \[26, Exercise II.8.5\] the relative canonical divisor is $`K_{X^{}/X}=(c1)E`$, where $`E`$ is the exceptional divisor of the blowup. Then, using Proposition 2.6 in its simplest incarnation we compute
$$\begin{array}{cc}\hfill _{𝒥_{\mathrm{}}(X^{})}𝕃^{\mathrm{ord}_{K_{X^{}/X}}}𝑑\mu _X^{}& =_{𝒥_{\mathrm{}}(X^{})}𝕃^{\mathrm{ord}_{(c1)E}}𝑑\mu _X^{}\hfill \\ & =[X^{}E]+\frac{[E]}{[^c]}\hfill \\ & =[XY]+[Y]=[X],\hfill \end{array}$$
where we used that $`E`$ is a $`^c`$–bundle over $`Y`$ (by definition of blowup) and therefore $`[E]=[Y][^c]`$.
#### 3.0.1. The induced map on the arc space
The proper birational map $`f`$ induces a map $`f_{\mathrm{}}=𝒥_{\mathrm{}}(f):𝒥_{\mathrm{}}(X^{})\stackrel{}{}𝒥_{\mathrm{}}(X)`$. The first task will be to show that away from a set of measure zero $`f_{\mathrm{}}`$ is a bijection (of sets). Let $`\mathrm{\Delta }X^{}`$ be the locus where $`f`$ is not an isomorphism. We show that every arc $`\gamma :\mathrm{Spec}kt\stackrel{}{}X`$, which does not entirely lie in $`f(\mathrm{\Delta })`$ uniquely lifts to an arc in $`X^{}`$. For illustration consider the diagram
Observe that by assumption the generic point $`\mathrm{Spec}kt_{(0)}`$ of $`\mathrm{Spec}kt`$ does lie in $`Xf(\mathrm{\Delta })`$ and since $`f`$ is an isomorphism over $`X^{}\mathrm{\Delta }`$ it thus lifts to $`X^{}`$ uniquely (dashed arrow). Now the valuative criterion for properness (see \[26, Chapter 2, Theorem 4.7\]) yields the unique existence of the dotted arrow. Thus the map $`f_{\mathrm{}}:(𝒥_{\mathrm{}}(X^{})𝒥_{\mathrm{}}(\mathrm{\Delta }))\stackrel{}{}(𝒥_{\mathrm{}}(X)𝒥_{\mathrm{}}(f(\mathrm{\Delta })))`$ is a bijection of $`k`$-valued points. Since $`𝒥_{\mathrm{}}(\mathrm{\Delta })`$ has measure zero it can be safely ignored and we will do so in the following.
###### Exercise 3.1.
<sup>25</sup><sup>25</sup>25Solution: The valuative criterion of properness shows that any $`\gamma 𝒥_{\mathrm{}}(Z)`$ lies in the image (where $`ZX`$ is such that $`f`$ is an isomorphism from $`X^{}f^1(Z)\stackrel{}{}XZ`$. But for any $`\gamma _m𝒥_m(X)`$ the cylinder $`\pi _m^1(\gamma _m)`$ cannot be contained in $`𝒥_{\mathrm{}}(Z)`$ since the latter has measure zero.
Let $`f:X^{}\stackrel{}{}X`$ be a proper birational map (of smooth varieties). Then for every $`m`$ the map $`f:𝒥_m(X^{})\stackrel{}{}𝒥_m(X)`$ is surjective.
###### Proof of transformation rule.
The level sets $`C_e^{}=\mathrm{ord}_{K_{X^{}/X}}^1(e)`$ partition $`𝒥_{\mathrm{}}(X^{})`$, and cutting into even smaller pieces according to order of contact along $`f^1(D)`$ we define
$$C_{e,k}^{}=C_e^{}\mathrm{ord}_{f^1D}^1(k)\text{ and }C_{e,k}=f_{\mathrm{}}(C_{e,k}^{})$$
to get the following partitions (up to measure zero by 3.0.1) of the arc spaces
$$𝒥_{\mathrm{}}(X^{})=C_{e,k}^{}\text{ and }𝒥_{\mathrm{}}(X)=C_{e,k}.$$
The essence of the proof of the transformation rule is captured by the following two crucial facts.
1. $`C_{e,k}`$ are stable sets for all $`e,k`$.
2. $`\mu _X(C_{e,k})=\mu _X^{}(C_{e,k}^{})𝕃^e`$.
Two different proofs of these facts will occupy the remainder of this section. Using these facts, the transformation rule is a simple calculation:
$$\begin{array}{cc}\hfill _{𝒥_{\mathrm{}}(X)}𝕃^{\mathrm{ord}_D}𝑑\mu _X& =\underset{k}{}\mu (\mathrm{ord}_D^1(k))𝕃^k=\underset{k}{}\left(\underset{e}{}\mu (C_{e,k})\right)𝕃^k\hfill \\ & =\underset{e,k}{}\mu (C_{e,k}^{})𝕃^e𝕃^k\hfill \\ & =\underset{t}{}\left(\underset{e+k=t}{}\mu (C_{e,k}^{})\right)𝕃^t=\underset{t}{}\mu (\mathrm{ord}_{f^1D+K_{X^{}/X}}^1(t))𝕃^t\hfill \\ & =_{𝒥_{\mathrm{}}(X^{})}𝕃^{\mathrm{ord}_{f^1D+K_{X^{}/X}}}𝑑\mu _X^{}\hfill \end{array}$$
Besides facts (a) and (b) one uses that if a cylinder $`B`$ is written as a disjoint union of cylinders $`B_i`$ then the measure $`\mu (B)=\mu (B_i)`$. To check that this is correct one has to use the precise definition of the motivic measure which we have avoided until now. See the Section 4 for more details on this. ∎
We first treat properties (a) and (b) in a special case, namely the blowup at a smooth center. In the applications of motivic integration to birational geometry which we discuss below, we can always put ourselves in the favorable situation that the birational map in consideration is a sequence of blowing ups along smooth centers, hence already this simple version goes very far. Furthermore, using the Weak Factorization Theorem of one can make a general proof of the transformation rule by reducing to this case. However, despite the adjective *weak* in the Weak Factorization Theorem it is a very deep and difficult result and its use in the proof of the Transformation rule is overkill. Therefore we give in Appendix A.2 an essentially elementary (though at the first reading somewhat technical) proof following .
### 3.1. Images of cylinders under birational maps.
We start with some basic properties of the behaviour of cylinders under birational morphisms, see . These will be useful also in Section 5.
###### Proposition 3.2.
Let $`f:X^{}\stackrel{}{}X`$ be a proper birational map of smooth varieties. Let $`C^{}=(\pi _m^X^{})^1(B^{})𝒥_{\mathrm{}}(X^{})`$ be a cylinder such that $`B^{}`$ is a union of fibers of $`f_m`$, then $`C\stackrel{\mathrm{def}}{=}f_{\mathrm{}}(C^{})=(\pi _m^X)^1(f_mB^{})`$ is a cylinder in $`𝒥_{\mathrm{}}(X)`$.
###### Proof.
Clearly it is enough to show that $`C=\pi _m^1(B)`$ since $`B=f_m(B^{})`$ is constructible (being the image of a constructible set under a finite type morphism). The nontrivial inclusion is $`\pi _m^1(B)C`$. Let $`\gamma \pi _m^1(B)C`$ and consider for every $`pm`$, the cylinder
$$D_p\stackrel{\mathrm{def}}{=}(\pi _p^X^{})^1(f_p^1(\pi _p^X(\gamma )))$$
which is non empty by exercise 3.1. Clearly, $`D_pD_{p+1}`$ which implies that the $`D_p`$ form a decreasing sequence of nonempty cylinders. By Proposition 4.3 the intersection of all the $`D_p`$ is nonempty. Now let $`\gamma ^{}D_p`$ and clearly $`f_{\mathrm{}}\gamma ^{}=\gamma `$. Furthermore, since $`B^{}`$ is a union of fibers of $`f_m`$ we have $`C^{}D_m`$ and hence $`C^{}D_p`$ for all $`pm`$. Hence $`\pi _m(\gamma ^{})B^{}`$. ∎
###### Exercise 3.2.
<sup>26</sup><sup>26</sup>26Solution: The previous proposition shows that $`f_{\mathrm{}}(𝒥_{\mathrm{}}(X^{}))`$ is a cylinder. The surjectivity of $`f_m`$ which was showed in exercise 3.1 now implies the result (in fact even the surjectivity of $`f_0=f`$ is enough at this point).
Let $`f:X^{}\stackrel{}{}X`$ be a proper birational map of smooth varieties. Then $`f_{\mathrm{}}:𝒥_{\mathrm{}}(X^{})\stackrel{}{}𝒥_{\mathrm{}}(X)`$ is surjective.
#### 3.1.1. The key technical result for blowup at smooth center.
We now proceed to showing the key technical result used in the proof of the transformation formula.
###### Theorem 3.3 (Denef, Loeser, ).
Let $`f:X^{}\stackrel{}{}X`$ be a proper birational morphism of smooth varieties. Let $`C_e^{}=\mathrm{ord}_{K_{X^{}/X}}^1(e)`$ where $`K_{X^{}/X}`$ is the relative canonical divisor and let $`C_e\stackrel{\mathrm{def}}{=}f_{\mathrm{}}C_e^{}`$. Then, for $`m2e`$
1. the fiber of $`f_m`$ over a point $`\gamma _m\pi _mC_e`$ lies inside a fiber of $`\pi _{me}^m`$.
2. $`\pi _m(C_e^{})`$ is a union of fibers of $`f_m`$.
3. $`f_m:\pi _mC_e^{}\stackrel{}{}\pi _mC_e`$ is piecewise trivial $`𝔸^e`$–fibration.
###### Corollary 3.4.
With the notation as in the theorem, $`C_e`$ is a stable at level $`m2e`$ and $`[C_e^{}]=[C_e]𝕃^e`$.
###### Proof of Theorem 3.3.
First note (a’) implies (a) since $`C_e^{}`$ is stable at level $`e`$. Furthermore using the following Lemma 3.5 it is enough to show that a fibers of $`f_m`$ over a point in $`C_e`$ is affine space $`𝔸^e`$.
We give the proof here only in the case that $`X^{}=\mathrm{Bl}_YX\stackrel{}{}X`$ where $`Y`$ is a smooth subvariety of $`X`$, for the general argument see Appendix A.2 below. Since $`X`$ is smooth there is an étale morphism $`\phi :X\stackrel{}{}𝔸^n`$ such that $`\phi (Y)`$ is given by the vanishing of the first $`nc`$ coordinates on $`𝔸^n`$ and $`\phi ^1(\phi (Y))=Y`$. By the étale invariance of jet schemes (Proposition 2.2) one can hence assume that $`Y=𝔸^{nc}𝔸^n=X`$.<sup>27</sup><sup>27</sup>27This is not as straight froward as as it might seem. One has to check that $`\mathrm{Bl}_YXX\times _𝔸^n\mathrm{Bl}_{\phi (Y)}𝔸^n`$. Let $`I`$ be the ideal defining $`Y`$ and $`\phi _{}(I)`$ the ideal defining $`\phi (Y)`$. Locally one has $`\phi ^{}(\phi _{}(I))=I`$. The pullback of $`\phi _{}(I)`$ to $`\mathrm{Bl}_YX`$ is equal to the pullback of $`I`$ which is by construction of the blow-up a principal ideal sheaf. By the universal property of the blowup (applied to $`\mathrm{Bl}_{\phi (Y)}𝔸^n`$) we get a map $`\mathrm{Bl}_YX\stackrel{}{}\mathrm{Bl}_{\phi (Y)}𝔸^n`$ and hence a map
$$\mathrm{Bl}_YX\stackrel{}{}\mathrm{Bl}_{\phi (Y)}𝔸^n\times _{𝔸^n}X$$
The pullback of $`I`$ to $`\mathrm{Bl}_{\phi (Y)}𝔸^n\times _{𝔸^n}X`$ along the second projection is equal to the pullback of $`\phi _{}(I)`$ along the natural map to $`𝔸^n`$ which factors through the blowup. Hence this pullback is locally principal and again by the universal property of blowup (applied to $`\mathrm{Bl}_YX\stackrel{}{}X`$) we get a map in the opposite direction, which is easily verified to be inverse to the one given above. This implies in particular that $`\mathrm{Bl}_YX`$ is étale over $`\mathrm{Bl}_{\phi (Y)}𝔸^n`$ and thus
$$𝒥_m(\mathrm{Bl}_YX)\mathrm{Bl}_YX\times _{\mathrm{Bl}_{\phi (Y)}𝔸^n}𝒥_m(\mathrm{Bl}_{\phi (Y)}𝔸^n)X\times _{𝔸^n}𝒥_m(\mathrm{Bl}_{\phi (Y)}𝔸^n).$$
This allows one after fixing the base point $`xX`$ of $`\gamma `$ to really reduce to the case of affine space.
To further simplify notation (and notation only) we assume that $`n=3`$ and $`c=3`$, that is we only have to consider the blowup of the origin in $`𝔸^3`$. By definition
$$X^{}=\mathrm{Bl}_0𝔸^3𝔸^3\times ^2$$
is given by the vanishing of the $`2\times 2`$ minors of the matrix
$$\left(\begin{array}{ccc}x_0& x_1& x_2\\ y_0& y_1& y_2\end{array}\right)$$
where $`(x_0,x_1,x_2)`$ and $`(z_0,z_1,z_2)`$ are the coordinates on $`𝔸^3`$ and the homogeneous coordinates of $`^2`$ respectively. Due to the local nature of our question it is enough to consider on affine patch of $`X^{}`$, say the one determined by $`z_0=1`$. The equations of the minors then reduce to $`x_1=x_0z_1`$ and $`x_2=x_0z_2`$ such that on this patch the map $`X^{}\stackrel{}{}X`$ is given by the inclusion of polynomial rings
$$k[x_0,x_0z_1,x_0z_2]\stackrel{𝑓}{}k[x_0,z_1,z_2].$$
The exceptional divisor $`E`$ is hence given by the vanishing of $`x_0`$. The relative canonical divisor $`K_{X^{}/X}`$ is equal to $`2E`$ since $`det(\mathrm{Jac}(f))`$ is easily computed to be $`x_0^2`$.
Now let $`\gamma ^{}C_e^{}`$. In our local coordinates, $`\gamma _m^{}=\pi _m(\gamma ^{})`$ is uniquely determined by the three truncated powerseries
$$\begin{array}{cc}\hfill \gamma _m^{}(x_0)& =t^{e/2}\underset{i=0}{\overset{me/2}{}}a_it^i\text{with }a_00\hfill \\ \hfill \gamma _m^{}(z_1)& =\underset{i=0}{\overset{m}{}}b_it^i\text{and}\hfill \\ \hfill \gamma _m^{}(z_2)& =\underset{i=0}{\overset{m}{}}c_it^i,\hfill \end{array}$$
where the special shape of the first one comes from the condition that $`\gamma ^{}`$ has contact order with $`K_{X^{}/X}=2E`$ precisely equal to $`e`$ (we only have to consider $`e`$ which are divisible by $`c1=2`$ since otherwise $`C_e^{}`$ is empty). Its image $`\gamma _m=f_m(\gamma _m^{})=\gamma _m^{}f_m`$ is analogously determined by the three truncated powerseries
(4)
$$\begin{array}{cc}\hfill \gamma _m(x_0)& =t^{e/2}\underset{i=0}{\overset{me/2}{}}a_it^i\text{with }a_00\hfill \\ \hfill \gamma _m(x_0z_1)& =\gamma _m^{}(x_0)\gamma _m^{}(z_1)=t^{e/2}\underset{i=0}{\overset{me/2}{}}a_it^i\underset{i=0}{\overset{m}{}}b_it^imodt^{m+1}\hfill \\ \hfill \gamma _m(x_0z_2)& =\gamma _m^{}(x_0)\gamma _m^{}(z_1)=t^{e/2}\underset{i=0}{\overset{me/2}{}}a_it^i\underset{i=0}{\overset{m}{}}c_it^imodt^{m+1}\hfill \end{array}$$
Expanding the product of the sums in the last two equations of (4) we observe that (due to the occurence of $`t^{e/2}`$) the coefficients $`b_{m\frac{e}{2}+1},\mathrm{},b_m`$ are not visible in $`\gamma _m(x_0z_1)`$ since they only appear as coefficients of some $`t^k`$ for $`k>m`$. Analogously, $`\gamma _m(x_0z_2)`$ does not depend on $`c_{m\frac{e}{2}+1},\mathrm{},c_m`$. Conversely, given
$$\gamma _m(x_0z_1)=t^{e/2}\underset{i=0}{\overset{m\frac{e}{2}}{}}\beta _it^i$$
and knowing all $`a_i`$’s (which equally show up in $`\gamma _m^{}`$ and $`\gamma _m`$) we can inductively recover the $`b_i`$’s:
$$\begin{array}{cc}\hfill b_0& =(\beta _0)/a_0\text{(note: }a_00\text{)}\hfill \\ \hfill b_1& =(\beta _1a_1b_0)/a_0\hfill \\ \hfill b_2& =(\beta _2(a_2b_0+a_1b_1))/a_0\hfill \\ & \mathrm{}\hfill \\ \hfill b_t& =(\beta _t(a_tb_0+a_{t1}b_1+\mathrm{}+a_1b_{t1}))/a_0\hfill \end{array}$$
and this works until $`t=m\frac{e}{2}`$ since $`a_t`$ is known for $`tm\frac{e}{2}`$. The analogous statements of course hold also for the $`c_i`$’s.
Summing up these observations we see that the fiber of $`f_m`$ over $`\gamma _m`$ is an affine space of dimension $`e=2\frac{e}{2}`$, namely it is spanned by the last $`\frac{e}{2}`$ of the coefficients $`b_i`$ and $`c_i`$. This proves part (b). Furthermore, any two $`\gamma _m^{}`$ and $`\gamma _m^{\prime \prime }`$ mapping via $`f_m`$ to $`\gamma _m`$ only differ in these last $`\frac{e}{2}`$ coefficients, hence they become equal after further truncation to level $`me`$.<sup>28</sup><sup>28</sup>28Even to level $`m\frac{e}{2}`$ in this case of blowup of a point in $`𝔸^3`$. In general, for a blowup of $`c`$ codimensional smooth center it is truncation to level $`m\frac{e}{c1}`$ which suffices. Hence uniformly it is truncation to level $`me`$ which works. This shows $`(a^{})`$ and the proof is finished. ∎
###### Lemma 3.5.
Let $`\phi :V\stackrel{}{}W`$ be a morphism of finite type schemes such that all fibers $`\phi ^1(x)𝔸^e\times k(x)`$, then $`\phi `$ is a piecewise trivial $`𝔸^e`$–fibration.<sup>29</sup><sup>29</sup>29In fact, it follows from Hilbert’s Theorem 90 that a piecewise trivial $`𝔸^e`$ fibration is actually locally trivial. EXPLAIN!
###### Proof.
We may assume that $`W`$ is irreducible. Then the fiber over the generic point $`\eta `$ of $`W`$ is by assumption isomorphic to $`𝔸^e`$. This means that there is an open subset $`UW`$ such that $`f^1(U)𝔸^e\times U`$.<sup>30</sup><sup>30</sup>30The isomorphism $`𝔸^e(k(\eta ))\phi ^1(\eta )=V\times _W\mathrm{Spec}k(\eta )`$ is defined via some finitely many rational functions on $`W`$. For any $`UW`$ such that these are regular we get $`\phi ^1(U)𝔸^e\times U`$. Now $`\phi `$ restricted to the complement of $`U`$ is a map of the same type but with smaller dimensional base and we can finish the argument by induction. Even though we did not make this explicit in the proof of Theorem 3.3 the statement about the fibers being equal to $`𝔸^e`$ holds for all fibers (not only fibers over closed points). ∎
###### Exercise 3.3.
<sup>31</sup><sup>31</sup>31We may assume that $`C`$ is irreducible. Let $`e`$ be smallest such that $`C\mathrm{Cont}_{K_{X^{}/X}}^e0`$. This intersection is open and dense in $`C`$, and hence we may replace $`C`$ by $`C\mathrm{Cont}_{K_{X^{}/X}}^e`$. Now Theorem 3.3 applies.
Use Proposition 3.2 and 3.3 to show that if $`C𝒥_{\mathrm{}}(X^{})`$ is a cylinder, then the closure of $`f_{\mathrm{}}(C)𝒥_{\mathrm{}}(X)`$ is a cylinder, where $`f:X^{}\stackrel{}{}X`$ is a proper birational morphism.
### 3.2. Proof of transformation rule using Weak Factorization
With the proof given so far we have the transformation formula available for a large class of proper birational morphisms, namely the ones which are obtained as a sequence of blowups along smooth centers. So in particular we have the result for all resolutions of singularities, which is the only birational morphism we will consider in our applications later.
Let me finish by outlining how using the Weak Factorization Theorem one can make a full proof out of this. Let us first recall the statement:
###### Theorem 3.6 (Weak Factorization Theorem ).
Let $`\phi :X^{}X`$ be a birational map between smooth complete varieties over $`k`$ of characteristic zero. Then $`\phi `$ can be factored
such that all indicated maps are a blowup at a smooth center. Furthermore, there is an index $`i`$ such that the rational maps to $`X^{}`$ to the left ($`X^{}X_j`$ for $`ji`$) and the rational maps to $`X`$ to the right ($`X_jX`$ for $`ji`$) of that index are in fact regular maps.
One should point out that the second part of the Theorem about the regularity of the maps is crucial in many of its applications, and particularly in the application that follows next.
###### Proof of Theorem 3.1.
So far we have proved the Transformation rule for the blowup along a smooth center. Given a proper birational morphism of smooth varieties $`f:X^{}\stackrel{}{}X`$ we can factor it into a chain as in the Weak Factorization Theorem. In particular, for each birational map in that chain, the Transformation rule holds. The second part of the Factorzation Theorem together with the assumption that $`X^{}\stackrel{}{}X`$ is a morphism implies that $`X_{n1}\stackrel{}{}X`$ is also a morphism. Part (b) of the following Exercise shows that for this morphism $`X_{n1}\stackrel{}{}X`$ the Transformation rule holds. Now the shorter chain ending with $`X_{n2}\stackrel{}{}X`$ is again a chain such that for each map the Transformation rule holds, by part (a) of that exercise. By induction we can conclude that the transformation rule holds for $`f`$ itself. ∎
###### Exercise 3.4.
<sup>32</sup><sup>32</sup>32Solution: Part (a) is easy. Not sure if part (c) is really feasible. Part (d) is straightforward from (c). Let me outline part (b): Recall that $`K_{X^{}/X}=K_{X^{}/X^{}^{}}+f^{}{}_{}{}^{}{}_{}{}^{}K_{X^{}{}_{}{}^{}/X}^{}`$ such that
$$\begin{array}{cc}\hfill _X𝕃^{\mathrm{ord}_D}𝑑\mu _X& =_X^{}𝕃^{\mathrm{ord}_{f_{}^{}{}_{}{}^{}DK_{X^{}/X}}}𝑑\mu _X^{}\hfill \\ & =_X^{}𝕃^{\mathrm{ord}_{f_{}^{}{}_{}{}^{}Df^{}{}_{}{}^{}{}_{}{}^{}K_{X^{}{}_{}{}^{}/X}^{}+K_{X^{}/X^{}^{}}}}𝑑\mu _X^{}=_X^{}^{}𝕃^{\mathrm{ord}_{f^{}DK_{X^{}{}_{}{}^{}/X}}}𝑑\mu _X^{}^{}\hfill \end{array}$$
where the first and last equality is the Transformation rule for $`f^{}^{}`$ and $`f`$.
Suppose one is given a commuting diagram of proper birational morphisms
1. If the Transformation rule holds for $`f^{\prime \prime }`$ and $`f`$, then also for $`f^{}=f^{\prime \prime }f`$.
2. If the Transformation rule holds for $`f^{\prime \prime }`$ and $`f^{}`$, then it also holds for $`f`$.
3. The same as (a) and (b) but with “the Transformation rule” replaced by “the conclusions of Theorem 3.3”. (This part is more difficult than the others)
4. Using (c) and the Weak Factorization Theorem produce a proof of Theorem 3.3 building on the case of the blowup at a smooth center we considered above.
## 4. Brief outline of a formal setup for the motivic measure.
In this section we fill in some details that were brushed over in our treatment of motivic integration so far. First this concerns some basic properties and the well–definedness of the motivic measure and integral.
### 4.1. Properties of the motivic measure
For simplicity we still assume that $`X`$ is a smooth $``$–variety. Recall that we defined for a stable set $`C=\pi _m^1(B)`$ a volume by setting
$$\mu _X(C)=[B]𝕃^{nm}_𝓀.$$
It is easy to verify that on stable sets (i.e. cylinders if $`X`$ is smooth) the measure is additive on finite disjoint unions. Furthermore, for stable sets $`CC^{}`$ one has $`dim\mu _X(C)dim\mu _X(C^{})`$.<sup>33</sup><sup>33</sup>33Check these assertions as an exercise. We begin with a rigorous definition of what is a measurable set extending the above definition.
###### Definition 4.1.
A subset $`C𝒥_m(X)`$ is called *measurable* if for all $`n`$ there is a stable set $`C_n`$ and stable sets $`D_{n,i}`$ for $`i`$ such that
$$C\mathrm{\Delta }C_n\underset{i}{}D_{n,i}$$
and $`dim\mu (D_{n,i})n`$ for all $`i`$. Here $`C\mathrm{\Delta }C_n=(CC_n)(C_nC)`$ denotes the symmetric difference of two sets. In this case the *volume* of $`C`$ is
$$\mu _X(C)=\underset{n\stackrel{}{}\mathrm{}}{lim}\mu _X(C_n)_𝓀.$$
This limit converges and is independent of the $`C_n`$’s.
The key point in proving the claims in the definition<sup>34</sup><sup>34</sup>34In a more restrictive definition of measurable is used, namely he requires that $`dimD_{n,i}(n+i)`$. This has the advantage that one does not require the field to be uncountable to conclude the well definedness of the measure. Essentially, he uses that if $`DD_i`$ are cylinders with $`limdimD_i=\mathrm{}`$, then $`D`$ is already contained the union of finitely many of the $`D_i`$. This is true if $`k`$ is infinite, and, if $`k`$ is uncountable, even true without the assumption on the $`D_i`$. The advantage of Looijenga’s setup is that one is not bound to an uncountable field, but unfortunatly I was not able to verify that in his setup $`𝒥_{\mathrm{}}(Y)`$ is measurable and has zero volume. is the so called Baire property of constructible subsets of a $``$–variety which crucially uses the fact that $``$ is uncountable, see \[22, Corolaire 7.2.6\].
###### Proposition 4.2.
Let $`K_1K_2K_3\mathrm{}`$ be an infinite sequence of nonempty constructible subsets of a $``$–variety $`X`$. Then $`_iK_i`$ is nonempty.
For cylinder sets this implies the following Proposition, which will be used below in a version which asserts that a cylinder $`C`$ which is contained in the union of countably many cylinders $`C_i`$ is already contained in the union of finitely many of these.
###### Proposition 4.3.
Let $`C_1C_2C_3\mathrm{}`$ be an infinite sequence of nonempty cylinders in $`𝒥_{\mathrm{}}(X)`$ where $`X`$ is a smooth $``$–variety. Then $`_iC_i`$ is nonempty.
###### Proof.
By definition of a cylinder and Chevalley’s theorem $`\pi _0(C_i)`$ is a constructible subset of $`X`$. Thus we can apply Proposition 4.2 to the sequence $`\pi _0(C_1)\pi _0(C_2)\mathrm{}`$ to obtain an element $`x_0\pi _0(C_i)`$. Now consider the sequence of cylinders $`C_i^{}\stackrel{\mathrm{def}}{=}C_i\pi _0^1(x_0)`$ and repeat the argument for the sequence of constructible sets $`\pi _1(C_1^{})\pi _1(C_2^{})\pi _1(C_2^{})\mathrm{}`$ to obtain an element $`x_1\pi _1(C_i^{})`$. Repeating this procedure we successively lift $`x_0X`$ to $`x_1𝒥_1(X)`$, $`x_2𝒥_2(X)`$ and so forth. The limit of these $`x_i`$ gives an element $`x𝒥_{\mathrm{}}(X)`$ which is in the intersection of all $`C_i`$. Thus this intersection is nonempty. ∎
On a possibly singular $``$–variety the statement (and proof) is true if “cylinder” is replaced by “stable set” above. If $`X`$ is a $`k`$–variety with $`k`$ at most countable these statements might be false. Nevertheless, if one views the constructible sets (resp. cylinders) not as subsets of the $`k`$–rational points but rather as certain sub–arragnements (something weaker than a subfunctor) of the functor of points represented by $`X`$ the above is essentially true. This is the point of view of Denef and Loeser and is carried out in .
###### Justification of Definition 4.1.
For both claims it suffices to show $`dim(\mu _X(C_i)\mu _X(C_j^{}))i`$ for all $`ji`$, where the prime indicates a second set of defining data as is in the definition. We have
$$C_iC_j^{}(C\mathrm{\Delta }C_i)(C\mathrm{\Delta }C_j^{})\underset{m}{}C_{i,m}\underset{m}{}C_{j,m}^{}.$$
Since $`(C_iC_j^{})`$ and all terms on the right are cylinders the previous proposition applies and $`C_iC_j^{}`$ is contained in finitely many of the cylinders of the right hand side. This implies that $`dim\mu _X(C_iC_j^{})i`$. The same applies to $`dim\mu _X(C_j^{}C_i)`$. Using that $`C_i=(C_iC_j^{})(C_iC_j^{})`$ and $`C_j^{}=(C_iC_j^{})(C_j^{}C_i)`$ and that $`\mu _X`$ is additive on finite disjoint unions of cylinders we get
$$\begin{array}{cc}\hfill dim(\mu _X(C_i)\mu _X(C_j^{}))& =dim(\mu _X(C_iC_j^{})\mu _X(C_j^{}C_i))\hfill \\ & \mathrm{max}\{dim\mu _X(C_iC_j^{}),dim\mu _X(C_j^{}C_i)\}i.\hfill \end{array}$$
For once this shows that the $`\mu _X(C_i)`$ form a Cauchy sequence, thus the limit exists as claimed. Secondly it immediately follows that this limit does not depend on the chosen data. ∎
We summarize some basic properties of the measure and measurable sets.
###### Proposition 4.4.
The measurable sets form an algebra of sets. $`\mu _X`$ is additive on disjoint unions, thus $`\mu _X`$ is a pre–measure in classical terminology. If $`C_i`$ are a infinite disjoint sequence of measurable sets such that
$$\underset{i\stackrel{}{}\mathrm{}}{lim}\mu _X(C_i)=0$$
then $`C=C_i`$ is measurable and $`\mu _X(C)=\mu _X(C_i)`$.
###### Proof.
The verification of all parts is quite easy. As an example we only show that the complement of a measurable set is also measurable and leave the rest as an exercise.
If $`C`$ is measurable then cylinders $`C_i`$ and $`C_{i,j}`$ can be chosen with the properties as in Definition 4.1. The complements $`C_i^c=𝒥_{\mathrm{}}(X)C_i`$ are cylinders. Since $`C^c\mathrm{\Delta }C_j^c=C\mathrm{\Delta }C_j`$ it follows at once that the complement of $`C`$ is measurable. ∎
The following proposition was one of the missing ingredients for the setup of motivic integration we outlined so far. It ensures that our typical functions $`𝕃^{\mathrm{ord}_Y}`$ are in fact measurable (the missing part was the level set at infinity $`𝒥_{\mathrm{}}(Y)`$ which we owe the proof that it is measurable with measure zero).
###### Proposition 4.5.
Let $`YX`$ be a locally closed subvariety. Then $`𝒥_{\mathrm{}}(Y)`$ is a measurable subset of $`𝒥_{\mathrm{}}(X)`$ and if $`dimY<dimX`$ the volume $`\mu _X(𝒥_{\mathrm{}}(Y))`$ is zero.
###### Proof.
<sup>35</sup><sup>35</sup>35In \[2, Proposition 6.22\] Batyrev claims that with the previous results one can reduce the proof to the case that $`Y`$ is a smooth divisor, where it is easily verified – we discussed this on page 21. Unfortunately I was not able to follow Batyrevs argument, thus the somewhat not so self contained proof is included here. For another another proof see Proposition 5.7.
The proof of this result relies on a fundamental result of Greenberg which, for our purpose is best phrased as follows:
###### Proposition 4.6.
Let $`Y`$ be a variety. Then there exists a positive integer $`c1`$. such that
$$\pi _{\mathrm{}{\scriptscriptstyle \frac{m}{c}}\mathrm{}}𝒥_{\mathrm{}}(Y)=\pi _{\mathrm{}{\scriptscriptstyle \frac{m}{c}}\mathrm{}}^m𝒥_m(Y)$$
for all $`m0`$.
This, in particular, implies that the image $`\pi _n(𝒥_{\mathrm{}}(Y))`$ is a constructible subset of $`𝒥_n(Y)`$. It can be shown that its dimension is just the expected one, namely $`dim\pi _n(𝒥_{\mathrm{}}(Y))=(n+1)dimY`$, see \[13, Lemma 4.3\]. From these observations we obtain a bound for the dimension of $`𝒥_m(Y)`$ as follows. We can work locally and may assume that $`YX`$ for a smooth $`X`$. Then we have
$$\begin{array}{cc}\hfill dim𝒥_m(Y)& dim(\pi _{\mathrm{}{\scriptscriptstyle \frac{m}{c}}\mathrm{}}^m𝒥_m(Y))+(m\mathrm{}\frac{m}{c}\mathrm{})dimX\hfill \\ & =dim(\pi _{\mathrm{}{\scriptscriptstyle \frac{m}{c}}\mathrm{}}𝒥_{\mathrm{}}(Y))+(m\mathrm{}\frac{m}{c}\mathrm{})dimX\hfill \\ & =(\mathrm{}\frac{m}{c}\mathrm{}+1)dimY+(m+1)dimX(\mathrm{}\frac{m}{c}\mathrm{}+1)dimX\hfill \\ & =(m+1)dimX(\mathrm{}\frac{m}{c}\mathrm{}+1)(dimXdimY).\hfill \end{array}$$
Thus, If the codimension of $`Y`$ in $`X`$ is greater or equal to 1, then $`dim𝒥_m(Y)𝕃^{mdimX}`$ approaches $`\mathrm{}`$ as $`m`$ approaches $`\mathrm{}`$. This implies that $`𝒥_{\mathrm{}}(Y)`$ is measurable that its measure $`\mu _X(𝒥_{\mathrm{}}(Y))`$ is zero. ∎
#### 4.1.1. Comparison with Lesbeque integration
To guide ones intuition a comparison of the motivic measure with more classical measures such as the Lesbeque measure or $`p`$–adic measures is sometimes helpful. We discuss here the similarities with the Lesbeque measure on $`𝔸^n`$ since this is wellknown.<sup>36</sup><sup>36</sup>36The analogy with $`p`$–adic measures is even more striking, see for example the discussions in and .
For convenience we consider the case that $`X=𝔸^n`$ in which case we identify the $`k`$–points of $`𝒥_{\mathrm{}}(𝔸_k^n)`$ with $`n`$–tuples of power series with coefficients in $`k`$. That is we identify $`𝒥_{\mathrm{}}(𝔸_k^n)(kt)^n=𝔸_{kt}^n`$. Since $`kt`$ is a discrete valuation domain (DVR) we can use the valuation to define a norm on $`(kt)^n`$ by defining
$$\tau \frac{1}{m}i:\tau _i(t^m)kt$$
for $`\tau =(\tau _1,\mathrm{},\tau _n)`$ a tuple of power series in $`𝒥_{\mathrm{}}(𝔸_k^n)`$. It is easy to check that this defines a (non–archimedian) norm. In Table 1 the similarities between Lesbeques and motivic measure are summarized.
### 4.2. Motivic integration on singular varieties
In the preceding discussion we used the assumption that your spaces are non-singular in several places in an essential way. Roughly speaking we used the fact that if $`X`$ is smooth then every cylinder is a stable set, thus can be endowed with a measure in a natural way. The fact that the resulting algebra of measurable sets to include the cylinders was essential for the setup since the level sets of the functions $`\mathrm{ord}_Y`$ are cylinders in a natural way.
If $`X`$ is singular however, many things one might got used to from the smooth case fail. Most prominently, the truncation maps are no longer surjective and a cylinder is in general not stable. Thus one has to work somewhat harder to obtain an algebra of measurable sets which include the cylinders. In order to be able to setup an integration theory one expects from this algebra of measurable sets the properties asserted in the following Proposition.
###### Proposition 4.7.
Let $`X`$ be a $`k`$–variety. Then there is an algebra of measurable subsets of $`𝒥_{\mathrm{}}(X)`$ and a measure $`\mu _X`$ on that algebra such that
1. If $`A`$ is stable, then $`A`$ is measurable and $`\mu _X(A)=[\pi _mA]𝕃^{nm}`$ for $`m0`$.
2. A cylinder $`C`$ is measurable and $`\mu _X(C)=lim_m[\pi _mC]𝕃^{nm}`$.
3. The measure is additive on finite disjoint unions.
4. If $`AB`$ are measurable, then $`dim\mu _X(A)dim\mu _X(B)`$.
To achieve this, one starts with the stable sets to which we know how to assign a measure. Then one proceeds just as in the smooth case using Definition 4.1, and replacing “cylinder” by “stable set” whenever necessary the same proof holds as well. The critical point now is to show that a cylinder is measurable with volume as claimed above. Even for the cylinder $`𝒥_{\mathrm{}}(X)=\pi ^1(X)`$ it is not clear a priori that it is measurable and what it’s measure should be (in fact, the measure of $`X`$ to be constructed leads to new birational invariants of $`X`$).
The point is that one has to partition $`𝒥_{\mathrm{}}(X)`$ according to intersection with the singular locus $`\mathrm{Sing}X`$, defined by the $`n`$th Fitting ideal of $`\mathrm{\Omega }_X^1`$. Then we can write
$$𝒥_{\mathrm{}}(X)=\underset{e0}{}𝒥_{\mathrm{}}^{()}(𝒳)$$
where $`𝒥_{\mathrm{}}^{()}(𝒳)=\mathrm{ord}_{\mathrm{Sing}𝒳}^\mathcal{1}()`$. As it turns out, the sets $`𝒥_{\mathrm{}}^{()}(𝒳)`$ are in fact stable at level $`e`$. The method is analogous to partitioning the cylinders of $`X^{}`$ according to intersection with $`K_{X^{}/X}`$ (defined by the $`0`$th Fitting ideal of $`\omega _{X^{}/X}`$) in the proof of the transformation rule. If one treats everything subordinate to this partition according to intersection with the singular locus one can construct a working theory in the singular case, see .
## 5. Birational invariants via motivic integration
As an illustration of the theory we discuss some applications of geometric motivic integration to birational geometry, namely we give a description of the log canonical threshold of a pair $`(X,Y)`$, where $`Y`$ is a closed subscheme of the smooth scheme $`X`$, in terms of the asymptotic behavior of the dimensions of the jet schemes $`𝒥_m(Y)`$. These results are due to Mustaţǎ and his collaborators . The precise statement is as follows:
###### Theorem 5.1.
Let $`YX`$ be a subscheme of the smooth variety $`X`$. Then the log canonical threshold of the pair $`(X,Y)`$ is
$$c(X,Y)=dimX\underset{m}{sup}\left\{\frac{dim𝒥_m(Y)}{m+1}\right\}.$$
The definition of the log canonical threshold requires the introduction of some more notation from birational geometry which will be done shortly. It is an invariant which can be read of from the data of a log resolution of the pair $`(X,Y)`$. The proof of the above is a very typical application of motivic integration as its strategy is to express the quantity one is interested in (say $`dim𝒥_m(Y)`$), in terms of a motivic integral. Then one uses the transformation rule (Theorem 3.1) to reduce to an integral over a normal crossing divisor which can be explicitly computed, similarly as Formula 2.6.
This result is in line with the earliest investigations of jet spaces by Nash who conjectured an intimate correspondence between the jet spaces of a singular space and the divisors appearing in a resolution of singularities. Even though his conjecture was disproved recently by Kollar and Ishii in general, there are important cases where his prediction was true, for example in the case of toric varieties.
### 5.1. Notation from birational geometry
Throughout this section we fix the following setup. Let $`X`$ be a smooth $`k`$–variety of dimension $`n`$ and let $`k`$ be of characteristic zero. Let $`YX`$ be a closed subscheme. Let $`f:X^{}\stackrel{}{}X`$ be a log resolution of the pair $`(X,Y)`$, that is a proper birational map such that
1. $`X^{}`$ is smooth and
2. denoting by $`F\stackrel{\mathrm{def}}{=}f^1Y=_{i=1}^sa_iD`$ and $`K\stackrel{\mathrm{def}}{=}K_{X^{}/X}=_{i=1}^sb_iD`$ for some $`a_i,b_i`$ and prime divisors $`D_i`$, the divisors $`F`$, $`K`$ and $`F+K`$ have simple normal crossing support.
The existence of log resolutions is a consequence of Hironaka’s resolution of singularities. Now one defines:
###### Definition 5.2.
Let $`X`$ and $`Y`$ and the datum of a log resolution be as above and let $`q0`$ be a rational number. Then we say that
1. $`(X,qY)`$ is *Kawamata log terminal* (KLT) if and only if $`b_iqa_i+1>0`$ for all $`i`$.
2. $`(X,qY)`$ is *log canonical* (LC) if and only if $`b_iqa_i+10`$ for all $`i`$.
We point out (without proof) that these notions are independent of the chosen log resolution and therefore well defined, see for details.
###### Remark 5.3.
These notions can be expressed in terms of the multiplier ideal $`(𝓆_𝒴)`$ of $`I_Y`$, the sheaf of ideals which cuts out $`Y`$ on $`X`$, as follows:
$`(X,qY)\text{ is KLT }`$ $`(_𝒴^𝓆)=𝒪_𝒳`$
$`(X,qY)\text{ is LC }`$ $`(_𝒴^𝓆)=𝒪_𝒳𝓆^{}<𝓆.`$
To see this observe that by definition
$$(_𝒴^𝓆)=𝒻_{}𝒪_𝒳^{}(\mathrm{}𝒦𝓆\mathrm{})=𝒻_{}𝒪_𝒳^{}(\mathrm{}𝒷_𝒾𝓆𝒶_𝒾\mathrm{}𝒟_𝒾)$$
which is equal to $`𝒪_𝒳`$ if and only if $`\mathrm{}b_iqa_i\mathrm{}0`$ for all $`i`$. as the upper corners denote the round up of an integer, this is equivalent to $`b_iqa_i+1>0`$ for all $`i`$ as required.
Now we proceed to the definition of the log canonical threshold, which is just the largest $`q`$ such that the pair $`(X,qY)`$ is Kawamata log terminal.
###### Definition 5.4.
The *log canonical threshold* of the pair $`(X,Y)`$ is
$$\begin{array}{cc}\hfill \mathrm{lct}(X,Y)& =sup\{q|(X,qY)\text{ is KLT }\}\hfill \\ & =sup\{q|b_iqa_i+1>0i\}\hfill \\ & =\underset{i}{\mathrm{min}}\left\{\frac{b_i+1}{a_i}\right\}\hfill \end{array}$$
Note that clearly one has $`\mathrm{lct}(X,qY)=q^1\mathrm{lct}(X,Y)`$ so that we can restrict to the case $`q=1`$ in the definition of the log canonical threshold. The formula for the log canonical threshold in terms of the jet schemes which we are aiming to proof in this section is an immediate consequence of the following theorem.
###### Theorem 5.5.
Let $`X`$ and $`Y`$ be as before. Then
$`(X,qY)\text{ is KLT}`$ $``$ $`dim𝒥_m(Y)<(m+1)(nq)\text{ for all }m`$
$`(X,qY)\text{ is LC}`$ $``$ $`dim𝒥_m(Y)(m+1)(nq)\text{ for all }m.`$
From this the proof of Theorem 5.1 follows immediately.
###### Proof of Theorem 5.1.
By definition we have
$$\begin{array}{cc}\hfill \mathrm{lct}(X,Y)& =sup\{q|(X,qY)\text{ is KLT }\}\hfill \\ & =sup\{q|dim𝒥_m(Y)<(m+1)(nq)m\}\hfill \\ & =n\underset{m}{sup}\left\{\frac{dim𝒥_m(Y)}{m+1}\right\}\hfill \end{array}$$
Moreover, the proof of Theorem 5.5 will reveal that the supremum in Theorem 5.5 is actually obtained by infinitely many $`m`$, namely whenever $`m+1`$ is divisible by all the coefficients $`a_i`$ of $`D_i`$ in $`f^1Y=a_iD_i`$ one has $`\mathrm{lct}(X,Y)=dimX\frac{dim𝒥_m(Y)}{m+1}`$.
Before proceeding to the proof we derive some elementary properties of the log canonical threshold.
###### Proposition 5.6.
Let $`YY^{}X`$ closed subvarieties of $`X`$. Then
1. $`\mathrm{lct}(X,Y)\mathrm{lct}(X,Y^{})`$.
2. $`0<\mathrm{lct}(X,Y)\mathrm{codim}(Y,X)`$ with equality if $`(X,Y)`$ is log canonical.
3. $`\mathrm{lct}(X,Y)`$ is independent of $`X`$ of fixed dimension.
4. Let $`(X^{},Y^{})`$ be another pair, then $`\mathrm{lct}(X\times X^{},Y\times Y^{})=\mathrm{lct}(X,Y)+\mathrm{lct}(X^{},Y^{})`$.
###### Proof.
For (1) note that $`YY^{}`$ implies that $`dim𝒥_m(Y)dim𝒥_m(Y^{})`$ then apply Theorem 5.5.
For (2) recall that in any case $`𝒥_m(Y)`$ contains $`𝒥_m(Y_{\text{reg}})`$ the jet scheme over the regular locus of $`Y`$. The latter has dimension $`(m+1)dimY`$. Therefore $`dim𝒥_m(Y)(m+1)dimY`$ and we finish by applying Theorem 5.5.
(3) is immediate since in the formula for $`\mathrm{lct}(X,Y)`$ the only feature of $`X`$ that appears is its dimension.
The formation of jet schemes preserves products<sup>37</sup><sup>37</sup>37By definition the functor $`𝒥_m(\underset{¯}{})`$ has a left adjoint thus commutes with direct products. and hence it follows that $`dim𝒥_m(Y\times Y^{})=dim𝒥_m(Y)+dim𝒥_m(Y^{})`$ from which (4) is implied immediately. ∎
### 5.2. Proof of threshold formula
Now we present the proof of Theorem 5.5. For this we first recall the transformation rule in a slightly more general form than stated above. Let $`A𝒥_{\mathrm{}}(X)`$ be a measurable subset (for example a stable subset) and let $`g:𝒥_{\mathrm{}}(X)\stackrel{}{}`$ be a function whose level sets are measurable. Then
$$_A𝕃^g𝑑\mu _X=_{f_{\mathrm{}}^1(A)}𝕃^{gf_{\mathrm{}}\mathrm{ord}_{K_{X^{}/X}}}𝑑\mu _X^{}$$
provided that one of the integrals exists, which then implies the existence of the other. The new features are minor. Clearly, it is allowed to integrate only over a measurable subset $`A`$ as long as we also only integrate over its measurable image $`f^{\mathrm{}}(A)`$ as well (that this image is measurable can be deduced from the proof of the transformation rule). In order to make the expression $`𝕃^q`$ for rational $`q`$ defined we have to adjoin roots of $`𝕃`$ to the already huge ring $`_𝓀`$ and define the dimension of $`𝕃^q`$ as $`q`$. For the following application it is in fact enough to adjoin a single root of $`𝕃`$ such that the new value ring of the integral is $`\widehat{}_k[𝕃^{1/n}]`$ for a sufficiently big $`n`$.
We apply this result with $`g=q\mathrm{ord}_Y=\mathrm{ord}_{qY}`$ and $`A=\mathrm{ord}_Y^1(m+1)`$ so that $`f_{\mathrm{}}^1(A)=\mathrm{ord}_F^1(m+1)`$ (up to measure zero) and $`gf_{\mathrm{}}\mathrm{ord}_{K_{X^{}/X}}=\mathrm{ord}_{K_{X^{}/X}qF}`$. Thus we get:
(5)
$$_{\mathrm{ord}_Y^1(m+1)}𝕃^{q\mathrm{ord}_Y}𝑑\mu _X=_{\mathrm{ord}_F^1(m+1)}𝕃^{\mathrm{ord}_{K_{X^{}/X}qF}}𝑑\mu _X^{}$$
Now, the left hand side of this equation contains information about the dimension of the $`m`$th jet scheme $`J_m(Y)`$, whereas the right hand side allows us to express this dimension in terms of the data of the log resolution. Together this will lead to a proof of Theorem 5.5.
Recall from Section 2.4 that
$$\mathrm{ord}_Y^1(m+1)=\pi _m^1(𝒥_m(Y))\pi _{m+1}^1(𝒥_{m+1}(Y))$$
and thus for the measure we have
$$\mu _X(\mathrm{ord}_Y^1(m+1))=[𝒥_m(Y)]𝕃^{nm}[𝒥_{m+1}(Y)]𝕃^{n(m+1)}.$$
Now computing the left hand side of equation (5) one gets
$$\begin{array}{cc}\hfill S_m& \stackrel{\mathrm{def}}{=}_{\mathrm{ord}_Y(m+1)}𝕃^{q\mathrm{ord}_Y}𝑑\mu _X=\mu _X(\mathrm{ord}_Y^1(m+1))𝕃^{q(m+1)}\hfill \\ & =([𝒥_m(Y)][𝒥_{m+1}(Y)]𝕃^n)𝕃^{nm+q(m+1)}\hfill \end{array}$$
Recalling that
(6)
$$dim[𝒥_m(Y)]+ndim[𝒥_{m+1}(Y)]$$
this implies, while evaluating at the dimension (see section 2.1), that
(A)
$$dimS_mdim𝒥_m(Y)nm+q(m+1)$$
with “$`<`$” holding only if we have equality in (6). This is a consequence of the property of the dimension function which says that $`dim(A+B)\mathrm{max}\{dimA,dimB\}`$ with equality as soon as $`dimAdimB`$.
Now we turn to computing the right hand side of equation (5):
$$\begin{array}{cc}\hfill _{\mathrm{ord}_F^1(m+1)}𝕃^{\mathrm{ord}_{KqF}}𝑑\mu _X^{}& =\underset{i=1}{\overset{\mathrm{}}{}}\mu _X^{}(\mathrm{ord}_F^1(m+1)\mathrm{ord}_{KqF}^1(i))𝕃^i\hfill \\ & =\underset{rA_m}{}\mu _X^{}(\underset{i}{}\mathrm{ord}_{D_i}^1(r_i))𝕃^{{\scriptscriptstyle (b_iqa_i)r_i}}\hfill \end{array}$$
where the last equality relies on a partitioning of $`\mathrm{ord}_F^1(m+1)`$, according to intersection with each component $`D_i`$ of the occurring normal crossing divisors:
$$A_m=\{r=(r_1,\mathrm{},r_s)|r_i0,a_ir_i=m+1\}\text{(this ensures }\mathrm{ord}_F=m+1\text{)}$$
$$\mathrm{ord}_F^1(m+1)=\underset{rA_m}{}(\underset{i}{}\mathrm{ord}_{D_i}^1(r_i))$$
Clearly, this refines the partition $`\mathrm{ord}_F^1(m+1)=_i(\mathrm{ord}_F^1(m+1)\mathrm{ord}_{KqF}^1(i))`$ and we have for $`\gamma \mathrm{ord}_{D_i}^1(r_i)`$ that $`\mathrm{ord}_{KqF}(\gamma )=(b_iqa_i)r_i`$ which justifies the above computation. The point now is that $`\mu _X^{}(_i\mathrm{ord}_{D_i}^1(r_i))`$ was computed explicitly in Lemma 2.7 to be equal to $`[D_{\mathrm{supp}r}^0](𝕃1)^{|\mathrm{supp}r|}𝕃^{{\scriptscriptstyle r_i}}`$. Hence
$$\begin{array}{cc}\hfill S_m& =_{\mathrm{ord}_F^1(m+1)}𝕃^{\mathrm{ord}_{KqF}}𝑑\mu _X^{}\hfill \\ & =\underset{rA_m}{}[D_{\mathrm{supp}r}^0](𝕃1)^{|\mathrm{supp}r|}𝕃^{{\scriptscriptstyle (b_iqa_i+1)r_i}}\hfill \end{array}$$
Note that the dimension of each summand is equal to $`n|\mathrm{supp}r|+|\mathrm{supp}r|(b_iqa_i+1)r_i=n(b_iqa_i+1)r_i`$, and since the coefficient of the highest dimensional part of each summand has a positive sign, there is no cancellation of highest dimensional parts in the sum. Thus we get
(B)
$$dimS_m=\underset{rA_m}{\mathrm{max}}\{n(b_iqa_i+1)r_i\}$$
With the formulas (A) and (B) at hand the proof of Theorem 5.5 follows easily.
###### Proof of Theorem 5.5.
It is enough to prove the first equivalence of Theorem 5.5, the second one being a limiting case of the first. Slightly reformulating and using the definition of KLT we have to show that
$$b_iqa_i+1>0idim𝒥_m(Y)nm+q(m+1)<nm$$
Let us first treat the implication “$``$”: In fact, we only need to assume the right hand side for one $`m_0`$ such that $`m_0+1`$ is divisible by each $`a_i`$. Then using equation (A) we have $`dimS_{m_0}dim𝒥_{m_0}(Y)nm_0+q(m_0+1)<n`$. Since $`b_iqa_i+1>0`$ holds trivially for $`a_i=0`$ we only need to consider $`i`$ such that $`a_i0`$. In this case (for fixed $`i`$) define the tuple $`r=(r_1,\mathrm{},r_s)`$ by setting $`r_i=\frac{m_0+1}{a_i}`$ and $`r_j=0`$ otherwise. Clearly $`rA_{m_0}`$. Now, equation (B) implies
$$n>dimS_{m_0}n\underset{j=1}{\overset{s}{}}(b_jqa_j+1)r_j=n(b_iqa_i+1)r_i$$
which says nothing but that $`b_iqa_i+1>0`$ as required.
Now we treat the converse “$``$”: Assuming that the pair $`(X,qY)`$ is KLT (i.e. $`dimS_m<n`$ for all $`m`$ by equation (B)) and that
(7)
$$dim𝒥_m(Y)nm+q(m+1)n$$
for some $`m`$ we seek a contradiction. These two assumptions together imply that the inequality (A) is strict, that is we have $`dimS_m<dim𝒥_m(Y)nm+q((m+1)`$. This, as we argued before, happens only if
$$dim𝒥_m(Y)=dim𝒥_{m+1}(Y)n.$$
Substituting this last equality into (7) we get $`(\text{7})`$ for $`m`$ replaced by $`m+1`$. Repeating this we obtain
$$dim𝒥_{m+i}(Y)=dim𝒥_m(Y)+in$$
for all $`i`$ which contradicts Proposition 5.7 below. ∎
We include here a more elementary and more explicit version of 4.5. This is also due to Mustaţǎ, we only sketch the proof here and refer the reader to \[35, Lemma 3.7\].
###### Proposition 5.7.
Let $`X`$ be smooth of dimension $`n`$ and let $`YX`$ be a subvariety. Let $`a=\mathrm{mult}_yY`$ be the local multiplicity of $`Y`$ at the point $`yY`$, then
$$dim(\pi _Y^m)^1(y)mdimX\mathrm{}\frac{m}{a}\mathrm{}.$$
Thus if $`a`$ is now the maximum of all local multiplicities of $`Y`$ one has
$$dim𝒥_m(Y)dimY+mdimX\mathrm{}\frac{m}{a}\mathrm{}(m+1)dimX\mathrm{}\frac{m}{a}\mathrm{}.$$
It follows that if $`Y`$ is nowhere dense in $`X`$, then $`\mu _X(𝒥_{\mathrm{}}(Y))=0`$.
###### Proof.
The second statement clearly follows from the first which in turn immediately reduces to the case that $`YX`$ is a hypersurface. Since $`X`$ is smooth it is étale over $`𝔸^n`$ with $`y`$ mapping to $`0`$. Since $`(\pi _X^{\mathrm{}})^1(y)`$ gets thereby identified with $`(\pi _{𝔸^n}^{\mathrm{}})^1(0)`$ we may assume that $`Y𝔸^n`$ is given by the vanishing of $`fk[x]=k[x_1,\mathrm{},x_n]`$ and the point $`y`$ is the origin.
The condition that the local multiplicity of $`Y`$ at $`0`$ is equal to $`a`$ means that the smallest degree monomial of $`f`$ has degree $`a`$ in $`x_0,\mathrm{},x_n`$. For simplicity we assume now that $`f`$ is homogeneous of degree $`a`$, in general one can combine the following proof with a deformation argument to reduce to this case along the way, see for this general case.
Exercise 2.5 states that $`𝒥_m(Y)𝒥_m(𝔸^n)`$ is given by $`m+1`$ equations $`f^{(0)},\mathrm{},f^{(m)}`$ in the coordinates of $`𝒥_m(𝔸^n)`$ described in Example 2.1. Concretely, $`f^{(i)}k[x^{(0)},\mathrm{},x^{(i)}]`$ is given as the coefficient of $`t^i`$ in the power series
$$f(\underset{i}{}x_1^{(i)}t^i,\mathrm{},\underset{i}{}x_n^{(i)}t^i).$$
With the notation $`f_0^{(i)}\stackrel{\mathrm{def}}{=}f^{(i)}(0,x^{(1)},x^{(2)},\mathrm{},x^{(i)})`$ (we abbreviated the tuples $`x_1^{(i)},\mathrm{},x_n^{(i)}`$ by $`x^{(i)}`$) the fiber $`(\pi _Y^m)^1(0)`$ is given by the vanishing of $`f_0^{(1)},\mathrm{},f_0^{(m)}`$ in the fiber $`(\pi _X^m)^1(0)\mathrm{Spec}k[x^{(1)},\mathrm{},x^{(m)}]`$. Recall that we need to show that the dimension of $`(\pi _Y^m)^1(0)`$ is at most $`mn\mathrm{}\frac{m}{a}\mathrm{}`$. Thus it is enough to show that the dimension of the variety given by the vanishing of the ideal $`I_p=(f_0^{(a)},f_0^{(2a)},\mathrm{},f_0^{(pa)})`$ is at most $`(pa)np`$.
To show this we turn to an initial ideal of $`I_p`$ with respect to the degree reverse lexicographic order on $`k[x^{(1)},\mathrm{},x^{(pa)}]`$ where the underlying ordering of the variables $`x_i^{(j)}`$ is first by upper index and then by lower.<sup>38</sup><sup>38</sup>38That is we order the variables according to $`x_i^{(j)}x_i^{}^{(j^{})}`$ iff $`j>j^{}`$ or $`j=j^{}`$ and $`i<i^{}`$. The degree reverse lexicographic order on a polynomial ring $`k[x_1,\mathrm{},x_n]`$ with $`x_1>x_2>\mathrm{}>x_n`$ is given by
$$A=x_1^{a_1}\mathrm{}x_n^{a_n}>x_1^{b_1}\mathrm{}x_n^{b_n}=B$$
if $`\mathrm{deg}A>\mathrm{deg}B`$ or $`\mathrm{deg}A=\mathrm{deg}B`$ and $`a_i<b_i`$ for the last index $`k`$ for which $`a_ib_i`$. Roughly speaking, a monomial in degree rev lex is big if it contains fewest of the cheap variables. For example $`x_2^4>x_1^3x_3>x_1^2x_2x_3`$. Consult \[18, Chapter 15\] for details. Now it is a matter of unravelling the definitions<sup>39</sup><sup>39</sup>39Show that for two monomials $`A<B`$ in $`k[x]`$ one has $`\mathrm{in}_<(A_0^{(j)})<\mathrm{in}_<(B_0^{(j)})`$. This reduces to the case that $`f=\mathrm{in}_<(f)`$ is a monomial. Then observe that the monomial $`f(x^{(j)})`$ appears in $`f^{(ja)}`$. Since the rev-lex-ordering is such that a term is large if it contains fewest variables which are small it follows that $`f(x^{(j)})`$ is in fact the leading term as claimed. (of the $`f^{(i)}`$, of the order …) to see that
$$\mathrm{in}_<(f_0^{(aj)})=\mathrm{in}_<(f)(x^{(j)})\text{ for }j=1,\mathrm{},p.$$
As these are $`p`$ many nontrivial equations in disjoint variables it follows that the dimension of the vanishing locus of the initial ideal of $`I_p`$ is at most $`(pa)np`$. Thus the same upper bound holds for the dimension of the vanishing locus of $`I_p`$ itself, and a forteriori for the dimension of $`(\pi _Y^{pa})^1(0)`$. ∎
### 5.3. Bounds for the log canonical threshold
In this section we show how the just derived description of the log canonical threshold in terms of the dimension of the jet spaces lead to some interesting bounds for $`\mathrm{lct}(X,Y)`$.
###### Proposition 5.8.
Let $`a`$ be the maximal local multiplicity of a point in $`Y`$. Then
$$\frac{1}{a}\mathrm{lct}(X,Y)\frac{dimX}{a}.$$
###### Proof.
We may assume that $`Y`$ is neither $`\mathrm{}`$ nor all of $`X`$ since these cases are trivial. Let $`p`$ be a point with maximal multiplicity $`a`$. The second part of Exercise 5.1 shows that $`dim𝒥_{a1}(Y)dim(\pi ^{a1})^1(p)dimX(a1)`$. Hence by Theorem 5.1
$$\mathrm{lct}(X,Y)dimX\frac{dim𝒥_{a1}(Y)}{a}dimXdimX\frac{a1}{a}=\frac{dimX}{a}$$
For the lower bound Proposition 5.7 gives $`dim𝒥_m(Y)dimY+mdimX\mathrm{}\frac{m}{a}\mathrm{}`$ and again using Theorem 5.1 this yields
$$\mathrm{lct}(X,Y)\frac{dimXdimY}{m+1}+\frac{\mathrm{}\frac{m}{a}\mathrm{}}{m+1}$$
for all $`m`$. For sufficiently divisible and sufficiently big $`m`$ this implies that $`\mathrm{lct}(X,Y)\frac{1}{a}`$ as claimed. ∎
The next application of the arc space techniques is a bound for the log canonical threshold of a homogeneous hypersurface. This uses the following exercise as a key ingredient.
###### Exercise 5.1.
<sup>40</sup><sup>40</sup>40This exercise is solved by explicitly writing down the equations which define $`(\pi _Y^m)^1(0)`$ within $`(\pi _{𝔸^n}^m)^1(0)`$. These turn out to be the same as the defining equations of the right hand side; just the variables are different.
Let $`Y𝔸^n`$ be a homogeneous hypersurface of degree $`d`$. Show that one has an isomorphism
$$(\pi _Y^m)^1(0)𝒥_{md}(Y)\times 𝔸^{n(d1)}$$
for all $`md1`$, where we set $`𝒥_1(Y)`$ to be a point.
Drop the assumption *homogeneous* and assume instead that the local multiplicity of $`Y`$ at $`p`$ be equal to $`a`$. Show that $`(\pi _Y^{a1})^1(p)𝔸^{n(a1)}`$.
###### Proposition 5.9.
Let $`Y𝔸^n`$ be a homogeneous hypersurface of degree $`d`$. Then
$$\mathrm{lct}(𝔸^n,Y)\mathrm{min}\{\frac{nr}{d},1\}$$
where $`r=dim\mathrm{Sing}Y`$.
###### Proof.
One key ingredient is the observation of the previous Exercise 5.1 that for $`md1`$
$$(\pi _Y^m)^1(0)𝒥_{md}(Y)\times 𝔸^{n(d1)}$$
By semicontinuity of $`dim(\pi _Y^m)^1(p)`$ the inequality
$$dim(\pi _Y^m)^1(p)dim𝒥_{md}(Y)+n(d1)$$
holds for all $`p`$, and in particular for the $`p\mathrm{Sing}Y`$. Hence all together we get the estimate
$$dim𝒥_m(Y)\mathrm{max}\{dim𝒥_{md}(Y)+n(d1),(n1)(m+1)\}$$
where $`(n1)(m+1)`$ is equal to $`dim𝒥_m(Y\mathrm{Sing}Y)`$ which is always a lower bound for $`dim𝒥_m(Y)`$. Now set $`m=pd1`$ and apply the inequality repeatedly to get
$$dim𝒥_{pd1}(Y)\mathrm{max}\{(ndn+r)p,(n1)pd\}$$
which amounts to
$$\mathrm{lct}(𝔸^n,Y)=n\frac{dim𝒥_{pd1}(Y)}{pd}\mathrm{min}\{\frac{nr}{d},1\}$$
since the log canonical threshold is computed via the dimension of the jet spaces $`𝒥_{pd1}(Y)`$ for sufficiently divisible $`pd`$. ∎
In the main achievement is to characterize the extremal case as follows: In the case that $`\mathrm{lct}(𝔸^n,Y)1`$ one has $`\mathrm{lct}(𝔸^n,Y)=\frac{nr}{d}`$ if and only if $`YY^{}\times 𝔸^r`$ for some $`Y^{}𝔸^{nr}`$ a hypersurface.
### 5.4. Inversion of adjunction
One of the most celebrated applications of motivic integration to birational geometry is a much improved understanding of the Inversion of Adjunction conjecture of Shokurov and Kollar . The conjecture describes how certain invariants of singularities of pairs behave under restriction. The following proposition goes in this direction as it shows that under restriction to a smooth hypersurface the log canonical threshold can only decrease, that is the singularities cannot get better under restriction.
###### Proposition 5.10.
Let $`(X,Y)`$ be a pair and $`H`$ a smooth hypersurface in $`X`$. Then
$$c_H(X,Y)\mathrm{lct}(H,YH)$$
where $`c_H(X,Y)`$ is the log canonical threshold of the pair $`(X,Y)`$ *around* $`H`$, that is the minimum $`\mathrm{lct}(U,YU)`$ over all open $`UX`$ with $`UH\mathrm{}`$.
###### Proof.
The proof uses a straightforward extension of the formula for the log canonical threshold to include this more general case of $`c_H(X,Y)`$. In fact the same proof as above shows that one has
$$c_H(X,Y)=dimX\underset{m}{sup}\left\{\frac{dim_H𝒥_m(Y)}{m+1}\right\}$$
with equality for sufficiently divisible $`m+1`$. Here $`dim_H𝒥_m(Y)`$ is the *dimension of $`𝒥_m(Y)`$ along $`H`$*, that is the maximal dimension of an irreducible component $`T`$ of $`𝒥_m(Y)`$ such that $`\pi ^m(T)H\mathrm{}`$.
Let $`T`$ be an irreducible component of $`𝒥_m(Y)`$ such that $`\pi ^m(T)H\mathrm{}`$. This implies that $`T𝒥_m(YH)`$ is also nonempty since the projection $`\pi _{YH}^m`$ is surjective. Since $`HX`$ is locally given by one equation the same is true for $`HYY`$. Hence $`𝒥_m(HY)𝒥_m(Y)`$ is given by at most $`m+1`$ equations (see Exercise 2.5). Hence $`dim_H𝒥_m(Y)dim𝒥_m(YH)+(m+1)`$ which shows
$$\underset{m}{sup}\left\{\frac{dim_H𝒥_m(Y)}{m+1}\right\}\underset{m}{sup}\left\{\frac{dim_H𝒥_m(YH)}{m+1}\right\}+1$$
which implies the claimed inequality $`c_H(X,Y)\mathrm{lct}(H,YH)`$. ∎
The *Inversion of Adjunction* Conjecture of Kollár and Shokurov describes how the singularities of pair behave under restriction to a Cartier divisor. More precisely, let $`(X,Y)`$ be a pair where we allow $`Y=q_iY_i`$ to be any formal integer combination (rational or real combination even) of closed subvarieties of $`X`$. With the notation $`f:X^{}\stackrel{}{}X`$ of a log resolution as above (in particular $`f^1Y=a_iE_i`$ and $`K_{X^{}/X}=b_iE_i`$) and a subvariety $`WX`$ fixed we define the *minimal log discrepancy*
$$\mathrm{mld}(W;X,Y)\stackrel{\mathrm{def}}{=}\{\begin{array}{cc}\mathrm{min}\{b_ia_i+1|f(E_i)W\}\hfill & \text{if this minimum is non-negative}\hfill \\ \mathrm{}\hfill & \text{otherwise.}\hfill \end{array}$$
It follows that $`(X,Y)`$ is log canonical on an open subset containing $`W`$ iff $`\mathrm{mld}(W;X,Y)\mathrm{}`$. The inversion of adjunction conjecture now states:
###### Conjecture 5.11.
With $`(X,Y)`$ and as above, let $`D`$ be a normal effective Cartier divisor on $`X`$ such that $`DY`$ and let $`WD`$ a proper closed subset. Then we have
$$\mathrm{mld}(W;X,Y+D)=\mathrm{mld}(W;D,Y|_D).$$
The inequality “$``$” is the *adjunction* part and is well known to follow from the adjunction formula $`K_D=(K_X+D)|_D`$. The reverse inequality “$``$” is the critical part of this conjecture.
In the conjecture was proved in the case that $`X`$ is smooth and $`Y`$ is effective. In this was established even for $`X`$ a complete intersection (and $`Y`$ effective). The proof of these results use a description of the minimal log discrepancies in terms of dimensions of certain cylinders of the jet spaces of $`X`$, analogous to the one for the log canonical threshold.
###### Proposition 5.12.
With the notation as above and for $`X`$ smooth and $`Y`$ effective,
$$\mathrm{mld}(W;X,Y)\tau $$
$$$$
$$\mathrm{codim}_{𝒥_{\mathrm{}}(X)}(\mathrm{Cont}_Y^\nu \pi _0^1(W))q_i\nu _i+\tau \text{ for all multi-indices }\nu \text{.}$$
where $`\mathrm{Cont}_Y^\nu =_i\mathrm{ord}_{Y_i}^1(\nu _i)`$.
The proof of this is not more than a technical complication of the proof of the log canonical threshold formula we gave above. With this characterization of $`\mathrm{mld}(W;X,Y)`$ the proof of inversion of adjunction becomes a matter of determining the co-dimensions of the cylinders involve. In the case that $`X`$ and $`D`$ are both smooth this is quite easy (and could be done as an exercise). In general ($`X`$ a complete intersection) the combinatorics involved can become quite intricate, cf. .
One should point out that after these results on inversion of adjunction were obtained, Takagi found an alternative approach using positive characteristic methods.
### 5.5. Geometry of arc spaces without explicit motivic integration.
As it should have become apparent by now, the applications of motivic integration to birational geometry are by means of describing certain properties of a variety $`X`$ in terms of (mostly simpler) properties of its jet spaces $`𝒥_m(X)`$. Motivic integration serves as the path to make this connection. However, due to some combinatiorial difficulties one encounters along this path, one can ask if there is a more direct relationship. This is indeed the case and it is the content of the paper of of Ein, Lazarsfeld and Mustaţǎ which is the source of the material in this section.
Their point is that instead of using the birational transformation rule to control the dimension of components of the jet spaces one uses the Key Theorem 3.3 of its proof to more directly get to the desired information.
We keep the notation of a subvariety $`YX`$ of a smooth variety $`X`$ and define.
###### Definition 5.13.
Let $`Y`$ be a subvariety of $`X`$ and $`p0`$ an integer define the *contact locus* to be the cylinder
$$\mathrm{Cont}_Y^p\stackrel{\mathrm{def}}{=}\mathrm{ord}_Y^1(p)𝒥_{\mathrm{}}(X)$$
The aim is to understand the components (or at least the dimension) of the cylinders $`\mathrm{Cont}_Y^p`$ in terms of the data coming from a log resolution of the pair $`(X,Y)`$. Using the notation of Section 5.1 we fix a log resolution $`f:X^{}\stackrel{}{}X`$ of the pair $`(X,Y)`$ and denote $`f^1Y=_1^ka_iE_i`$ and $`K_{X^{}/X}=_1^kb_iE_i`$ where the support of $`E_i`$ is a simple normal crossing divisor.
###### Definition 5.14.
Given $`E=_1^kE_i`$, a simple normal crossing divisor of $`X^{}`$, and a multi-index $`\nu =(\nu _1,\mathrm{},\nu _k)`$ define the *multi contact locus*
$$\mathrm{Cont}_E^\nu \stackrel{\mathrm{def}}{=}\left\{\gamma ^{}𝒥_{\mathrm{}}(X^{})\right|\mathrm{ord}_{E_i}(\gamma ^{})=\nu _i\text{ for }i=1\mathrm{}k\}.$$
###### Definition 5.15.
For every cylinder $`C𝒥_{\mathrm{}}(X)`$ there is a well defined notion of codimension, namely
$$\mathrm{codim}C\stackrel{\mathrm{def}}{=}\mathrm{codim}(𝒥_m(X),\pi _mC)$$
for $`m0`$.
Of course one must check that this is independent of the chosen $`m0`$. This however is immediately clear from the definition of cylinder.
###### Exercise 5.2.
For $`C𝒥_{\mathrm{}}(X)`$ a cylinder, show that $`\mathrm{codim}C=dimXdim\mu _X(C)`$.
The following proposition replaces in this new setup the computation of the motivic integral of a normal crossing divisor in Proposition 2.6. Note the comparative simplicity!
###### Proposition 5.16.
For $`E=_1^kE_i`$ a simple normal crossing divisor $`\mathrm{Cont}_E^\nu `$ is a smooth irreducible cylinder of codimension
$$\mathrm{codim}\mathrm{Cont}_E^\nu =\underset{1}{\overset{k}{}}\nu _i$$
provided $`\mathrm{Cont}_E^\nu `$ is nonempty.
###### Proof.
This is a computation in local coordinates. Assume $`E`$ is locally given by the vanishing of the first $`k`$ of the coordinates $`x_1=\mathrm{}=x_k=0`$. Then, for an arc $`\gamma `$, which is determined by $`\gamma (x_i)=\gamma _i^{(j)}t^j`$ for $`i=1\mathrm{}n`$, to have the prescribed contact order with the $`E_i`$ means precisely that $`\gamma _i^{(j)}=0`$ for $`j<\nu _i`$, and $`\gamma _i^{(\nu _i)}0`$. Hence for $`m0`$ we have
$$\pi _m(\mathrm{Cont}_E^\nu )(𝔸^1\{0\})^n\times 𝔸^{m\nu _i}$$
and therefore $`\mathrm{Cont}_E^\nu `$ is smooth irreducible and of codimension $`\nu _i`$. ∎
The central result (replacing the transformation rule) is the following Theorem
###### Theorem 5.17.
With the notation as above one has for all $`p>0`$ a finite partition
$$\mathrm{Cont}_Y^p=\underset{v}{}f_{\mathrm{}}\mathrm{Cont}_E^n$$
where the disjoint union is over all multi-indices $`\nu `$ such that $`\nu _ia_i=p`$.
For every multi-index is the set $`f_{\mathrm{}}\mathrm{Cont}_E^v`$ an irreducible cylinder of codimension
$$\nu _i(b_i+1).$$
In particular, for each irreducible component $`Z`$ of $`\mathrm{Cont}_Y^p`$ there is a unique multi-index $`\nu `$ such that $`\mathrm{Cont}_E^\nu `$ dominates $`Z`$.
###### Proof.
As in the proof of the transformation rule the key ingredient is Theorem 3.3. With this at hand the proof is not difficult.
The condition $`\nu _ia_i=p`$ ensures that $`f_{\mathrm{}}\mathrm{Cont}_E^\nu \mathrm{Cont}_Y^p`$. The surjectivity of $`f_{\mathrm{}}`$ (cf. Exercise 3.2) on the other hand implies the reverse inclusion. The disjoined-ness of the union follows from the fact that there is a one-to-one map between
$$𝒥_{\mathrm{}}(X^{})𝒥_{\mathrm{}}(E)\stackrel{11}{}𝒥_{\mathrm{}}(X)𝒥_{\mathrm{}}(Y)$$
(which is an implication of the valuative criterion for properness as explained in Section 3.0.1) and the observation that each $`\mathrm{Cont}_E^\nu `$ is contained in the left hand side.
Since $`\mathrm{Cont}_E^\nu \mathrm{Cont}_{K_{X^{}/X}}^{{\scriptscriptstyle b_iv_i}}=\mathrm{ord}_{K_{X^{}/X}}^1`$ it follows from Theorem 3.3 (a) and Proposition 3.2 that $`f_{\mathrm{}}\mathrm{Cont}_E^\nu `$ is a cylinder. Part (b) of Theorem 3.3 shows that
$$f_m:\mathrm{Cont}_E^\nu \stackrel{}{}f_m\mathrm{Cont}_E^\nu $$
is a piecewise trivial $`𝔸^{{\scriptscriptstyle b_i\nu _i}}`$–fibration. By Exercise 5.2 is the codimension of $`\mathrm{Cont}_E^\nu `$ equal to $`\nu _i`$. Hence the codimension of its image under $`f_{\mathrm{}}`$ is
$$\mathrm{codim}(f_{\mathrm{}}\mathrm{Cont}_E^\nu )=\nu _i+\nu _ib_i=\nu _i(b_i+1)$$
as claimed. ∎
As an illustration of this result we recover a very clean proof of the log canonical threshold formula of Theorem 5.1.
###### Proof of Theorem 5.1.
Let $`V_m`$ be an irreducible component of $`𝒥_m(Y)`$. For some $`pm+1`$ the set $`\mathrm{Cont}_Y^p(\pi _m^X)^1V_m`$ is open in $`(\pi _m^X)^1V_m`$, namely for the smallest $`p`$ (automatically $`m+1`$ since each arc in $`V_m`$ has contact order $`m+1`$ with $`Y`$) such that $`\mathrm{Cont}_Y^p(\pi _m^X)^1V_m\mathrm{}`$. Hence there is an irreducible component $`W`$ of $`\mathrm{Cont}_Y^p`$ such that the closure $`\overline{W}`$ contains $`(\pi _m^X)^1V_m`$. By Theorem 5.17 there is a unique multi-index $`\nu `$ (necessarily $`\nu _ib_i=p`$) such that $`\mathrm{Cont}_E^\nu `$ dominates $`W`$.
By definition of the log canonical threshold we have $`b_i+1\mathrm{lct}(X,Y)a_i`$ for all $`i`$ such that we obtain the following inequalities:
$$\begin{array}{cc}\hfill \mathrm{codim}(V_m,𝒥_m(X))& \mathrm{codim}W\hfill \\ & =\mathrm{codim}f_{\mathrm{}}\mathrm{Cont}_E^\nu \hfill \\ & \nu _i(b_i+1)\hfill \\ & \nu _i\mathrm{lct}(X,Y)a_i\hfill \\ & =\mathrm{lct}(X,Y)p=\mathrm{lct}(X,Y)(m+1)\hfill \end{array}$$
As this holds for every irreducible component of $`𝒥_m(Y)`$ we get
$$\mathrm{lct}(X,Y)\frac{\mathrm{codim}(𝒥_m(Y),𝒥_m(X))}{m+1}.$$
To see that there is equality for some $`m`$ we pick an index $`i`$ such that $`\mathrm{lct}(X,Y)=\frac{b_i+1}{a_i}`$ and $`m+1`$ divisible by $`a_i`$. Let $`\nu `$ be the multi-index which is zero everywhere except at the $`i`$th spot, where it is $`\frac{m+1}{a_i}`$. Then $`f_{\mathrm{}}\mathrm{Cont}_E^\nu \mathrm{Cont}_Y^{m+1}(\pi _m^X)^1𝒥_m(Y)`$ and by Theorem 5.17 the codimension $`f_{\mathrm{}}\mathrm{Cont}_E^\nu `$ is equal to $`\frac{m+1}{a_i}(b_i+1)=\mathrm{lct}(X,Y)(m+1)`$. Hence in particular $`\mathrm{codim}(𝒥_m(Y),𝒥_m(X))\mathrm{lct}(X,Y)(m+1)`$ for this chosen $`m+1`$. This finishes the argument. ∎
In summary, the above agument shows that the irreducible components $`V`$ of $`𝒥_m(Y)`$ of maximal possible dimension, that is the ones that compute the log canonical theshold as $`\mathrm{lct}(X,Y)=\mathrm{codim}(𝒥_m(X),V)`$ are dominated by multi-contact loci $`\mathrm{Cont}_E^\nu `$ with $`\nu _i0`$ for all the indices $`i`$ such that $`E_i`$ computes the the log canonical theshold (meaning $`\mathrm{lct}(X,Y)=\frac{b_i+1}{a_i}`$).
We want to finish these notes with Mustaţǎ’s characterization of rational singularities for complete intersections in terms of arc spaces. This was indeed the first application of motivic integration to characterizing singularities. With the just developed viewpoint this result is not too difficult anymore.
###### Theorem 5.18.
Let $`YX`$ be a reduced and irreducible locally complete intersection subvariety of codimension $`c`$. Then the jet spaces $`𝒥_m(Y)`$ are irreducible for all $`m`$ if and only if $`Y`$ has rational singularities.
###### Proof.
Let $`f:X^{}\stackrel{}{}X`$ be a log resolution of $`(X,Y)`$ which dominates the blowup of $`X`$ along $`Y`$. Keeping the previous notation we may assume that $`E_1`$ is the exceptional divisor of this blowup. In Theorem 2.1 it is shown that $`Y`$ has at worst rational singularities if and only if $`b_ica_i`$ for every $`i2`$. Hence we must show
$$𝒥_m(Y)\text{ is irreducible for all }m1b_ica_i\text{ for }i2$$
Assume that $`𝒥_m(Y)`$ is not irreducible, that is we have a component $`V𝒥_m(Y)`$ other than the main component $`𝒥_m(Y\mathrm{Sing}Y)`$. As in the previous proof we have $`W\mathrm{Cont}_Y^p`$ with $`pm+1`$ whose closure contains $`\pi _m^1(V)`$. By Theorem 5.17 this component is dominated by some multi-contact locus $`\mathrm{Cont}_E^\nu `$ for $`\nu (m+1,0,\mathrm{},0)`$ since the latter is the multi-index corresponding to the multi-contact locus dominating $`\pi _m^1(𝒥_m(Y\mathrm{Sing}Y))`$. Since $`YX`$ is a local complete intersection of codimension $`c`$ we have $`\mathrm{codim}(V,𝒥_m(X))(m+1)c`$. To arrive at a contradiction assume now that $`Y`$ has rational singularities, that is assume that $`b_ica_i`$ for $`i2`$. Then
$$\begin{array}{cc}\hfill (m+1)c& \mathrm{codim}(W)\hfill \\ & =\nu _1c\underset{i2}{}\nu _i(b_i+1)\hfill \\ & c\underset{i1}{}\nu _ia_i+\underset{2}{}\nu _i\text{ (since }b_ica_i\text{)}\hfill \\ & =cp+\underset{i2}{}\nu _ic(m+1)+\underset{i2}{}\nu _i\hfill \end{array}$$
Hence for $`i2`$ we must have $`\nu _i=0`$, a contradiction.
Conversely, suppose $`b_i<ca_i`$ for some $`i2`$. Setting $`v`$ to be the multi-index with all entries zero except the $`i`$th equal to $`1`$. Let $`(m+1)=a_i`$, then the $`\mathrm{Cont}_E^\nu `$ maps to an irreducible subset $`W\mathrm{Cont}_Y^{m+1}`$ of codimension $`(m+1)c`$. Hence $`\pi _{m+1}(W)`$ is an irreducible component of $`𝒥_m(Y)`$ of codimension $`(m+1)c`$ which is not the component $`\overline{𝒥_m(Y\mathrm{Sing}Y)}`$. Hence $`𝒥_m(Y)`$ is not irreducible. ∎
## Appendix A An elementary proof of the Transformation rule.
We present Looijenga’s elementary proof of Theorem 3.3 which then leads to a proof of the transformation formula avoiding weak factorization. For this we have to investigate more carefully the definition of the relative canonical divisor $`K_{X^{}/X}`$ and suitably interpret the contact multiplicity of an arc $`\gamma `$ with $`K_{X/X^{}}`$.
### A.1. The relative canonical divisor and differentials
Let us consider the first fundamental exact sequence for Kähler differentials, as it plays a pivotal role in all that follows.
The morphism $`f:X^{}\stackrel{}{}X`$ induces a linear map, its derivative, $`f^{}\mathrm{\Omega }_X\stackrel{df}{}\mathrm{\Omega }_X^{}`$ which is part of the first fundamental exact sequence for Kähler differentials:
(8)
$$0\stackrel{}{}f^{}\mathrm{\Omega }_X\stackrel{df}{}\mathrm{\Omega }_X^{}\stackrel{}{}\mathrm{\Omega }_{X^{}/X}\stackrel{}{}0$$
Note that by our assumption of smoothness, the $`𝒪_𝒳^{}`$–modules $`f^{}\mathrm{\Omega }_X`$ and $`\mathrm{\Omega }_X^{}`$ are locally free of rank $`n=dimX`$. Since, by birationality of $`f`$, $`\mathrm{\Omega }_{X^{}/X}`$ has rank zero, the first map is injective as well. Taking the $`n`$th exterior power we obtain the map
$$0\stackrel{}{}f^{}\mathrm{\Omega }_X^n\stackrel{^ndf}{}\mathrm{\Omega }_X^{}^n$$
of locally free $`𝒪_𝒳^{}`$ modules of rank $`1`$. If we set $`\omega =\mathrm{\Omega }^n`$ and tensor the above sequence with the invertible sheaf $`\omega _X^{}^1`$ we obtain
$$f^{}\omega _X\omega _X^{}^1𝒪_𝒳^{}$$
thus identifying $`f^{}\omega _X\omega _X^{}^1`$ with a locally principal ideal in $`𝒪_𝒳^{}`$, which we shall denote by $`J_{X^{}/X}`$ (so, by definition, $`J_{X^{}/X}`$ is the 0-th Fitting ideal of $`\mathrm{\Omega }_{X^{}/X}`$). Now define $`K_{X^{}/X}`$ to be the Cartier divisor which is locally given by the vanishing of $`J_{X^{}/X}`$. It is important to note that $`K_{X^{}/X}`$ is defined as an effective divisor and not just as a divisor class. By choosing bases for the free $`𝒪_𝒳^{}`$-modules $`f^{}\mathrm{\Omega }_X`$ and $`\mathrm{\Omega }_X^{}`$ the map $`df`$ is given by a $`n\times n`$ matrix with entries in $`𝒪_𝒳^{}`$. Its determinant is a local defining equation for $`K_{X^{}/X}`$.
Let $`L`$ be an extension field of $`k`$ and let $`\gamma :\mathrm{Spec}Lt\stackrel{}{}X`$ be a $`L`$-rational point of $`𝒥_{\mathrm{}}(X^{})`$, and assume that $`\mathrm{ord}_{K_{X^{}/X}}(\gamma )=e`$. By definition of contact order, this means that $`(t^e)=\gamma ^{}(J_{X^{}/X})Lt`$. As $`J_{X^{}/X}`$ is locally generated by $`detdf𝒪_𝒳^{}`$ (well defined up to unit) we obtain that $`(t^e)=det(\gamma ^{}(df))`$. The pullback of the sequence (8) along $`\gamma `$ illustrates the situation:
(9)
$$0\stackrel{}{}(f\gamma )^{}\mathrm{\Omega }_X\stackrel{\gamma ^{}df}{}\gamma ^{}\mathrm{\Omega }_X^{}\stackrel{}{}\gamma ^{}\mathrm{\Omega }_{X^{}/X}\stackrel{}{}0$$
Since $`Lt`$ is a PID, we can choose bases of $`(f\gamma )^{}\mathrm{\Omega }_X`$ and $`\gamma ^{}\mathrm{\Omega }_X^{}`$ such that $`\gamma ^{}(df)`$, a map of free $`Lt`$ modules of rank $`n`$, is given by a diagonal matrix. With respect to this basis the exact sequence (8) takes the form
(10)
$$0\stackrel{}{}Lt^n\stackrel{\left(\begin{array}{ccc}t^{e_1}& & 0\\ & \mathrm{}& \\ 0& & t^{e_n}\end{array}\right)}{}Lt^n\stackrel{}{}\frac{Lt}{(t^{e_i})}\stackrel{}{}0$$
The condition that $`\mathrm{ord}_{K_{X^{}/X}}(\gamma )=e`$ translates into $`_{i=1}^ne_i=e`$ or, equivalently, into saying that the rightmost module is torsion of length $`e`$.
### A.2. Proof of Theorem 3.3
We start by recalling the statement of Theorem 3.3 we want to proof slightly reformulated in order to set up the notation that is used in its proof below.
###### Theorem A.1.
Let $`f:X^{}\stackrel{}{}X`$ be a proper birational morphism of smooth $`k`$-varieties. Let $`C_e^{}=\mathrm{ord}_{K_{X^{}/X}}^1(e)`$ where $`K_{X^{}/X}`$ is the relative canonical divisor and let $`C_e\stackrel{\mathrm{def}}{=}f_{\mathrm{}}C_e^{}`$. Let $`\gamma C_e^{}`$ an $`L`$-point of $`𝒥_{\mathrm{}}(𝒳^{})`$, that is a map $`\gamma ^{}:𝒪_𝒳^{}\stackrel{}{}𝓉`$, satisfying $`\gamma ^{}(J_{X^{}/X})=(t^e)`$, with $`Lk`$ a field extension. Then for $`m2e`$ one has:
1. For all $`\xi 𝒥_{\mathrm{}}(𝒳)`$ such that $`\pi _m^X(\xi )=f_m(\pi _m^X^{}(\gamma ))`$ there is $`\gamma ^{}𝒥_{\mathrm{}}(𝒳^{})`$ such that $`f_{\mathrm{}}(\gamma ^{})=\xi `$ and $`\pi _{me}^X^{}(\gamma ^{})=\pi _{me}^X^{}(\gamma )`$. In particular, the fiber of $`f_m`$ over $`f_m(\gamma _m)`$ lies in the fiber of $`\pi _{me}^m`$ over $`\gamma _{me}`$.
2. $`\pi _m(C_e^{})`$ is a union of fibers of $`f_m`$.
3. The map $`f_m:\pi _m^X^{}(C_e^{})\stackrel{}{}C_e`$ is a piecewise trivial $`𝔸^e`$ fibration.
###### Proof.
To ease notation we will denote truncation by lower index, i.e. write $`\gamma _m`$ as shorthand for $`\pi _m^X(\gamma )`$. We already pointed out before that (a’) implies (a):
The proof of (b) can be divided into two steps. First we show that the fiber of $`f_m`$ over $`f_m(\gamma _m)`$ can be naturally identified with $`\mathrm{Der}_{𝒪_𝒳}(𝒪_𝒳^{},\frac{𝓉}{(𝓉^{𝓂+\mathcal{1}})})`$. Then we have to show that the latter is an affine space of dimension $`e`$. As this is easy let’s do it first: Immediately preceding this proposition we noted that the cokernel of $`\gamma ^{}(df)`$ is torsion of length $`e`$ as a $`Lt`$-module. This cokernel is $`\gamma ^{}\mathrm{\Omega }_{X^{}/X}`$. Since $`m>e`$ the dual, $`\mathrm{Hom}_{Lt}(\gamma ^{}\mathrm{\Omega }_{X^{}/X},\frac{Lt}{(t^{m+1})})`$, is also torsion of length $`e`$. Using adjointness of $`\gamma ^{}`$ and $`\gamma _{}`$ this $`\mathrm{Hom}`$ is just $`\mathrm{Hom}_{𝒪_𝒳^{}}(\mathrm{\Omega }_{X^{}/X},\gamma _{}\frac{Lt}{(t^{m+1})})`$, which is equal to $`\mathrm{Der}_{𝒪_𝒳}(𝒪_𝒳^{},\gamma _{}\frac{𝓉}{(𝓉^{𝓂+\mathcal{1}})})`$ essentially by definition of $`\mathrm{\Omega }_{X^{}/X}`$. This shows that $`\mathrm{Der}_{𝒪_𝒳}(𝒪_𝒳^{},\gamma _{}\frac{𝓉}{(𝓉^{𝓂+\mathcal{1}})})`$ is isomorphic to $`𝔸_L^e`$. Thus we are left to show the identification $`(\mathrm{}\mathrm{}\mathrm{})`$ of the following diagram the last line of which is the first exact sequence for derivations, analogous to the above exact sequence of Kähler differentials.
The identification $`(\mathrm{})`$ is given by sending $`\gamma _m^{}`$ to $`\gamma _m^{}\gamma _m`$ which, since $`m2e`$, can easily<sup>41</sup><sup>41</sup>41Fix an homomorphism $`\gamma :R\stackrel{}{}S`$ of $`k`$-algebras which makes $`S`$ into an $`R`$-algebra. For any ideal $`IS`$ with $`I^2=0`$ one has a map (11)
$$\{\gamma ^{}\mathrm{Hom}_{k\text{alg}}(R,S)|\mathrm{Im}(\gamma ^{}\gamma )I\}\stackrel{}{}\mathrm{Der}_k(R,I)$$
by sending $`\gamma ^{}`$ to $`\gamma ^{}\gamma `$. To check that $`(\gamma ^{}\gamma )`$ is indeed a derivation one has to make the following calculation verifying the Leibniz rule (note that the $`R`$ algebra structure on $`S`$ is via $`\gamma `$):
$$\begin{array}{cc}\hfill (\gamma ^{}\gamma )(xy)& ((\gamma ^{}\gamma )(x)\gamma (y)+\gamma (x)(\gamma ^{}\gamma )(y))\hfill \\ & =\gamma ^{}(x)\gamma ^{}(y)\gamma (x)\gamma (y)\gamma ^{}(x)\gamma (y)+\gamma (x)\gamma (y)\gamma (x)\gamma ^{}(y)+\gamma (x)\gamma (y)\hfill \\ & =\gamma ^{}(x)(\gamma ^{}(y)\gamma (y))\gamma (x)(\gamma ^{}(y)\gamma (y))\hfill \\ & =(\gamma ^{}\gamma )(x)(\gamma ^{}\gamma )(y)=0\hfill \end{array}$$
The last line is zero by the assumption that $`\mathrm{Im}(\gamma ^{}\gamma )I`$ and $`I^2=0`$. The obvious inverse map sending a derivation $`\delta `$ to $`\gamma +\delta `$ shows that the two sets in (11) are equal. This setup clearly applies in our situation: $`R=𝒪_𝒳^{}`$, $`S=Lt/t^{m+1}`$, $`I=(t^{m+1e})`$ where $`m2e`$ ensures that $`I^2=0`$. be checked to define an $`L`$-derivation $`𝒪_𝒳^{}\stackrel{}{}\frac{(𝓉^{𝓂+\mathcal{1}})}{(𝓉^{𝓂+\mathcal{1}})}`$. In this way (and using (a’)) we think of $`f_m^1(f_m(\gamma _m))`$ as a subspace of $`\mathrm{Der}_L(𝒪_𝒳^{},\frac{(𝓉^{𝓂+\mathcal{1}})}{(𝓉^{𝓂+\mathcal{1}})})`$. As this is the $`t^e`$-torsion part of $`\mathrm{Der}_L(𝒪_𝒳^{},\frac{𝓉}{(𝓉^{𝓂+\mathcal{1}})})`$ and since we just observed that $`\mathrm{Der}_{𝒪_𝒳}(𝒪_𝒳^{},\frac{𝓉}{(𝓉^{𝓂+\mathcal{1}})})`$ is torsion of lenght $`e`$ the inclusion $`(\mathrm{}\mathrm{})`$ is also clear, and thus $`(\mathrm{}\mathrm{}\mathrm{})`$ becomes a statement about subsets of $`\mathrm{Der}_L(𝒪_𝒳^{},\frac{(𝓉^{𝓂+\mathcal{1}})}{(𝓉^{𝓂+\mathcal{1}})})`$.
Let $`(\gamma _m^{}\gamma _m)\mathrm{Der}_L(𝒪_𝒳^{},\frac{(𝓉^{𝓂+\mathcal{1}})}{(𝓉^{𝓂+\mathcal{1}})})`$. The image of $`(\gamma _m^{}\gamma _m)`$ in $`\mathrm{Der}_L(f^{}𝒪_𝒳,\frac{𝓉}{(𝓉^{𝓂+\mathcal{1}})})`$ is $`\gamma _m^{}f\gamma _mf`$. This is zero (i.e. $`\gamma _m^{}\mathrm{Der}_{𝒪_𝒳}(𝒪_𝒳^{},\frac{𝓉}{(𝓉^{𝓂+\mathcal{1}})})`$) if and only if $`f_m(\gamma _m^{})=f_m(\gamma _m)`$, that is if and only if $`\gamma _m^{}`$ is in the fiber of $`f_m`$ over $`f_m(\gamma _m)`$. This concludes the proof of (b).
In order to come by the element $`\gamma ^{}𝒥_{\mathrm{}}(X^{})`$ as claimed in (a’) we construct a sequence of arcs $`\gamma ^k𝒥_{\mathrm{}}(X^{})`$ satisfying the following two properties for all $`km`$:
1. $`\pi _k(f_{\mathrm{}}(\gamma ^k))=\pi _k(\xi )`$ and
2. $`\pi _{k1e}(\gamma ^k)=\pi _{k1e}(\gamma ^{k1})`$ and $`\pi _{me}(\gamma ^k)=\pi _{me}(\gamma )`$.
Clearly, setting $`\gamma ^{m1}=\gamma ^m=\gamma `$ these conditions hold for $`k=m`$. Furthermore, the second condition implies that the limit $`\gamma ^{}\stackrel{\mathrm{def}}{=}lim_k\gamma ^k`$ exists and that $`\pi _{me}(\gamma ^{})=\pi _{me}(\gamma )`$. The first condition shows that $`f_{\mathrm{}}(\gamma ^{})=\xi `$. Thus we are left with constructing the sequence $`\gamma ^k`$. This is done inductively. As we already verified the solution for $`k=m`$ we assume to have $`\gamma ^k`$ and $`\gamma ^{k1}`$ as claimed – now $`\gamma ^{k+1}`$ is constructed as follows:
Since $`\pi _k(f(\gamma ^k))=\pi _k(\xi )`$ we can view their difference as a derivation $`\delta =\xi f\gamma ^k\mathrm{Der}_L(𝒪_𝒳,\frac{(𝓉^{𝓀+\mathcal{1}})}{(𝓉^{𝓀+\mathcal{2}})})`$ which we identify with $`\mathrm{Hom}_{Lt}(\gamma ^kf^{}\mathrm{\Omega }_X,\frac{(t^{k+1})}{(t^{k+2})})`$. The latter module appears in $`\mathrm{Hom}_{Lt}(\underset{¯}{},\frac{Lt}{(t^{k+2})})`$ applied to the sequence (9), where $`\gamma ^k`$ takes the place of $`\gamma `$:
(12)
In order to understand this better we turn to the same sequence, but with respect to the basis as in sequence (10), where it takes this form:
Now it becomes clear that $`\delta `$ lies in the image of $`df`$ since $`e`$ and therefore all $`e_i`$ are less than $`m+1k+1`$. Furthermore, any pre-image $`\delta ^{}`$ must lie in $`\mathrm{Hom}(\gamma ^k\mathrm{\Omega }_X^{},\frac{(t^{k+1e})}{(t^{k+2})})`$ by the shape of the matrix and the fact that $`e_ie`$ for all $`e`$. Now pick any such pre-image $`\delta ^{}`$ and define $`\gamma ^{k+1}\stackrel{\mathrm{def}}{=}\delta ^{}+\gamma ^k`$. This is an arc in $`X^{}`$ with $`\pi _{ke}(\gamma ^{k+1})=\pi _{ke}(\gamma ^k)`$. Furthermore since $`df(\delta ^{})=\delta `$ we get
$$\gamma ^{k+1}f\gamma ^kf=\delta =\xi \gamma ^kfmod(t^{k+2})$$
and thus $`\pi _{k+1}(\gamma ^{k+1}f)=\pi _{k+1}(f_{\mathrm{}}(\gamma ^{k+1}))=\pi _{k+1}(\xi )`$. ∎
With this proof of Theorem 3.3 at hand a proof of the Transformation rule follows immediately as indicated in Section 3. |
warning/0507/math0507378.html | ar5iv | text | # 1 Statement
## 1 Statement
Let $`V`$ be a real vector space of dimension $`d`$, and let $`f_1,\mathrm{},f_n`$ be a collection of nonconstant affine linear functions on $`V`$ such that the associated linear forms span the dual vector space $`V^{}`$. Let $`𝒜`$ denote the collection of affine hyperplanes $`H_1,\mathrm{},H_nV`$, where $`H_j=f_j^1(0)`$ for all $`j\{1,\mathrm{},n\}`$. Let
$$(𝒜):=V^{}\underset{j=1}{\overset{n}{}}H_j^{}$$
be the complement of the complexification of $`𝒜`$. Our goal is to introduce a new space which is which is naturally the base of a principal bundle with structure group $`V^{}`$ and total space $`(𝒜)`$. Note that a principal bundle with structure group equal to a vector space is the same as an affine bundle whose associated vector bundle is trivial.
Our new space, which we call $`𝒵(𝒜)`$, will be a real algebraic prevariety of dimension $`d`$, where “prevariety” means that $`𝒵(𝒜)`$ will not be Hausdorff in the analytic topology. Though this may sound nasty, our description will be quite simple. We construct $`𝒵(𝒜)`$ by gluing together a collection of vector spaces $`\{V_C\}`$, each isomorphic to $`V`$, where the index $`C`$ ranges over the chambers of $`𝒜`$, by which we mean the connected components of the complement of $`𝒜`$ in $`V`$. These vector spaces are attached to each other along open sets, according to the following rule:
$$V_CV_C^{}V\underset{\begin{array}{c}H_j\text{ separates}\\ C\text{ from }C^{}\end{array}}{}H_j.$$
In Figure 1 we illustrate the example of $`n`$ distinct points on a line; here $`𝒵(𝒜)`$ is a real line with $`n`$ double points.
More generally, $`𝒵(𝒜)`$ admits a natural map to $`V`$, and the number of points of $`𝒵(𝒜)`$ lying above a point $`qV`$ is equal to the number of chambers $`C`$ which contain $`q`$ in their closures.
In this particular example, it is clear that there is a map from $`(𝒜)`$ to $`𝒵(𝒜)`$ given by smooshing the imaginary axis, and that this map is a weak homotopy equivalence (see Remark 1.2). In the following theorem, which is our main result, we strengthen this observation, and generalize it to arbitrary arrangements.
###### Theorem 1.1
The complexified complement $`(𝒜)`$ may be realized as a principal bundle over $`𝒵(𝒜)`$ in the category of real analytic prevarieties, with structure group naturally isomorphic to the additive group of the dual vector space $`V^{}`$.
###### Remark 1.2
Since $`(𝒜)`$ is isomorphic to the total space of a locally trivial fiber bundle over $`𝒵(𝒜)`$ with contractible fibers, the projection from $`(𝒜)`$ to $`𝒵(𝒜)`$ is a weak homotopy equivalence, i.e. it induces isomorphisms on all homotopy and homology groups. This map is not an honest homotopy equivalence because it does not have a homotopy inverse; in particular, it admits no section.
###### Remark 1.3
Theorem 1.1 bears a strong similarity to the celebrated theorem of Salvetti \[Sa\] which exhibits a simplicial complex $`\mathrm{Sal}(𝒜)`$ that is homotopy equivalent to $`(𝒜)`$. Salvetti’s theorem offers some obvious advantages over Theorem 1.1 – the space in question is a simplicial complex, one gets a homotopy equivalence rather than a weak homotopy equivalence, and the Salvetti complex is an invariant of the oriented matroid associated to $`𝒜`$. On the other hand, Theorem 1.1 has some advantages of its own. First, the space $`𝒵(𝒜)`$ is quite easy to visualize, perhaps more so than $`\mathrm{Sal}(𝒜)`$. Second, we are able to work real analyticlly, rather than simply topologically. Finally, there is a canonical map from $`(𝒜)`$ to $`𝒵(𝒜)`$, while the homotopy equivalence in Salvetti’s theorem requires some arbitrary choices.
###### Remark 1.4
In the special case where the arrangement $`𝒜`$ is central and defined over the rational numbers, Theorem 1.1 may be interpreted in the world of hypertoric varieties, originally introduced by Bielawski and Dancer \[BD\]. A real hypertoric variety is a real variety of dimension $`2d`$ associated to a rational hyperplane arrangement of rank $`d`$, and it carries an action of the real algebraic gtorus $`(^\times )^d`$. Let $`𝒴(𝒜)`$ denote the open subset of this variety on which $`(^\times )^d`$ acts freely. One may show that $`𝒴(𝒜)/_2^d`$ is isomorphic to $`(𝒜)`$, and $`𝒴(𝒜)/(^\times )^d`$ is isomorphic to $`𝒵(𝒜)`$. Hence $`(𝒜)`$ is a principal bundle over $`𝒵(𝒜)`$ with structure group $`(^\times )^d/_2^d`$, which is analytically isomorphic to the additive group $`^{}`$. Though we will not use the language of hypertoric varieties in our proof of Theorem 1.1, those familiar with hypertoric varieties will observe that the proof is closely guided by this interpretation.
###### Remark 1.5
Since $`𝒵(𝒜)`$ comes with an open cover $`\{V_C\}`$, it is natural to consider the Mayer-Vietoris spectral sequence associated to this cover. All multiple intersections of open sets are complements of real hyperplane arrangements, and therefore have cohomology only in degree zero. It follows that the cohomology groups of $`𝒵(𝒜)`$ are equal to the homology groups of the complex $`E_1^{0,}`$, where $`E_1^{0,q}`$ is the direct sum of the $`0`$-th cohomology groups of all $`(q+1)`$-fold intersections. Theorem 1.1 tells us that the cohomology groups of $`𝒵(𝒜)`$ are isomorphic to those of $`(𝒜)`$, which are well understood by \[OS\]. It would be interesting to see the complex $`E_1^{0,}`$ arise in some independent combinatorial context.
###### Example 1.6
We conclude the first section with a detailed analysis of the simplest case, that of a single point on a line, which is the $`n=1`$ case in Figure 1. Here $`(𝒜)`$ is isomorphic to $`^\times `$, and $`𝒵(𝒜)`$ is a real line with a double point at the origin. We will write down explicitly the action of the additive group $``$ on $`(𝒜)`$, see that the orbit space is isomorphic to $`𝒵(𝒜)`$, and check that the projection from $`(𝒜)`$ to $`𝒵(𝒜)`$ is locally trivial. This will serve both as an illustration of Theorem 1.1, and as an important tool to apply toward the proof of the general case in Section 2.
First consider the action of $`^\times `$ on $`^2\{0\}`$ given by the formula
$$\lambda (x,y)=(\lambda x,\lambda ^1y).$$
(1)
The map from $`^2\{0\}`$ to $``$ taking $`(x,y)`$ to $`xy`$ is surjective and $`^\times `$-invariant. The fiber of this map over a nonzero number consists of a single $`^\times `$ orbit, while the fiber over zero consists of two orbits, namely the two coordinate axes. Hence the quotient of $`^2\{0\}`$ by $`^\times `$ is a line with a double point at the origin, which we will call $`𝔻`$. The restriction of the action to the complement of either of the two coordinate axes in $`^2`$ induces a trivial $`^\times `$ bundle over the quotient $``$, therefore $`^2\{0\}`$ is a principal $`^\times `$ bundle over $`𝔻`$, trivialized over the two copies of $``$ in $`𝔻`$.
Now consider the quotient of this entire picture by the subgroup $`_2^\times `$. We now obtain an action of
$$^\times /_2_+$$
on
$$\left(^2\{0\}\right)/_2^\times /_2^\times ,$$
where the isomorphism between $`_+`$ and $``$ is given by the logarithm, and the isomorphism between $`^\times /_2`$ and $`^\times `$ is given by the map taking $`\pm z`$ to $`iz^2`$. (We include the factor of $`i`$ for technical reasons to simplify notation in the proof of the general case of Theorem 1.1.) Thus we obtain $`^\times `$ as a principal $``$ bundle over $`𝔻`$. In fact, the structure group of the bundle is naturally dual to the original line, hence we will denote the structure group $`^{}`$.
For the sake of concreteness, let us write this action down in coordinates. To make clear the distinction between the action of $`^{}`$ on $`^\times `$ and the action of $`^\times `$ on $`^2\{0\}`$ given by Equation (1), we will denote the $`^{}`$ action with the symbol $``$. Given an element $`z^\times `$, choose a complex number $`x+iy`$ such that $`z=i(x+iy)^2`$. Then for $`\xi ^{}`$, we have
$$\xi z=i(e^\xi x+e^\xi yi)^2=2xy+i(e^{2\xi }x^2e^{2\xi }y^2)=\mathrm{Re}(z)+i(e^{2\xi }x^2e^{2\xi }y^2).$$
Hence the action of $`^{}`$ changes the imaginary part of $`z`$, and leaves the real part fixed. If the real part of $`z`$ is nonzero, then $`x`$ and $`y`$ are both nonzero, and the orbit $`^{}z`$ is equal to the vertical real line through $`z`$. On the other hand, if $`z`$ is purely imaginary, then either $`x`$ or $`y`$ is zero, and the orbit $`^{}z`$ is equal to the component of the punctured imaginary axis containing $`z`$.
###### Remark 1.7
It is easy to become confused by the changes of variables in Example 1.6. We will always adhere to the convention that when we use the letters $`x`$, $`y`$, and $`z`$, we have $`z=i(x+iy)^2`$. If we need letters to refer to the real and imaginary parts of $`z`$, as we will in Lemma 2.1, we will use the notation $`z=a+ib`$.
## 2 Proof
We now turn to the proof of the general case of Theorem 1.1. Consider the affine linear map
$$f:V^n$$
given by the functions $`f_i`$ which define the hyperplanes of $`𝒜`$. The complexification $`f^{}`$ of $`f`$ induces a closed embedding of $`(𝒜)`$ into $`(^\times )^n`$. By Example 1.6, $`(^\times )^n`$ is a principal bundle over $`𝔻^n`$ with structure group $`(^n)^{}`$. Let $`\pi `$ denote the projection from $`(^\times )^n`$ to $`𝔻^n`$.
###### Lemma 2.1
The image of $`(𝒜)`$ under $`\pi `$ is isomorphic to $`𝒵(𝒜)`$.
###### Proof.
Let us write $`𝔻=^\times \{p,m\}`$, where $`p`$ and $`m`$ stand for plus and minus. Then the projection from $`^\times `$ to $`𝔻`$ is given by sending $`a+ib`$ to $`a`$ if $`a`$ is nonzero, and otherwise to $`p`$ or $`m`$, depending on the sign of $`b`$.
Let $`\stackrel{~}{a}=(\stackrel{~}{a}_1,\mathrm{},\stackrel{~}{a}_n)`$ be a point of $`𝔻^n`$ lying over a point $`a=(a_1,\mathrm{},a_n)^n`$. Then $`\stackrel{~}{a}`$ lifts to an element $`a+ib`$ of $`(𝒜)`$ if and only if $`a`$ lies in the image of $`f`$, and there exists $`b=(b_1,\mathrm{},b_n)^n`$ satisfying the following conditions:
$`(i)`$ $`b+f(0)`$ is in the image of $`f`$
$`(ii)`$ if $`\stackrel{~}{a}_j=p`$, then $`b_j>0`$
$`(iii)`$ if $`\stackrel{~}{a}_j=m`$, then $`b_j<0`$.
Consider the map $`\stackrel{~}{f}_C:V_C𝔻^n`$ lifting the map $`f:V^n`$, defined by the property that if $`f_j(q)=0`$, then the $`j^{\text{th}}`$ coordinate of $`\stackrel{~}{f}_C(q)`$ is determined by the sign of $`f_j`$ on $`C`$. These maps glue together to define an inclusion
$$\stackrel{~}{f}:𝒵(𝒜)𝔻^n.$$
Furthermore, $`\stackrel{~}{a}`$ is in the image of $`f_C`$ if and only if $`a`$ is in the image of $`f`$ and $`C`$ is contained in the set
$$S(\stackrel{~}{a}):=\underset{\stackrel{~}{a}_j=p}{}f_j^1(_+)\underset{\stackrel{~}{a}_j=m}{}f_j^1(_{}).$$
It follows that $`\stackrel{~}{a}`$ is in the image of $`\stackrel{~}{f}`$ if and only if $`a`$ is in the image of $`f`$ and $`S(\stackrel{~}{a})`$ is nonempty. Suppose that $`a=f(q_0)`$ for some point $`q_0V`$. Then the map
$$qf(q)f(q_0)=f(qq_0)f(0)$$
takes $`S(\stackrel{~}{a})`$ bijectively onto the set of points $`b^n`$ satisfying Conditions $`(i),(ii),`$ and $`(iii)`$ above, hence $`S(\stackrel{~}{a})`$ is nonempty if and only if such a $`b`$ exists. This proves that the image of $`\stackrel{~}{f}`$ coincides with the image of $`\pi `$ restricted to $`(𝒜)`$. ∎
Let $`W(^n)^{}`$ be the set of linear forms on $`^n`$ which are constant on the affine subspace $`V`$. (Equivalently, $`W`$ is the kernel of the dual of the linearization of $`f`$.) Consider an element $`z(^\times )^n`$, and choose $`x,y^n`$ such that $`z_j=i(x+iy)^2`$ for all $`j\{1,\mathrm{},n\}`$. We define a function $`\rho _z:W`$ by the formula
$$\rho _z(\xi )=e^\xi (x,y)^2,$$
where $`e^\xi \mathrm{exp}(W)(^\times )^n`$ acts coordinatewise as in Equation (1), and the norm is the Euclidean norm on $`^{2n}`$. Given a nonzero element $`\xi W`$, we compute the directional derivative of $`\rho _z`$ along $`\xi `$:
$$\frac{d}{dt}\mathrm{exp}(t\xi )(x,y)^2=\underset{j=1}{\overset{n}{}}\frac{d}{dt}(e^{t\xi _j}x_j,e^{t\xi _j}y_j)^2=\underset{j=1}{\overset{n}{}}2\xi _j\left(e^{2t\xi _j}x_j^2e^{2t\xi _j}y_j^2\right).$$
(2)
Evaluating at $`t=0`$, we obtain
$$\frac{d}{dt}\rho _z(\mathrm{exp}(t\xi ))|_{t=0}=\mathrm{\hspace{0.17em}\hspace{0.17em}2}\xi ,\mathrm{Im}(z),$$
(3)
where $`,`$ is the natural pairing between $`(^n)^{}`$ and $`^n`$.
Let
$$𝒩(𝒜)=\pi ^1\left(\pi ((𝒜))\right)=(^n)^{}(𝒜)(^\times )^n.$$
By Lemma 2.1, $`𝒩(𝒜)`$ is a principal bundle over $`𝒵(𝒜)`$ with structure group $`(^n)^{}`$. We now prove a pair of lemmas that are analogous to the main theorem of \[KN\], which lays the groundwork for the equivalence of quotients in symplectic and algebraic geometry. This idea is central to the perspective alluded to in Remark 1.4.
###### Lemma 2.2
If $`z𝒩(𝒜)`$, then $`\xi `$ is a critical point of $`\rho _z`$ if and only if $`e^\xi z(𝒜)`$.
###### Proof.
Using the fact that $`\rho _z(\xi )=\rho _{\xi z}(\mathrm{id})`$, we may immediately reduce to the case where $`\xi =0`$. An element $`z`$ of $`𝒩(𝒜)`$ automatically has its real part contained in the image of $`f`$, hence it lies in $`(𝒜)`$ if and only if $`\mathrm{Im}(z)`$ lies in the image of the linearization of $`f`$. This in turn is the case if and only if the imaginary part of $`z`$ pairs trivially with every $`\xi W`$ (this is how $`W`$ is defined). Equation 3 tells us that this happens if and only if zero is a critical point of $`\rho _z`$. ∎
###### Lemma 2.3
If $`z𝒩(𝒜)`$, then $`\rho _z`$ has a unique critical point.
###### Proof.
By differentiating Equation 2 we see that $`\rho _z`$ is convex, hence any critical point must be unique. What remains is to prove existence.
Let us consider the behavior of the right hand side of Equation 2 as $`t`$ becomes very large. If $`\xi _j>0`$, then the $`j^{\text{th}}`$ term approaches positive infinity provided that $`x_j0`$, otherwise it remains bounded. Similarly, if $`\xi _j<0`$, then the $`j^{\text{th}}`$ term approaches positive infinity provided that $`y_j0`$, and is otherwise bounded. Hence the directional derivative of $`\rho _z`$ along $`\xi `$ at $`t\xi z`$ is positive for large $`t`$ provided that there exists an index $`j`$ for which either
$`\xi _j>0`$ and $`x_j0`$ or $`\xi _j<0`$ and $`y_j0`$. (4)
If this condition is satisfied, then when $`t`$ is large, the gradient of $`\rho _z`$ restricted to the sphere of radius $`t`$ always points outward, hence $`\rho _z`$ must have a critical point somewhere on $`\mathrm{exp}`$ of the ball of radius $`t`$. Hence it will suffice to prove that there exists an index $`j`$ satisfying Equation 4.
Recall that $`z_j=i(x_j+iy_j)^2`$. If $`x_j=0`$, then $`z_j`$ lies on the negative part of the imaginary axis; if $`y_j=0`$, then it lies on the positive part of the imaginary axis. Since $`z𝒩(𝒜)`$, there exists $`w(𝒜)`$ lying in the same $`(^n)^{}`$ orbit as $`z`$. In particular, we have $`\mathrm{Im}(w_j)<0`$ whenever $`x_j=0`$, and $`\mathrm{Im}(w_j)>0`$ whenever $`y_j=0`$. Suppose that Condition (4) fails for all $`j`$. Then $`\xi _j`$ and $`w_j`$ have opposite signs whenever $`\xi _j0`$, hence $`\xi ,w0`$. This contradicts the fact that $`\xi W`$ and $`w(𝒜)`$. ∎
Lemmas 2.2 and 2.3 combine to tell us that each $`W`$ orbit in $`𝒩(𝒜)`$ contains a unique element of $`(𝒜)`$. Thus
$$(𝒜)𝒩(𝒜)/W𝒩(𝒜)\times _{(^n)^{}}V^{},$$
which is the principal bundle over $`𝒵(𝒜)`$ induced from $`𝒩(𝒜)`$ by the surjection
$$(^n)^{}(^n)^{}/WV^{}.$$
This completes the proof of Theorem 1.1.
Acknowledgments. The author is grateful to Columbia University and the City of New York for their hospitality during the writing of this paper. |
warning/0507/hep-lat0507035.html | ar5iv | text | # DESY 05-124 HU-EP-05/31 SFB/CPP-05-32 Non–perturbative quark mass renormalization in two–flavor QCD
## 1 Introduction
Lattice QCD provides a definition of quark masses from first principles. For example given a hadronic input accessible to experiments, such as the K–meson mass $`m_\mathrm{K}`$ and decay constant $`F_\mathrm{K}`$, it is possible to compute the strange quark mass on the lattice. A convenient quantity to consider is the renormalization group invariant (RGI) strange quark mass, which is independent of the renormalization scheme if the renormalization conditions are imposed at zero quark mass \[?\].
If one looks at the most recent lattice results for the strange quark mass, from simulations with two \[?,?,?\] and two plus one \[?,?\] dynamical quark-flavours, there is a large spread of values ranging from 68$`\mathrm{MeV}`$ to 132$`\mathrm{MeV}`$. There are several sources of systematic errors in these computations, such as the use of perturbative renormalization (except for \[?\], where non–perturbative renormalization is done in the RI–MOM scheme) and the values of the lattice spacing affordable nowadays.
To perform a completely controlled computation we start from the bare PCAC mass $`m_i(g_0)`$ for a given quark flavor $`i`$ and determine the RGI mass
$$M_i=Z_\mathrm{M}(g_0)m_i(g_0).$$
(1.1)
The bare coupling $`g_0`$ is in a one–to–one relation to the lattice spacing and the continuum limit of $`M_i`$ exists (and should be taken). As an intermediate step we first define the running mass
$$\overline{m}_i(\mu )=\frac{Z_\mathrm{A}(g_0)}{Z_\mathrm{P}(g_0,a\mu )}m_i(g_0),$$
(1.2)
where $`\overline{m}`$ is non–perturbatively defined in the Schrödinger Functional renormalization scheme and hence is well–defined also for low energies $`\mu `$. This corresponds to splitting $`Z_\mathrm{M}`$ into two factors
$$Z_\mathrm{M}(g_0)=\frac{M}{\overline{m}(\mu )}\times \frac{Z_\mathrm{A}(g_0)}{Z_\mathrm{P}(g_0,a\mu )},$$
(1.3)
which will be computed with full non–perturbative precision following the strategy of \[?\]. Note that all renormalization factors are flavor–independent and we hence omit the subscript $`i`$ in mass ratios. In the first factor of the above splitting the universal continuum limit is understood to have been taken. Its computation is the main objective of this paper. This result can then be used for any action and only the second factor needs to be redetermined. This becomes tractable by making the universal factor available for $`\mu `$ in the range of hadronic energies of $`\mathrm{O}(12\mathrm{GeV})`$. As for the running coupling \[?\] our method thus avoids the need to treat a multi–scale problem in a large volume in this step.
First results on the $`\mu `$–dependence of $`\overline{m}`$ in the $`N_\mathrm{f}=2`$ theory have already appeared in Refs. \[?,?,?\]. Beyond a finalization of these results, we here compute the second factor — and hence $`Z_\mathrm{M}`$ — using non–perturbative improvement \[?\] and renormalization \[?\] of the axial current in the theory with two flavors of O($`a`$) improved Wilson quarks and plaquette gauge action.
As an application, one wants to determine the light (up, down, strange) quark masses.<sup>1</sup><sup>1</sup>1Later, when lattice spacings are small enough, the charm quark mass can be determined \[?,?\]. Also the beauty quark mass computed in HQET through the strategy of \[?\], is based on Eq. (1.1) in QCD. For practical reasons this is at present not yet possible for us in the most straight-forward way by simulations at the physical parameters. We need a hierarchy of additional approximations to compute for instance the strange quark mass.
First of all our simulation algorithm is at the moment still restricted to pairs of degenerate flavours and we include one such pair ($`N_\mathrm{f}=2`$). In the large volume simulations needed to determine $`m_i(g_0)`$ in Eq. (1.1) our masses can at present not be taken to values small enough for the up– and down–flavors, for instance by tuning the pseudoscalar states to the physical pion mass. Instead we shall determine the RGI quark mass $`M_{\mathrm{ref}}`$ that is associated with a “Kaon” made from two degenerate flavours.
The analogous computation of $`M_{\mathrm{ref}}`$ has been performed in the quenched $`N_\mathrm{f}=0`$ theory \[?\] with a result that agrees within errors with the new one. We therefore assume that at the present level of accuracy it also applies to a hypothetical QCD with three degenerate flavors. A relation between this model and real QCD is finally established by chiral perturbation theory \[?,?\] supplemented with some knowledge \[?,?,?\] of the phenomenologically inaccessible \[?\] low energy constants. The conclusion of \[?\] is that the ratio of quark masses is close to the one given at lowest order in chiral perturbation theory. Given this “fact” — but keeping in mind that it should be scrutinized in future lattice QCD computations — it is then sufficient to compute any one quark mass from lattice QCD. In particular we can connect our $`M_{\mathrm{ref}}`$ with the strange mass $`M_\mathrm{s}`$. The result that corresponds to the Gell-Mann–Oakes–Renner formula \[??\] is $`M_\mathrm{s}=48/25M_{\mathrm{ref}}`$.
Finally, the conversion of the RGI mass to the conventionally cited $`\overline{\mathrm{MS}}`$ mass at 2 $`\mathrm{GeV}`$ renormalization scale is of course based on perturbation theory, which does however look very well behaved, see Table 6.
## 2 The renormalization scheme
QCD is a theory which has as free parameters the bare gauge coupling $`g_0`$ and the bare quark masses $`m_i,i=1,\mathrm{},N_\mathrm{f}`$. The hadronic scales like $`F_\mathrm{K}`$ or $`m_\mathrm{K}`$ are connected to the perturbative high energy regime of QCD via the running of renormalized couplings
$`\overline{g}^2(\mu )=Z_gg_0^2,`$ $`\overline{m}_i(\mu )=Z_mm_i,`$ (2.4)
where $`\mu `$ is the renormalization scale. In the following we assume that the renormalization conditions are imposed at zero quark mass (mass–independent schemes) \[?\]. Introducing a regularization prescription, e.g. a lattice spacing $`a`$, the renormalization factors are functions of $`g_0`$ and $`a\mu `$
$`Z_g=Z_g(g_0,a\mu ),`$ $`Z_m=Z_m(g_0,a\mu ).`$ (2.5)
In renormalized quantities the regulator can be removed, e.g. the continuum limit $`a0`$ can be taken, yielding a finite result. The advantage of mass–independent renormalization schemes is that in all such schemes the ratios of renormalized quark masses for different flavors are scale and scheme independent constants \[?,?\].
The running of the renormalized couplings is described by the renormalization group equations (RGE)
$`\mu {\displaystyle \frac{\mathrm{d}\overline{g}}{\mathrm{d}\mu }}=\beta (\overline{g}),`$ $`\mu {\displaystyle \frac{\mathrm{d}\overline{m}_i}{\mathrm{d}\mu }}=\tau (\overline{g})\overline{m}_i.`$ (2.6)
The $`\beta `$ and $`\tau `$ functions are non-perturbatively defined if this is true for $`\overline{g}`$ and $`\overline{m}_i`$. Their perturbative expansions are
$`\beta (\overline{g})`$ $`_{\stackrel{}{\overline{g}0}}`$ $`\overline{g}^3\{b_0+b_1\overline{g}^2+b_2\overline{g}^4+\mathrm{}\},`$ (2.7)
$`\tau (\overline{g})`$ $`_{\stackrel{}{\overline{g}0}}`$ $`\overline{g}^2\{d_0+d_1\overline{g}^2+\mathrm{}\}.`$ (2.8)
The coefficients
$`b_0={\displaystyle \frac{1}{(4\pi )^2}}\left(11{\displaystyle \frac{2}{3}}N_\mathrm{f}\right),b_1={\displaystyle \frac{1}{(4\pi )^4}}\left(102{\displaystyle \frac{38}{3}}N_\mathrm{f}\right),d_0={\displaystyle \frac{8}{(4\pi )^2}},`$ (2.9)
are scheme independent. A physical quantity $`P`$ is a quantity for which the total dependence on the renormalization scale $`\mu `$ vanishes, i.e. it is a renormalization group invariant (RGI)
$`\mu {\displaystyle \frac{\mathrm{d}}{\mathrm{d}\mu }}P(\mu ,\overline{g},\{\overline{m}_i\})`$ $`=`$ $`0.`$ (2.10)
Examples are the $`\mathrm{\Lambda }`$–parameter and the RGI quark masses
$`\mathrm{\Lambda }`$ $`=`$ $`\mu (b_0\overline{g}^2)^{b_1/2b_0^2}\mathrm{e}^{1/(2b_0\overline{g}^2)}\mathrm{exp}\left\{{\displaystyle _0^{\overline{g}}}dx\left[{\displaystyle \frac{1}{\beta (x)}}+{\displaystyle \frac{1}{b_0x^3}}{\displaystyle \frac{b_1}{b_0^2x}}\right]\right\},`$ (2.11)
$`M_i`$ $`=`$ $`\overline{m}_i(2b_0\overline{g}^2)^{d_0/2b_0}\mathrm{exp}\left\{{\displaystyle _0^{\overline{g}}}dx\left[{\displaystyle \frac{\tau (x)}{\beta (x)}}{\displaystyle \frac{d_0}{b_0x}}\right]\right\},`$ (2.12)
where $`\overline{g}=\overline{g}(\mu )`$ and $`\overline{m}_i=\overline{m}_i(\mu )`$. The $`\mathrm{\Lambda }`$–parameter and the RGI quark masses are defined independent of perturbation theory and their connections between different mass independent renormalization schemes can be given in a simple and exact way \[?\]. In particular the RGI quark masses $`M_i`$ are scheme independent. Any physical quantity $`P`$ can be considered to be a function of $`\mathrm{\Lambda }`$ and $`M_i`$, i.e. there exist a function $`\widehat{P}`$ such that \[?\]
$`P(\mu ,\overline{g},\{\overline{m}_i\})`$ $`=`$ $`\widehat{P}(\mathrm{\Lambda },\{M_i\}).`$ (2.13)
For this reason $`\mathrm{\Lambda }`$ and $`M_i`$ seem preferable as the fundamental parameters of QCD.
### 2.1 Quark masses in the Schrödinger Functional
In QCD a renormalized mass is defined through the partially conserved axial current (PCAC) relation, which involves the renormalized axial current $`(A_\mathrm{R})_\mu (x)`$ and the renormalized pseudoscalar density $`P_\mathrm{R}(x)`$
$`_\mu (A_\mathrm{R})_\mu `$ $`=`$ $`(\overline{m}_i+\overline{m}_j)P_\mathrm{R},`$ (2.14)
$`(A_\mathrm{R})_\mu (x)`$ $`=`$ $`Z_\mathrm{A}\overline{\psi }_i(x)\gamma _\mu \gamma _5\psi _j(x),`$ (2.15)
$`P_\mathrm{R}(x)`$ $`=`$ $`Z_\mathrm{P}\overline{\psi }_i(x)\gamma _5\psi _j(x).`$ (2.16)
The renormalization constant $`Z_\mathrm{A}`$ can be calculated non–perturbatively, using chiral Ward identities and does not depend on the renormalization scale.
The renormalization constant $`Z_\mathrm{P}`$ can be conveniently determined in the Schrödinger Functional (SF) renormalization scheme \[?,?\]. There QCD is formulated in a finite box of spatial size $`L`$ and temporal extent $`T`$. The fields are subject to Dirichlet boundary conditions in time, which provide an infrared cutoff to the frequency spectrum of quarks and gluons. This allows to perform simulations at zero quark mass and thus to use the SF as a mass–independent renormalization scheme. Our renormalization scheme is further specified by setting $`T=L`$ (see below). The renormalization conditions are then naturally imposed at the scale $`\mu =1/L`$.
The presence of boundary values for the fields in the SF formulation of a field theory requires in general additional (compared to the case without boundaries) renormalization \[?,?\]. The renormalizability of QCD with SF boundaries has been studied in Refs. \[?,?,?,?\] where it was shown that no additional counterterms are needed except for one boundary term which amounts to a rescaling of the boundary values of the fermion fields by a logarithmically divergent factor.
The renormalization condition for $`Z_\mathrm{P}`$ that we employ is discussed in Refs. \[?,?,?\]. It uses a correlation function $`f_\mathrm{P}(x_0)`$, which is a matrix element of the pseudoscalar density inserted at time distance $`x_0`$ from a pseudoscalar boundary state, the other boundary state having vacuum quantum numbers. To cancel the multiplicative renormalization of the boundary quark fields the boundary–to–boundary correlation function $`f_1`$ is used. The renormalization constant $`Z_\mathrm{P}`$ is then defined through
$`Z_\mathrm{P}={\displaystyle \frac{\sqrt{3f_1}}{f_\mathrm{P}(L/2)}}\text{at}m_i=0,i=1,\mathrm{},N_\mathrm{f}.`$ (2.17)
The correlation functions are schematically represented in Fig. 1.
They are computed at zero quark masses $`m_i`$. The definition Eq. (2.17) is such that $`Z_\mathrm{P}=1`$ at tree level of perturbation theory. The renormalization condition for $`Z_\mathrm{P}`$ is further specified by setting
$`T=L,`$ $`C=C^{}=0,`$ $`\theta =\mathrm{\hspace{0.17em}0.5},`$ (2.18)
where $`C,C^{}`$ are the boundary gauge fields in the Lie algebra at $`x_0=0`$ and $`x_0=T`$ and $`\theta `$ the parameter controlling the spatial boundary conditions of the fermion fields. For more details of the calculation we refer to \[?,?\].
A rigorous definition of the renormalized mass can be given in the lattice regularization of QCD. In our case we work with $`N_\mathrm{f}=2`$ mass–degenerate flavors of O($`a`$) improved \[?,?\] Wilson fermions. The massless theory is defined in the bare parameter space along the line $`\kappa =\kappa _c(g_0)`$ where the PCAC mass
$`m(g_0,\kappa )`$ $`=`$ $`{\displaystyle \frac{\frac{1}{2}(_0^{}+_0)f_\mathrm{A}(x_0)+c_\mathrm{A}a_0^{}_0f_\mathrm{P}(x_0)}{2f_\mathrm{P}(x_0)}}|_{x_0=T/2}`$ (2.19)
vanishes. Here, $`_0`$ and $`_0^{}`$ are the forward and backward lattice derivatives, respectively. The O($`a`$) improvement coefficient $`c_\mathrm{A}`$ of the axial current has been computed non–perturbatively in Ref. \[?\]. The correlation function $`f_\mathrm{A}(x_0)`$ of the axial current $`A_0`$ is defined analogously to $`f_\mathrm{P}(x_0)`$. The renormalized PCAC mass at the scale $`\mu =1/L`$ can then be written as
$`\overline{m}(\mu )=\underset{a0}{lim}Z_m(g_0,a\mu )m(g_0,\kappa )|_{u=\overline{g}^2(L)},Z_m(g_0,a\mu )={\displaystyle \frac{Z_\mathrm{A}(g_0)}{Z_\mathrm{P}(g_0,L/a)}},`$ (2.20)
where the renormalized gauge coupling $`\overline{g}^2(L)`$ \[?\] is kept fixed. Since the coupling runs with the box size $`L`$, keeping $`\overline{g}^2`$ fixed means keeping the renormalization scale $`\mu `$ fixed. The renormalization factor $`Z_\mathrm{A}`$ of the axial current has been determined non–perturbatively in Ref. \[?\]. The precise definition of $`Z_\mathrm{P}`$ (differing at order $`a^2`$ from Eq. (2.17)) is
$`Z_\mathrm{P}(g_0,L/a)`$ $`=`$ $`c{\displaystyle \frac{\sqrt{3f_1}}{f_\mathrm{P}(L/2)}}\text{at}\kappa =\kappa _c,`$ (2.21)
where the factor $`c(L/a)`$ is chosen such that $`Z_P(0,L/a)=1`$ and can be found in \[?\].
In order that the continuum limit is reached with cutoff effects strictly proportional to $`a^2`$, the $`O(a)`$ improvement factor $`(1+(b_\mathrm{A}b_\mathrm{P})am_\mathrm{q})`$ should be included in Eq. (2.20) \[?\]. Here $`am_\mathrm{q}=(1/\kappa 1/\kappa _c)/2`$ is the bare subtracted quark mass. Results in perturbation theory \[?\] and in the quenched approximation \[?\] show that the difference of improvement coefficients $`b_\mathrm{A}b_\mathrm{P}`$ is small. The quark masses in our simulations will also be relatively small and we expect corrections due to $`b_\mathrm{A}b_\mathrm{P}`$ at the per mille level, which we neglect.
## 3 The running of the mass in the SF–scheme
The running of the renormalized quark mass $`\overline{m}(\mu )`$ in the SF scheme as specified in Section 2.1 with $`N_\mathrm{f}=2`$ mass–degenerate flavors can be computed on the lattice from the step scaling function of the renormalization factor $`Z_\mathrm{P}`$ extrapolated to the continuum
$`\sigma _\mathrm{P}(u)`$ $`=`$ $`\underset{a0}{lim}\mathrm{\Sigma }_\mathrm{P}(u,a/L),\mathrm{\Sigma }_\mathrm{P}(u,a/L)={\displaystyle \frac{Z_\mathrm{P}(g_0,2L/a)}{Z_\mathrm{P}(g_0,L/a)}}|_{u=\overline{g}^2(L)}.`$ (3.22)
From Eq. (2.20) it follows immediately that
$`\sigma _\mathrm{P}(u)`$ $`=`$ $`{\displaystyle \frac{\overline{m}(\mu )}{\overline{m}(\mu /2)}}\text{for }\mu =1/L,`$ (3.23)
i.e. the step scaling function $`\sigma _\mathrm{P}(u)`$ describes the running of the renormalized quark mass. We computed $`\mathrm{\Sigma }_\mathrm{P}(u,a/L)`$ at six values of the renormalized coupling $`u`$ corresponding approximately to a range of box sizes of the order $`L=10^2\mathrm{fm}\mathrm{}1\mathrm{fm}`$ (or equivalently $`\mu `$ of the order $`100\mathrm{GeV}\mathrm{}1\mathrm{GeV}`$). At each value of $`u`$ we simulated three lattice resolutions $`L/a=6,8,12`$ and the results for $`Z_\mathrm{P}`$ and $`\mathrm{\Sigma }_\mathrm{P}`$ are summarized in Table 7 in Appendix A. For the extrapolation to the continuum, we fitted to a constant the two values of $`\mathrm{\Sigma }_\mathrm{P}`$ on the finer lattices, separately for each coupling $`u`$. We then added linearly the difference between the fit and the $`L/a=8`$ result as a systematic error. The continuum estimates can be seen in Fig. 2. Our data do not show any significant dependence on the lattice spacing, as we could verify by trying different extrapolations (quadratic, linear in $`a`$). This statement is based on the statistical accuracy that we could achieve. We remark that also in the quenched approximation the cutoff effects were found to be small \[?\] and there $`\mathrm{\Sigma }_P`$ was computed at an even finer lattice resolution $`L/a=16`$. The continuum values of $`\sigma _\mathrm{P}(u)`$ are summarized in Table 1. In the last column we list the $`\chi ^2`$ divided by the number of degrees of freedom $`n_{\mathrm{df}}`$ of the fit. Their average is close to the expected value of one.
In perturbation theory the step scaling function $`\sigma _\mathrm{P}(u)`$ has an expansion $`\sigma _\mathrm{P}(u)=1\mathrm{ln}(2)d_0u+\mathrm{O}(u^2)`$. In Fig. 3 our non–perturbative data for $`\sigma _\mathrm{P}(u)`$ are conveniently plotted for comparison with perturbation theory. We relate $`\sigma _\mathrm{P}(u)`$ to $`\tau `$ and $`\beta `$ using Eq. (2.12)
$`\sigma _\mathrm{P}(u)`$ $`=`$ $`\left({\displaystyle \frac{u}{\sigma (u)}}\right)^{d_0/(2b_0)}\mathrm{exp}\left\{{\displaystyle _\sqrt{u}^{\sqrt{\sigma (u)}}}dx\left[{\displaystyle \frac{\tau (x)}{\beta (x)}}{\displaystyle \frac{d_0}{b_0x}}\right]\right\},`$ (3.24)
where $`\sigma (u)`$ is the step scaling function of the coupling and is determined by \[?\]
$`2\mathrm{ln}(2)`$ $`=`$ $`{\displaystyle _u^{\sigma (u)}}dx{\displaystyle \frac{1}{\sqrt{x}\beta (\sqrt{x})}}.`$ (3.25)
Using for the $`\tau `$–function the 2–loop expression with \[?\]
$`d_1`$ $`=`$ $`d_0(0.0271+0.0105N_\mathrm{f})`$ (3.26)
and for the $`\beta `$–function the 3–loop expression with (see \[?\])
$`b_2`$ $`=`$ $`(0.4830.275N_\mathrm{f}+0.0361N_\mathrm{f}^20.00175N_\mathrm{f}^3)/(4\pi )^3`$ (3.27)
we get from Eq. (3.24) and Eq. (3.25) the perturbative curve shown in Fig. 3. Our non–perturbative data do not show any significant deviation from the perturbative estimates.
We now take the continuum values for $`\sigma _\mathrm{P}(u)`$ in Table 1 (with added statistical and systematic errors) and for the step scaling function of the coupling $`\sigma (u)`$ the results from fits to constants in Table 4 of Ref. \[?\]. We solve the following joint recursion to evolve coupling and mass from a low energy scale $`1/L_{\mathrm{max}}`$ defined by
$`u_0=\overline{g}^2(L_{\mathrm{max}})`$ $`=`$ $`4.61`$ (3.28)
to the higher scales $`1/L_k`$, $`k=0,1,\mathrm{},8`$ (with $`L_0L_{\mathrm{max}}`$)
$`\{\begin{array}{c}u_0=\overline{g}^2(L_{\mathrm{max}})=4.61\hfill \\ \sigma (u_{k+1})=u_k\hfill \end{array}`$ $``$ $`u_k=\overline{g}^2(L_k),L_k=2^kL_{\mathrm{max}},`$ (3.31)
$`\{\begin{array}{c}w_0=1\hfill \\ w_k=\left[_{i=1}^k\sigma _\mathrm{P}(u_i)\right]^1\hfill \end{array}`$ $``$ $`w_k={\displaystyle \frac{\overline{m}(1/L_{\mathrm{max}})}{\overline{m}(1/L_k)}}.`$ (3.34)
We interpolate the values of $`\sigma (u)`$ and $`\sigma _\mathrm{P}(u)`$ through a polynomial ansatz
$`\sigma (u)`$ $`=`$ $`u+s_0u^2+s_1u^3+s_2u^4+s_3u^5+s_4u^6,`$ (3.35)
$`\sigma _\mathrm{P}(u)`$ $`=`$ $`1+p_0u+p_1u^2+p_2u^3,`$ (3.36)
where the coefficients $`s_0`$, $`s_1`$ \[?\] and $`p_0=\mathrm{ln}(2)d_0`$ are fixed to their perturbative values. The coefficients $`s_2,s_3,s_4`$ and $`p_1,p_2`$ are here fit parameters. The errors of the recursion coefficients are computed by error propagation.
Using the coefficients $`w_k`$ in Eq. (3.34) we compute
$`{\displaystyle \frac{M}{\overline{m}(1/L_{\mathrm{max}})}}`$ $`=`$ $`w_k^1{\displaystyle \frac{M}{\overline{m}(1/L_k)}},`$ (3.37)
where the factor $`M/\overline{m}(1/L_k)`$ is calculated from Eq. (2.12) with $`\overline{g}^2=u_k`$ by employing the perturbative expressions for the $`\tau `$– and $`\beta `$–functions at 2– respectively 3–loop order. The results are shown in Table 2. They have a remarkable stability in the coupling<sup>2</sup><sup>2</sup>2 The deviation in the case $`k=0`$ is due to the difference between the perturbative and the non–perturbative values of $`\sigma (u)`$ at large $`u`$ (see \[?\]). $`u_k`$ and we take $`k=6`$ as our result
$`{\displaystyle \frac{M}{\overline{m}(\mu )}}`$ $`=`$ $`1.297(16)\text{at}\mu =\mathrm{\hspace{0.17em}1}/L_{\mathrm{max}}.`$ (3.38)
We emphasize that from Eq. (2.12) it is evident that $`\sigma _\mathrm{P}(u)`$ and $`M/\overline{m}(\mu )`$ are flavor independent. Moreover, since the continuum limit has been taken any regularization dependence has been removed from the result Eq. (3.38).
Finally, in Fig. 4 we plot the non–perturbative running of the renormalized mass. For $`\mu /\mathrm{\Lambda }=\mathrm{\hspace{0.17em}1}/(L_k\mathrm{\Lambda })`$, $`k=0,1,\mathrm{},8`$, we plot the points $`\overline{m}(1/L_k)/M`$ obtained from Eq. (3.37) using the result Eq. (3.38). The physical scale $`\mathrm{\Lambda }`$ is here implicitly determined through $`\mathrm{ln}(\mathrm{\Lambda }L_{\mathrm{max}})=1.298(58)`$ obtained from the recursion Eq. (3.31). In the plot we neglect the overall uncertainties of $`\overline{m}(1/L_{\mathrm{max}})/M`$ and $`\mathrm{ln}(\mathrm{\Lambda }L_{\mathrm{max}})`$, since they would simply change the scales on the plot axes. The errors of the points in Fig. 4 come from the coefficients $`w_k`$. Together with the non–perturbative points we show the perturbative curves that are obtained from Eq. (2.11) and Eq. (2.12) by using the perturbative expressions for the $`\tau `$– and $`\beta `$–functions at 1–, 2– respectively 2–, 3–loop order. The non–perturbative and perturbative running are very close down to the smallest energies that were accessible in our simulations. We remark that this statement explicitly refers to the special SF renormalization scheme considered here.
## 4 Estimate of the strange quark mass
### 4.1 Complete $`Z_\mathrm{M}`$ for the improved Wilson discretization
We now derive the second factor in Eq. (1.3) for a few values of the lattice spacing or respectively the bare coupling. As emphasized before, this contribution is non–universal and in the form given it will be valid only for our action of non–perturbatively improved Wilson fermions with plaquette gauge action and $`c_{\mathrm{sw}}`$ as specified in \[?\]. For $`Z_\mathrm{A}`$ and $`c_\mathrm{A}`$ we employ the values recently given and parameterized in \[?\] and \[?\]. It remains to compute $`Z_\mathrm{P}(g_0,L_{\mathrm{max}}/a)`$ for the desired values of the bare coupling $`g_0`$, here given by $`\beta =5.2,5.29,5.4`$. The scale $`L_{\mathrm{max}}`$ is fixed by $`\overline{g}^2(L_{\mathrm{max}})=4.61`$, where the universal factor of $`Z_\mathrm{M}`$ is now known, Eq. (3.38). Our basis are the simulation results summarized in Table 3.
While the simulation at the largest bare coupling is exactly at the target value for $`\overline{g}^2`$, the two other series of simulations require a slight interpolation. This has been done using a fit ansatz motivated by Eq. (2.12)
$`\mathrm{ln}(Z_\mathrm{P})`$ $`=`$ $`c_1+c_2\mathrm{ln}(\overline{g}^2)`$ (4.39)
to interpolate $`Z_\mathrm{P}`$ between two values of $`\overline{g}^2`$ straddling $`4.61`$. The fit takes into account the (independent) errors of both $`Z_\mathrm{P}`$ and $`\overline{g}^2`$. The fit error is then augmented by the difference between the fit result from Eq. (4.39) and the result from a simple two point linear interpolation in $`\overline{g}^2`$. The values of the coefficient $`c_2`$ in the fit Eq. (4.39) are found to be $`0.369(25)`$ at $`\beta =5.29`$ and $`0.430(34)`$ at $`\beta =5.40`$, which are not far from $`d_0/(2b_0)=0.4138`$.
The resulting numbers for $`Z_\mathrm{P}`$ and $`Z_\mathrm{M}`$ are summarized in Table 4. The first error for $`Z_\mathrm{M}`$ comes from the error of the factor $`Z_\mathrm{A}/Z_\mathrm{P}`$. The second error is the 1.2% uncertainty in the universal factor $`M/\overline{m}`$ and should be added in quadrature to the quark mass error after the continuum limit has eventually been taken.
### 4.2 The reference quark mass
As announced in the introduction we next compute the reference quark mass $`M_{\mathrm{ref}}`$ producing a pseudoscalar with the mass of the Kaon in our simulated two-flavour theory. The hopping parameter $`\kappa _{\mathrm{ref}}`$ is tuned to keep the pseudoscalar mass $`m_{\mathrm{PS}}`$ in the relation
$`(r_0(\kappa _c)m_{\mathrm{PS}}(\kappa _{\mathrm{ref}}))^2`$ $`=`$ $`(r_0m_\mathrm{K})^2=\mathrm{\hspace{0.17em}1.5736},`$ (4.40)
where $`r_0`$ is the scale extracted from the static quark potential \[?\]. This value corresponds to the K–meson mass $`m_\mathrm{K}^2=(495\mathrm{MeV})^2`$ for $`r_0=0.5\mathrm{fm}`$ \[?\]. Using the data for $`r_0(\kappa )`$, $`m_{\mathrm{PS}}(\kappa )`$ and $`\kappa _c`$ available in Ref. \[?\] and the extrapolations to the chiral limit $`r_0(\kappa _c)`$ of Ref. \[?\], we can safely perform a slight extrapolation to the value $`\kappa =\kappa _{\mathrm{ref}}`$ defined in Eq. (4.40). We employ the following fit ansatz for $`r_0m_{\mathrm{PS}}`$ as a function of the bare subtracted quark mass $`am_\mathrm{q}`$
$`{\displaystyle \frac{(r_0(\kappa _c)m_{\mathrm{PS}}(\kappa ))^2}{am_\mathrm{q}}}`$ $`=`$ $`e_1+e_2am_\mathrm{q},`$ (4.41)
with fit coefficients $`e_1`$ and $`e_2`$. The results of the fits are $`\kappa _{\mathrm{ref}}=0.135680(30)`$ at $`\beta =5.20`$, $`\kappa _{\mathrm{ref}}=0.136018(27)`$ at $`\beta =5.29`$ and $`\kappa _{\mathrm{ref}}=0.136293(25)`$ at $`\beta =5.40`$.
The PCAC masses $`am_{\mathrm{ref}}`$ have been computed in simulations at $`\kappa =\kappa _{\mathrm{ref}}`$ and the results are summarized in Table 5. In these simulations the parameter $`\theta `$ in Eq. (2.18) has been set to zero. Barring cutoff effects the PCAC mass is independent of the time $`x_0`$ at which the right hand side of Eq. (2.19) is evaluated. Therefore in Table 5 we average over a range $`t_1:t_2`$ of $`x_0`$ values around $`x_0=L/2`$, where the PCAC mass has a plateau. We keep $`t_1t_2`$ roughly constant in physical units. We see a significant gain in statistical precision due to the averaging compared to taking only the time $`x_0=L/2`$. The error analysis of derived observables like the (averaged) PCAC mass and $`Z_\mathrm{P}`$ has been done with the method of Ref. \[?\].
The cutoff effects in $`m_{\mathrm{ref}}`$ also depend on the volume. In fact, at a lattice spacing of $`0.1\mathrm{fm}`$ this dependence is rather strong \[?\]. Thus one has to define the PCAC mass for a fixed volume in order to ensure that cutoff effects disappear smoothly as O($`a^2`$). Furthermore it is preferable to choose $`L`$ relatively large. We choose $`L1.5\mathrm{fm}`$, where the volume dependence can be neglected. Indeed, we compared results from a simulation at $`\beta =5.2`$, $`\kappa =0.1355`$, $`c_{\mathrm{sw}}=2.02`$ on a $`L/a=16`$ lattice with those from the JLQCD collaboration<sup>3</sup><sup>3</sup>3 We are grateful to Takashi Kaneko for providing us with data for this comparison. obtained at the same parameters on a $`20^3\times 48`$ lattice \[?\]. We could check that at the level of 1.6% statistical precision in our simulation, the PCAC masses agree. We conclude that our volume is large enough at $`\beta =5.2`$ for the volume dependence to be negligible. Therefore at the higher $`\beta `$ values we choose approximately matched (or larger) physical volumes.
The second to last column of Table 5 shows our results for the reference RGI quark mass
$$M_{\mathrm{ref}}=Z_\mathrm{M}m_{\mathrm{ref}}.$$
(4.42)
We do see large cutoff effects. At $`\beta =5.2`$ similarly large cutoff effects have been observed in other quantities as well \[?,?,?\]. Moreover the lattice spacing in our simulations changes by 30% only when going from $`\beta =5.2`$ to $`\beta =5.4`$. We therefore do not attempt (and discourage) any elaborate continuum extrapolation based on this range of lattice spacings. Instead we use a conservative estimate of the continuum value of $`M_{\mathrm{ref}}`$ by taking its value at our largest $`\beta =5.4`$ with the difference to the value at $`\beta =5.2`$ added as systematic error
$`M_{\mathrm{ref}}`$ $`=`$ $`72(3)(13)\mathrm{MeV}.`$ (4.43)
This result makes it clear that at present the systematic error is dominating over the statistical one.
The last column of Table 5 shows the results for an alternative estimate of $`M_{\mathrm{ref}}`$, which differs only by the use of the renormalization factor $`Z_\mathrm{A}^{\mathrm{con}}`$ \[?\]. We observe that the difference in $`M_{\mathrm{ref}}`$ between the two choices for $`Z_\mathrm{A}`$ is compatible with an $`a^2`$–behavior. Actually, the data obtained with $`Z_\mathrm{A}^{\mathrm{con}}`$ have a weaker $`a`$–dependence but this effect is essentially due to the change at the coarsest lattice spacing ($`\beta =5.2`$).
### 4.3 The strange quark mass
To make contact with physics we now assume that Eq. (4.43) also holds within errors in a theory with three degenerate flavors, which can then be related to the strange quark in QCD by chiral perturbation theory. This assumption is supported by the fact that the value of $`M_{\mathrm{ref}}`$ Eq. (4.43) is the same within errors as in the $`N_\mathrm{f}=0`$ theory \[?\].
At lowest order in chiral perturbation theory, disregarding the electromagnetic interaction, the formula \[??,?,?\]
$`m_\mathrm{K}^2={\displaystyle \frac{1}{2}}(m_{\mathrm{K}^+}^2+m_{\mathrm{K}^0}^2)`$ $`=`$ $`(\widehat{M}+M_\mathrm{s})B_{\mathrm{RGI}},`$ (4.44)
where $`\widehat{M}=1/2(M_\mathrm{u}+M_\mathrm{d})`$, holds. Here $`M_\mathrm{u}`$, $`M_\mathrm{d}`$, $`M_\mathrm{s}`$ are the up, down, strange RGI quark masses and $`B_{\mathrm{RGI}}`$ is a constant of the chiral Lagrangian. Eq. (4.44) implies for degenerate quarks of RGI mass $`M_{\mathrm{ref}}`$ (defined according to Eq. (4.40))
$`m_\mathrm{K}^2`$ $`=`$ $`2M_{\mathrm{ref}}B_{\mathrm{RGI}}.`$ (4.45)
Therefore the relation $`M_{\mathrm{ref}}=(\widehat{M}+M_\mathrm{s})/2`$ holds at lowest order in chiral perturbation theory and using $`M_\mathrm{s}/\widehat{M}=24.4(1.5)`$ \[?\] gives
$`M_\mathrm{s}`$ $``$ $`48/25M_{\mathrm{ref}}.`$ (4.46)
Corrections to Eq. (4.46) are expected to be small in chiral perturbation theory \[?,?,?,?\] and anyway below the accuracy that we will reach here for our result. In addition in the quenched approximation it was found that the dependence of pseudoscalar masses $`m_{\mathrm{PS}}`$, at fixed average quark mass, on the difference of quark masses is rather small \[?\]. We therefore assume Eq. (4.46) to hold also with dynamical quarks.
At this point we are ready to give an estimate of the continuum value of the RGI strange quark mass
$`M_\mathrm{s}`$ $`=`$ $`138(5)(26)\mathrm{MeV},`$ (4.47)
by combining Eq. (4.46) with Eq. (4.43). Equivalently to the determination of $`M_{\mathrm{ref}}`$, from Eq. (4.45) and Eq. (4.43) we get
$`B_{\mathrm{RGI}}`$ $`=`$ $`1.70(38)\mathrm{GeV}.`$ (4.48)
The renormalized strange quark mass in the $`\overline{\mathrm{MS}}`$ scheme at the renormalization scale $`\mu `$ is obtained by multiplying $`M_\mathrm{s}`$ with the conversion factor $`\overline{m}^{\overline{\mathrm{MS}}}(\mu )/M`$. The latter is computed perturbatively by numerical integration of Eq. (2.11) and Eq. (2.12) and is listed in Table 6 for different choices of $`\mu `$. The $`n=2,3,4`$–loop approximations of the $`\beta `$ and $`\tau `$ functions in the $`\overline{\mathrm{MS}}`$ scheme (for $`N_\mathrm{f}=2`$) are used. The coefficients at 4–loop have been computed in \[?,?,?,?\]. To set $`\mu `$ in physical units the result $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}^{(2)}=245(32)`$ of Ref. \[?\] is used. The uncertainty in the $`\mathrm{\Lambda }`$ parameter translates into a 2.8% uncertainty in $`\overline{m}^{\overline{\mathrm{MS}}}(\mu )/M`$ at $`\mu =2\mathrm{GeV}`$ and 4.9% at $`\mu =1\mathrm{GeV}`$ (with 4–loop evolution). Taking this into account our estimate for the $`\overline{\mathrm{MS}}`$ strange quark mass is
$`\overline{m}_\mathrm{s}^{\overline{\mathrm{MS}}}(\mu )`$ $`=`$ $`97(22)\mathrm{MeV}\text{at}\mu =2\mathrm{GeV}.`$ (4.49)
## 5 Conclusions and outlook
We have presented a fully non–perturbative renormalization of the quark mass in two flavor QCD, with renormalization conditions specified at zero quark mass in the Schrödinger Functional scheme. Simulations were performed with O($`a`$) improved Wilson quarks. Our main results are the running of the quark mass in Fig. 4 and the factor $`M/\overline{m}(\mu )`$ relating the quark mass at a specified low energy scale $`\mu `$ with the RGI quark mass in the continuum limit, Eq. (3.38).
In order to obtain the renormalized strange quark mass in physical units an appropriate hadronic scheme has to be defined. By using existing data on the pseudoscalar masses and the scale $`r_0`$ from Ref. \[?\] we were able to determine the RGI reference quark mass $`M_{\mathrm{ref}}`$ defined through Eq. (4.42) and Eq. (4.40), at three lattice spacings in the approximate range $`0.092\mathrm{}0.071\mathrm{fm}`$. Even in this small range we see distinct cutoff effects. At the largest lattice spacing large cutoff effects have been observed elsewhere \[?,?,?\]. These facts make it impossible to perform a systematic continuum extrapolation. Nevertheless we can give a conservative estimate of the continuum value in Eq. (4.43). We emphasize that the error in Eq. (4.43) is dominated by the systematic uncertainty of this step.
Using a number of additional reasonable assumptions we convert $`M_{\mathrm{ref}}`$ to the strange quark mass in the $`\overline{\mathrm{MS}}`$ scheme Eq. (4.49). A definite clarification that all of these lead to negligible errors will require some future work.
In Fig. 5 we summarize the most recent results for the strange quark mass from lattice simulations with two \[?,?,?\] and two plus one \[?,?\] dynamical quarks. The physical K–meson mass is always used as input although in some cases this is employed in a partially quenched setup where valence and sea quark masses differ, while we compute $`M_{\mathrm{ref}}`$ and then use chiral perturbation theory to connect to the physical theory. In the legend we emphasize whether nonperturbative or perturbative renormalization has been used as well as the adopted discretizations. It appears that perturbative renormalization leads to rather small values for the strange quark mass.<sup>4</sup><sup>4</sup>4 Note that the use of lowest order chiral perturbation theory in our work is not likely to be a significant source of difference to the other computations, since they do not report large deviations from lowest order chiral perturbation theory, see for example Ref. \[?\]. We plot the results of our present work using two renormalization conditions for $`Z_\mathrm{A}`$ that differ at O($`a^2`$). The red triangles are obtained with $`Z_\mathrm{A}^{\mathrm{con}}`$ \[?\] and are slightly displaced for clarity. This comparison gives a flavor of the quite large cutoff effects at $`\beta =5.2`$, our coarsest lattice spacing. The dotted lines in Fig. 5 mark the quenched result of Ref. \[?\]. Given the present status, illustrated in Fig. 5, it appears hard to claim a definite dependence of $`m_\mathrm{s}`$ on the number of dynamical fermions, even between $`N_\mathrm{f}=0`$ and $`N_\mathrm{f}=2`$.
Our present result for the renormalized strange quark mass should be improved by simulating at a finer lattice spacing. We are confident that this will be possible in the near future. There are promising improvements in algorithms \[?,?,?\], and a new machine apeNEXT is becoming available to us \[?,?\]. Moreover the determination of the lattice spacing can be improved by, for example, computing the renormalized strange quark mass in units of the K–meson decay constant $`F_\mathrm{K}`$. There are first indications that the SF is an efficient setup for this low–energy computation (see also \[?\]).
Acknowledgement. We are grateful to Gerrit Schierholz for communicating results of Ref. \[?\]. We further thank NIC/DESY for allocating computing resources on the APEmille machine to this project and the APE collaboration and the staff of the computer center at DESY, Zeuthen for their support.
The computation of the renormalized quark mass is part of project B2 of the SFB Transregio 9 “Computational Particle Physics” and has also been supported by the Deutsche Forschungsgemeinschaft (DFG) in the Graduiertenkolleg GK 271 as well as by the European Community’s Human Potential Programme under contract HPRN-CT-2000-00145.
## Appendix A Simulation results for $`Z_\mathrm{P}`$
In Table 7 we collect the bare parameters and results of our simulations to compute $`Z_\mathrm{P}`$. Simulations on $`L/a`$ and $`2L/a`$ lattices are required to extract the step scaling function $`\mathrm{\Sigma }_\mathrm{P}`$ Eq. (3.22). At the three lowest couplings $`\overline{g}^2(L)`$ simulations have been performed using the 1–loop value of $`c_\mathrm{t}`$ \[?\], except for $`L/a=6`$, $`\beta =7.5420`$ and $`L/a=8`$, $`\beta =7.7206`$. For the latter parameters and the larger couplings the 2–loop value of $`c_\mathrm{t}`$ \[?\] has been used. At the third lowest coupling $`u1.5`$ we checked at $`L/a=6`$ that there is no significant difference in $`\mathrm{\Sigma }_\mathrm{P}`$ using the 1– or 2–loop value for $`c_\mathrm{t}`$, as it is shown in Fig. 2.
To the statistical error of $`\mathrm{\Sigma }_\mathrm{P}`$ we added in quadrature the error due to the uncertainty in the coupling $`u`$. The latter was estimated using the 1–loop result $`\mathrm{ln}(2)d_0`$ for the derivative of $`\mathrm{\Sigma }_\mathrm{P}`$ with respect to $`u`$. This correction is tiny, it increases the errors of $`\mathrm{\Sigma }_\mathrm{P}`$ by at most $`5\%`$ at the largest coupling and lattice.
## Appendix B About the algorithm
Our simulations have been performed using the Hybrid Monte Carlo algorithm with two pseudo–fermion fields as proposed by M. Hasenbusch \[?,?,?\]. Here we extend the study on algorithmic precision presented in Ref. \[?\].
For our largest lattice volume $`L/a=24`$ at $`\beta =5.4`$ we investigated the issue of reversibility. We looked at the quantity $`|\mathrm{\Delta }H(\mathrm{cycle})|`$ which measures the difference of the Hamiltonian for a cyclic trajectory \[?\]. In Fig. 6 we present data for the average $`|\mathrm{\Delta }H(\mathrm{cycle})|`$ (empty symbols) and the maximum of $`|\mathrm{\Delta }H(\mathrm{cycle})|`$ (filled symbols) on three sets of configurations. Sticking to single precision arithmetic we found for a trajectory length $`\tau =0.5`$ (that we use in the production run) an average value for $`|\mathrm{\Delta }H(\mathrm{cycle})|`$ (black triangles) that is larger by almost a factor five with respect to the simulations at $`\beta =6.118`$ reported in Ref. \[?\] (blue circles). We therefore repeated the simulation at $`\beta =5.4`$ in double precision arithmetic. At the same time we reduced the parameter $`ϵ^2`$, defined<sup>5</sup><sup>5</sup>5 We adopt here the commonly used convention to define $`ϵ^2`$ as the square of the ratio between the norm of the residue vector and the norm of the source vector. It differs from the definition used in \[?\], where $`ϵ^2`$ was normalized by the norm of the solution vector. as the requested accuracy in the conjugate gradient iteration, from $`ϵ^2=10^{11}`$ (as used with single precision) to $`ϵ^2=10^{13}`$. The results for double precision are the red squares in Fig. 6. By comparing single with double precision, both the average and the maximal value of $`|\mathrm{\Delta }H(\mathrm{cycle})|`$ decrease by two orders of magnitude for double precision and these values are well below the values observed in Ref. \[?\]. Our statistics for the reversibility checks consists of 18 configurations analyzed at $`\beta =5.4`$ (both for single and double precision) and of 30 configurations analyzed at $`\beta =6.118`$. On a subset of the $`\beta =5.4`$ configurations we checked that within single precision by reducing $`ϵ^2`$ to $`10^{13}`$ (which one would expect to be the lower bound in single precision) the roundoff error in the Hamiltonian computation and consequently the value of $`|\mathrm{\Delta }H(\mathrm{cycle})|`$ remain unchanged.
Finally, we emphasize that we found no significant difference in the observables we computed at $`\beta =5.4`$ between the single and double precision simulations. The number quoted in Table 5 for $`m_{\mathrm{ref}}`$ is from the simulation in double precision. |
warning/0507/math0507582.html | ar5iv | text | # Internal Diffusion Limited Aggregation on discrete groups having exponential growth
## 1 Introduction
Let $`𝒢`$ be an infinite discrete set and let $`P`$ be a irreducible transition matrix on that set. The matrix $`P`$ allows to define a sequence of i.i.d. Markov chains $`(S^n())_n^{}`$ with a common starting point $`O`$ ($`S(0)=O`$). The Internal Diffusion Limited Aggregation (Internal DLA) is a Markov chain $`A(n)`$ $`(n)`$ of increasing subsets of $`𝒢`$ defined as follows:
* $`A(1)=\{O\}`$,
* $`A(n+1)=A(n)S^{n+1}(\tau _{A(n)^c})`$,
where $`\tau _A`$ is the hitting time of the set $`A`$.
Thus, the set $`A(n)`$ grows as follows. At each discrete time $`n`$, we start the Markov chain $`S^n()`$ (moving according to $`P`$) at $`O`$, wait until it leaves the previous set $`A(n1)`$, and add the first element visited outside $`A(n1)`$, to obtain $`A(n)`$. Note that the Markov chain $`A(n)`$ is well defined since, for all $`n`$, $`\tau _{A(n)^c}<\mathrm{}`$ almost surely, as $`A(n)`$ is finite and $`S^{n+1}()`$ is irreducible.
The name Internal DLA comes from the similitude with the (external) DLA. In the latter model, the Markov chains start “at infinity”, conditioned to hit the previous cluster, and are stuck just before hitting it. See for example Lawler for a precise definition and some properties.
The Internal DLA was introduced by Diaconis and Fulton in 1991 . They gave the first shape theorem on $``$ where $`S^n()`$ are Simple Random Walks (nearest neighbours random walks with uniform transition probability). Later, other shape theorems were proved on $`^d`$ (Lawler, Bramson & Griffeath , Lawler , Blachère ). These results set that, when the random walk is centered (under some moment condition), the model has a limiting shape (see one of these for precise definition). Roughly, that means that the shape of $`A(n)`$ tends to the shape of the balls of $`^d`$ for a norm associated to the transition matrix $`P`$.
In the present article, we focus on groups $`\mathrm{\Gamma }`$ having exponential growth. It means that the balls (associated to the word metric for any finite generating set) have cardinality growing exponentially in the radius. On such groups, we consider symmetric random walks with finitely supported increment. Our main result (Theorem 3.1) is a shape theorem for the Internal DLA. The meaning of a shape theorem is to compare the random shape of the Internal DLA with a deterministic exhaustion of the group. For that purpose, we define a new distance on groups, called the hitting distance, by
$$d(x,y)\stackrel{\text{def.}}{=}\mathrm{ln}(F(x,y)),$$
where $`F(x,y)`$ is the hitting probability of $`y`$ for a random walk starting at $`x`$. As the Green function $`G(,)`$ is proportional to the function $`F(,)`$, the balls for that distance are in fact the sets of the form,
$$\{x\mathrm{s}.\mathrm{t}.G(e,x)N\},$$
where $`e`$ is the identity of the group. These sets exhaust $`\mathrm{\Gamma }`$ when $`N`$ goes to $`0`$.
This distance can be compared to the word distance but they are not equivalent in general. Outside its link with the Internal DLA, an interesting property of that distance is that one can deduce the exact exponential growth rate of the balls (Proposition 2.3).
The shape theorem states that the Internal DLA is close to the balls $`B(n)`$ of radius $`Kn`$ for the hitting distance ($`K`$ is a constant that ensures that the ball $`B(n)`$ contains the boundary of the ball $`B(n1)`$). More precisely there exists some constants $`C_I`$ and $`C_O`$ such that at time $`V(n)`$ (=volume of $`B(n)`$) the sets of the IDLA satisfy:
$$B(nC_I\mathrm{ln}n)A(V(n))B(n+C_O\sqrt{n})$$
for every large $`n`$, with probability $`1`$. This result gives the upper bounds for the fluctuations between $`A(V(n))`$ and $`B(n)`$: the inner fluctuations are at most logarithmic, whence the outer fluctuations are at most of order square root of the radius. The proof of that result is adapted from the one on $`^d`$ (Lawler, Bramson & Griffeath and Lawler ) with new ingredients related to the hitting distance and some simplifications, partially due to the setting of groups with exponential growth.
For the Simple Random Walk of homogeneous trees (where the hitting distance and the word distance are proportional) such a result was already obtained by Etheridge and Lawler (unpublished). In that setting, they also announced (in ) that the lower bounds for the inner and outer fluctuations were of the same order than the upper ones, but according to Lawler, a written proof of that result was missing and the lower bound for the inner error could be considered an open question. Therefore, we prove (Propositions 4.1 and 4.2) that in the case of a homogeneous tree the order of the inner fluctuation cannot be smaller of $`\mathrm{ln}n`$, while the outer fluctuation is bigger than some constant times $`\sqrt{n}`$.
The article is constructed as follows. Section 2 defines the setting of symmetric random walks on groups with exponential growth, defines the “hitting distance” and compares it to the word distance. Section 3 gives the proof of our main result, the shape theorem for the Internal DLA with upper bounds for the fluctuations. Section 4 focus on Simple Random Walk on homogeneous trees and proves the sharpness of the bounds for these examples.
We would like to warmly acknowledge Wolfgang Woess for his help and support (both mathematical and financial) along this work.
## 2 Symmetric random walk on discrete groups with exponential growth
Let $`\mathrm{\Gamma }`$ be a finitely generated group. Let $`𝒮`$ be a finite symmetric set of generators. We write each element $`x`$ of $`\mathrm{\Gamma }`$ as a word composed by these generators (letters). The minimal length among the words representing $`x`$ is called the word distance $`|x|`$. Let $`e`$ be the identity of $`\mathrm{\Gamma }`$. According to this distance, we define the balls $`B_w(n)=B_w(e,n)\stackrel{\text{def.}}{=}\{x\mathrm{\Gamma }:|x|n\}`$ and their volume $`V_w(n)\stackrel{\text{def.}}{=}\mathrm{\#}B_w(n)`$. The exponential growth property corresponds to
$$V_w(n)\mathrm{exp}(cn),$$
for some constant $`c`$. Note that an exponential upper bound of the volume is always valid since the fastest volume growth occurs when there is no relation between the generators. Hence, $`V_w(n)1+(\mathrm{\#}𝒮1)^n`$. Note also that this property does not depend on the choice of the generating set.
Let $`\mu `$ be a symmetric probability law whose support is $`𝒮`$. Let $`(X_k)_k^{}`$ be a sequence of i.i.d. random variables whose common law is $`\mu `$. The process
$$S(k)\stackrel{\text{def.}}{=}xX_1X_2\mathrm{}X_k,$$
with $`S(0)=x`$, is a symmetric irreducible random walk on $`\mathrm{\Gamma }`$ starting at $`x`$. The process $`S()`$ is then a nearest neighbors symmetric random walk on the Cayley graph $`𝒢=(\mathrm{\Gamma },𝒮)`$ on which we study the Internal DLA (with $`A(0)=\{e\}`$). We denote $`^x`$ and $`𝔼^x`$, respectively, the probability and expectation related to a random walk starting at $`x`$. When $`x=e`$, the exponent will be omitted.
The Green function $`G(x,y)`$ is defined as the expected number of visits at $`y`$ for a random walk starting at $`x`$:
$$G(x,y)\stackrel{\text{def.}}{=}𝔼^x\left[\underset{k=0}{\overset{\mathrm{}}{}}\text{ 1 I}_{\{S(k)=y\}}\right]=\underset{k=0}{\overset{\mathrm{}}{}}^x[S(k)=y].$$
Since the group has an exponential growth, every random walk is transient and so the above summation is finite ($`G(x,y)`$ is actually bounded uniformly in $`x`$ and $`y`$).
For a random walk $`S()`$, let $`\tau _y`$ be the first hitting time of $`y`$:
$$\tau _y\stackrel{\text{def.}}{=}inf\{k0:S(k)=y\}.$$
When $`y`$ is never attained, let $`\tau _y=\mathrm{}`$. The hitting probability of $`y`$ starting at $`x`$ is denoted
$$F(x,y)\stackrel{\text{def.}}{=}^x[\tau _y<\mathrm{}].$$
We write $`G(x)=G(e,x)`$ and $`F(x)=F(e,x)`$. Note that $`F`$ and $`G`$ are symmetric and invariant by left multiplication. In particular, $`G(y,y)=G(e)`$. Therefore, the functions $`F`$ and $`G`$ are proportional since a straightforward computation shows
$$G(x,y)=G(y,y)F(x,y)=G(e)F(x,y).$$
(1)
We are now ready to define the distance we will use to study the Internal DLA. For all $`x,y\mathrm{\Gamma }`$, we define
$$d(x,y)\stackrel{\text{def.}}{=}\mathrm{ln}F(x,y).$$
Since $`F(x,y)`$ is the hitting probability of $`y`$ starting at $`x`$, we call $`d(,)`$ the hitting distance.
###### Lemma 2.1
$`d(,)`$ is a left invariant distance on $`\mathrm{\Gamma }`$.
Proof
As $`F(x,y)1`$, $`d(,)`$ is always non-negative.
The symmetry of the random walk and the left invariance of $`F(,)`$ yield the same properties for $`d(,)`$.
Remark also that since the random walk is symmetric and transient
$$xy,\mathrm{\hspace{0.33em}1}>^x[\tau _x^{}<\mathrm{}]^x[\tau _y<\mathrm{}]^y[\tau _x<\mathrm{}]=F(x,y)^2,$$
where $`\tau _x^{}\stackrel{\text{def.}}{=}inf\{k1:S(k)=x\}`$. Thus
$$d(x,y)=0F(x,y)=1x=y.$$
Finally,
$$^x[\tau _z<\mathrm{}]^x[\tau _y<\mathrm{}]^y[\tau _z<\mathrm{}]$$
leads to the triangular inequality. $`\mathrm{}`$
Let define
$$K\stackrel{\text{def.}}{=}\mathrm{max}\{d(e,x):x𝒮\}.$$
(2)
According to this distance, we define the ball of radius $`Kn`$,
$$B(x,n)\stackrel{\text{def.}}{=}\{y\mathrm{\Gamma }:d(x,y)Kn\}$$
and its external boundary
$$B(x,n)\stackrel{\text{def.}}{=}\{zB(x,n):yz\text{ and }yB(x,n)\}.$$
Here $`yz`$ means that $`y`$ and $`z`$ are neighbours, that is $`y^1z𝒮`$.
The reason why the constant $`K`$ appears in the definition of the balls is to be sure that $`B(e,n)B(e,n+1)`$. Indeed, let $`xB(e,n)`$ and $`yB(e,n)`$ such that $`yx`$. Then,
$$d(e,x)d(e,y)+d(y,x)Kn+\mathrm{max}\{d(e,z):z𝒮\}K(n+1).$$
We write $`d(x)=d(e,x)`$ and $`B(n)=B(e,n)`$. To have a better understanding of this distance, we will compare it to the word distance. First remark that, for any $`x\mathrm{\Gamma }`$, writing $`x=x_1\mathrm{}x_{|x|}`$ with all the $`x_i`$’s in $`𝒮`$, leads to
$$d(x)\underset{i=1}{\overset{|x|}{}}d(x_i)K|x|.$$
(3)
To get a lower bound for the hitting distance, we need asymptotic estimates of the Green function. These estimates are direct consequences of the following results on the asymptotic of the iterated transition probability:
1. An off-diagonal gaussian upper bound for symmetric random walks (with finite support) on generic groups (Varopoulos ; Carne ),
$$C_g>1\text{ s.t. }x,y\mathrm{\Gamma }\text{ and }k^x[S(k)=y]C_g\mathrm{exp}(C_g|x^1y|^2/k),$$
(4)
2. For random walks on groups having exponential growth (Varopoulos ; Carlen, Kusuoka & Stroock ),
$$C_e>1\text{ s.t. }x,y\mathrm{\Gamma }\text{ and }k^x[S(k)=y]C_e\mathrm{exp}(C_ek^{1/3}),$$
(5)
3. In the special case of non-ameanable groups, the latter estimates becomes (Kesten )
$$C_{na}>1\text{ s.t. }x,y\mathrm{\Gamma }\text{ and }k^x[S(k)=y]C_{na}\mathrm{exp}(C_{na}k).$$
(6)
For an overview of such estimates and precise references, see Woess or Pittet & Saloff-Coste .
###### Lemma 2.2
Let $`\mathrm{\Gamma }`$ be a group having exponential growth, then $`C>0`$ such that
$$G(x)\mathrm{exp}(C|x|^{1/2})x\mathrm{\Gamma }.$$
Proof
Using (4) and (5),
$`G(x)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}[S(k)=x]={\displaystyle \underset{k|x|^{3/2}}{}}[S(k)=x]+{\displaystyle \underset{k>|x|^{3/2}}{}}[S(k)=x]`$
$``$ $`{\displaystyle \underset{k|x|^{3/2}}{}}C_g\mathrm{exp}(C_g|x|^2/k)+{\displaystyle \underset{k>|x|^{3/2}}{}}C_e\mathrm{exp}(C_ek^{1/3})`$
$``$ $`C_g|x|^{3/2}\mathrm{exp}(C_g|x|^{1/2})+15|x|\mathrm{exp}(C_e|x|^{1/2})`$
$``$ $`\mathrm{exp}(C|x|^{1/2}),`$
for some constant $`C`$. The before last inequality comes from the fact that for any sufficiently large $`N`$
$`{\displaystyle \underset{k=N+1}{\overset{\mathrm{}}{}}}\mathrm{exp}(C_ek^{1/3})`$ $``$ $`{\displaystyle _N^+\mathrm{}}\mathrm{exp}(C_et^{1/3})𝑑t`$
$``$ $`{\displaystyle _N^+\mathrm{}}{\displaystyle \frac{d(6t^{2/3}\mathrm{exp}(C_et^{1/3}))}{dt}}𝑑t=6N^{2/3}\mathrm{exp}(C_eN^{1/3})`$
$`\mathrm{}`$
Actually, Revelle proved independently Lemma 2.2.
Note that when the group is non-ameanable, using (6), this upper bound of the Green function becomes:
$$G(x)\mathrm{exp}(C|x|).$$
(7)
The above lemma assures that $`G(x)`$ tends to $`0`$ as $`|x|`$ goes to infinity. By (1), so does $`F(x)`$. Lemma 2.2 leads also to
$$d(x)K_1|x|^{1/2},$$
(8)
for some constant $`K_1`$.
Actually, as soon as the Green function decreases exponentially (with respect to the word distance), these distances are equivalent. This is the case when $`\mathrm{\Gamma }`$ is non-ameanable (7) and also for some ameanable groups (for instance the Lamplighter group $`_2`$, Brofferio & Woess ).
Remark that (3) and (8) are sharp in the sense that examples exist where the Green function decreases as $`\mathrm{exp}(C|x|)`$ or $`\mathrm{exp}(C|x|^{1/2})`$ depending on the direction toward infinity (for instance $``$ , Revelle )
As $`B_w(n)B(n)`$, the size $`V(n)`$ of the balls $`B(n)`$ grows exponentially, but we prove that the constant in the exponential (lower and upper) bound is exactly the constant $`K`$ which appeared in the definition of the balls.
###### Proposition 2.3
There exist two constants $`C_a`$ and $`C_b`$ such that for every integer $`n`$,
$$C_a\mathrm{exp}(Kn)V(n)C_bn^3\mathrm{exp}(Kn).$$
Proof
The proof of the upper bound relies on the following estimate taken from Revelle : for any finite set $`A`$,
$$\underset{xA}{}G(x)(16C_e^2+1)(\mathrm{ln}\mathrm{\#}A)^3.$$
(9)
By sake of completeness, we give the computation that leads to (9). The estimate (5) (which introduces the constant $`C_e`$) yields
$`{\displaystyle \underset{xA}{}}G(x)`$ $``$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{xA}{}}[S(k)=x]`$
$`=`$ $`{\displaystyle \underset{kC_e^3(\mathrm{ln}\mathrm{\#}A)^3+1}{}}{\displaystyle \underset{xA}{}}[S(k)=x]+{\displaystyle \underset{k>C_e^3(\mathrm{ln}\mathrm{\#}A)^3+1}{}}{\displaystyle \underset{xA}{}}[S(k)=x]`$
$``$ $`C_e^3(\mathrm{ln}\mathrm{\#}A)^3+1+\mathrm{\#}A{\displaystyle \underset{k>C_e^3(\mathrm{ln}\mathrm{\#}A)^3+1}{}}C_e\mathrm{exp}(C_ek^{1/3})`$
$``$ $`(16C_e^2+1)(\mathrm{ln}\mathrm{\#}A)^3.`$
When $`A=B(n)`$, the definition of the ball $`B(n)`$ and (9) give
$$\mathrm{\#}B(n)\mathrm{exp}(Kn)\underset{xB(n)}{}G(x)(16C_e^2+1)(\mathrm{ln}\mathrm{\#}B(n))^3$$
(10)
which implies, since $`(\mathrm{ln}\mathrm{\#}B(n))^3(\mathrm{\#}B(n))^{1/2}`$, that $`\mathrm{\#}B(n)`$ is at most an exponential of $`n`$. Applying once more (10) leads to the $`C_bn^3\mathrm{exp}(Kn)`$ upper bound, for some positive constant $`C_b`$.
For the lower bound, remark that
$$\underset{xB(n1)}{}G(x)(\mathrm{\#}B(n1))\mathrm{exp}(K(n1)).$$
The left hand side is larger than (or equal to) $`1`$ since the random walk is transient and it must go through $`B(n1)`$. So with $`C_a\stackrel{\text{def.}}{=}\mathrm{exp}(K)`$,
$$V(n)\mathrm{\#}B(n1)C_a\mathrm{exp}(Kn).$$
$`\mathrm{}`$
###### Remark 2.4
For non-ameanable groups, the upper bound in the volume growth becomes $`C_bn\mathrm{exp}(Kn)`$. Indeed, using the specific estimate of the return probability (6), the upper bound in (9) becomes a constant times $`\mathrm{ln}\mathrm{\#}A`$.
## 3 Limiting shape for the Internal DLA
Let us now state our main result, a comparison between the shape of the Internal DLA and the balls for the hitting distance.
###### Theorem 3.1
Let $`A()`$ be the Internal DLA on a finitely generated group of exponential growth, associated to a symmetric random walk with finitely supported increment. Then, for any constants $`C_I>3/K`$ and $`C_O>2`$,
$$[n_0\text{ s.t. }n>n_0,B(nC_I\mathrm{ln}n)A(V(n))B(n+C_O\sqrt{n})]=1.$$
This result give an upper bound for the two random variables $`\delta _I(n)`$ and $`\delta _O(n)`$ describing the inner and outer fluctuation of $`A(n)`$ with respect to $`B(n)`$
$`\delta _I(n)`$ $`\stackrel{\text{def.}}{=}`$ $`ninf\{d(z)/K:zA(V(n))\},`$
$`\delta _O(n)`$ $`\stackrel{\text{def.}}{=}`$ $`sup\{d(z)/K:zA(V(n))\}n.`$
An alternative definition is that $`B(n\delta _I(n))`$ is the largest ball included in $`A(V(n))`$, and $`B(n+\delta _O(n))`$ is the smallest ball which contains $`A(V(n))`$. Theorem 3.1, states that for any constants $`C_I>3/K`$ and $`C_O>2`$, the events $`\{\delta _I(n)>C_I\mathrm{ln}n\}`$ and $`\{\delta _O(n)>C_O\sqrt{n}\}`$ occur only finitely often, with probability $`1`$.
We split the proof into two parts, one regarding the inner error and the other regarding the outer error.
Let us first give some definitions used throughout the proof. Let $`S^j()`$ be the $`\mathrm{j}^{\mathrm{th}}`$ random walk that defines the Internal DLA $`A(V(n))`$. Recall that all the $`S^j()`$ start from $`e`$. Moreover, we let this random walk running even after it adds a new point to the cluster $`A(j1)`$. This adding time will then be the stopping time:
$$\sigma ^j\stackrel{\text{def.}}{=}inf\{k0:S^j(k)A(j1)\}.$$
We also define the hitting time of a point $`z`$ by the (unstopped) random walk $`S^j`$:
$$\tau _z^j\stackrel{\text{def.}}{=}inf\{k0:S^j(k)=z\}.$$
### 3.1 Proof of the inner part
We will prove that there exists a constant $`C_I`$ such that, with probability $`1`$, for all $`n`$ large enough
$$B(nC_I\mathrm{ln}n)A(V(n)).$$
We start our proof by computing an upper bound for the probability that a given point $`zB(n)`$ does not belong to the Internal DLA at time $`V(n)`$. For that purpose, we define the following indicator random variable
$$N^j=N^j(z)\stackrel{\text{def.}}{=}\text{ 1 I}_{\{\tau _z^j\sigma ^j\}},$$
which can be decomposed as $`N_j=M_jL_j`$ where
$$M^j=M^j(z)\stackrel{\text{def.}}{=}\text{ 1 I}_{\{\tau _z^j<\mathrm{}\}}\text{and}L^j=L^j(z)\stackrel{\text{def.}}{=}\text{ 1 I}_{\{\sigma ^j\tau _z^j<\mathrm{}\}}.$$
Observe that $`z`$ does not belong to $`A(V(n))`$ if and only if $`N^j=0`$ for every $`j`$. So, for any $`\lambda >0`$,
$`[zA(V(n))]`$ $`=`$ $`\left[{\displaystyle \underset{j=1}{\overset{V(n)}{}}}N^j=0\right]𝔼\left[\mathrm{e}^{\lambda _{j=1}^{V(n)}N^j}\right]`$
$`=`$ $`𝔼\left[\mathrm{e}^{\lambda _{j=1}^{V(n)}M^j+\lambda _{j=1}^{V(n)}L^j}\right]`$
$``$ $`𝔼\left[\mathrm{e}^{2\lambda _{j=1}^{V(n)}M^j}\right]^{1/2}\times 𝔼\left[\mathrm{e}^{2\lambda _{j=1}^{V(n)}L^j}\right]^{1/2}.`$
As the $`M^j`$’s are i.i.d. indicator random variables,
$`𝔼\left[\mathrm{e}^{2\lambda _{j=1}^{V(n)}M^j}\right]`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{V(n)}{}}}𝔼\left[\mathrm{e}^{2\lambda M^j}\right]={\displaystyle \underset{j=1}{\overset{V(n)}{}}}(1(1\mathrm{e}^{2\lambda })F(z))`$
$`=`$ $`(1(1\mathrm{e}^{2\lambda })F(z))^{V(n)}\mathrm{exp}((1\mathrm{e}^{2\lambda })F(z)V(n)).`$
Observe that the $`L^j`$’s are not independent, but they can be bounded from above by random variables which are independent when knowing where the random walk $`S^j`$ adds to the cluster. More precisely, let define the hitting time of $`z`$ after the adding time $`\sigma ^j`$,
$$\stackrel{~}{\tau _z}^j\stackrel{\text{def.}}{=}inf\{k\sigma ^j\mathrm{s}.\mathrm{t}.S^j(k)=z\}.$$
Then,
$$L^j\stackrel{~}{L}^j\stackrel{\text{def.}}{=}\text{ 1 I}_{\{\stackrel{~}{\tau _z}^j<\mathrm{}\}}.$$
Let $`𝒢_n`$ be the $`\sigma `$-algebra generated by $`S^j(k\sigma ^j)`$ for all $`jV(n)`$ and $`k`$, that is the $`\sigma `$-algebra containing all the information on the $`V(n)`$ random walks before they add to the cluster. As the $`\stackrel{~}{L}^j`$’s are independent conditioning by $`𝒢_n`$,
$`𝔼\left[\mathrm{e}^{2\lambda _{j=1}^{V(n)}L^j}\right]`$ $``$ $`𝔼\left[\mathrm{e}^{2\lambda _{j=1}^{V(n)}\stackrel{~}{L}^j}\right]=𝔼\left[𝔼\left[\mathrm{e}^{2\lambda _{j=1}^{V(n)}\stackrel{~}{L}^j}|𝒢_n\right]\right]`$
$`=`$ $`𝔼\left[{\displaystyle \underset{j=1}{\overset{V(n)}{}}}𝔼\left[\mathrm{e}^{2\lambda \stackrel{~}{L}^j}|𝒢_n\right]\right].`$
Since the $`\stackrel{~}{L}^j`$’s are indicator random variables,
$$𝔼\left[\mathrm{e}^{2\lambda _{j=1}^{V(n)}L^j}\right]𝔼\left[\underset{j=1}{\overset{V(n)}{}}(1+(\mathrm{e}^{2\lambda }1)𝔼[\stackrel{~}{L}^j|𝒢_n])\right]𝔼\left[\mathrm{e}^{(\mathrm{e}^{2\lambda }1)_{j=1}^{V(n)}𝔼[\stackrel{~}{L}^j|𝒢_n]}\right].$$
Finally,
$`{\displaystyle \underset{j=1}{\overset{V(n)}{}}}𝔼\left[\stackrel{~}{L}^j|𝒢_n\right]`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{V(n)}{}}}𝔼\left[\text{ 1 I}_{\{\stackrel{~}{\tau _z}^j<\mathrm{}\}}|𝒢_n\right]={\displaystyle \underset{j=1}{\overset{V(n)}{}}}^{S^j(\sigma ^j)}[\stackrel{~}{\tau _z}^j<\mathrm{}]`$
$`=`$ $`{\displaystyle \underset{yA(V(n))}{}}F(y,z)={\displaystyle \underset{yA(V(n))}{}}F(z,y).`$
The definition of the hitting distance implies that for any $`yB(z,n)`$ and any $`y^{}B(z,n)`$, $`F(z,y)>F(z,y^{})`$. Thus, since $`\mathrm{\#}A(V(n))=V(n)=\mathrm{\#}B(z,n)`$,
$$\underset{yA(V(n))}{}F(z,y)\underset{yB(z,n)}{}F(z,y)=\underset{yB(n)}{}F(y).$$
(11)
Putting things together, we get that for any $`\lambda >0`$, $`n`$ and $`zB(n)`$,
$$[zA(V(n))]\mathrm{exp}\left(C_\lambda V(n)F(z)+C_\lambda ^{}\underset{yB(n)}{}F(y)\right).$$
(12)
with $`C_\lambda =\frac{1\mathrm{e}^{2\lambda }}{2}`$ and $`C_\lambda ^{}=\frac{\mathrm{e}^{2\lambda }1}{2}`$.
We remark that the inequality (12) is valid for every transient groups whatever its growth rate. For groups of exponential growth, Proposition 2.3 enables to estimate the last term in the exponential. Thus, if $`zB(nC_I\mathrm{ln}n)`$ (that is $`F(z)n^{KC_I}\mathrm{e}^{Kn}`$), using (9),
$$[zA(V(n))]\mathrm{e}^{C_\lambda V(n)n^{KC_I}\mathrm{e}^{Kn}+cC_\lambda ^{}(\mathrm{ln}V(n))^3}\mathrm{e}^{C_\lambda C_an^{KC_I}+c^{}C_\lambda ^{}n^3}\mathrm{e}^{C_\lambda C_an^3},$$
for $`C_I>3/K`$ and $`n`$ large enough. Then,
$`{\displaystyle \underset{zB(nC_I\mathrm{ln}n)}{}}[zA(V(n))]`$ $``$ $`V(nC_I\mathrm{ln}n)\mathrm{e}^{C_\lambda C_an^3}`$
$``$ $`C_b(nC_I\mathrm{ln}n)^3\mathrm{e}^{KnC_I\mathrm{ln}nC_\lambda C_an^3}`$
which is summable. Thus the result follows from Borel-Cantelli Lemma.
### 3.2 Proof of the outer part
To get the upper bound for the outer error in Theorem 3.1, we will estimate the decay of the expected number of points in $`A(V(n))`$ that lies within the spheres $`B(n+p)`$ for $`p1`$. Checking that this decay is exponential in $`p^2`$ will lead to the desired result.
We define
$$Z_p(j)\stackrel{\text{def.}}{=}\mathrm{\#}(A(j)B(n+p))=\underset{i=1}{\overset{j}{}}\text{ 1 I}_{S^i(\sigma ^i)B(n+p)},$$
and $`\nu _p(j)\stackrel{\text{def.}}{=}𝔼[Z_p(j)]`$. If a random walk adds to the cluster in $`B(n+p+1)`$, that implies it has stayed within the cluster before. It means that
$$\{S^i(\sigma ^i)B(n+p+1)\}\{xB(n+p)A(i1)\mathrm{s}.\mathrm{t}.\tau _x^i<\mathrm{}\}.$$
Therefore,
$`\nu _{p+1}(j)`$ $``$ $`{\displaystyle \underset{i=1}{\overset{j}{}}}[xB(n+p)A(i1)\mathrm{s}.\mathrm{t}.\tau _x^i<\mathrm{}]`$
$``$ $`{\displaystyle \underset{i=1}{\overset{j}{}}}𝔼\left[{\displaystyle \underset{xB(n+p)A(i1)}{}}\text{ 1 I}_{\tau _x^i<\mathrm{}}\right]`$
$``$ $`{\displaystyle \underset{i=1}{\overset{j}{}}}𝔼\left[Z_p(i1)\underset{xB(n+p)}{\mathrm{max}}[\tau _x<\mathrm{}]\right],`$
by independence between $`S^i`$ and $`A(i1)`$. Then,
$$\nu _{p+1}(j)\underset{i=1}{\overset{j}{}}\mathrm{exp}(K(n+p))\nu _p(i1).$$
(13)
For any fixed $`j`$, iterating (13) leads to
$$\nu _p(j)\mathrm{exp}(Kn(p1))\mathrm{exp}(Kp(p1)/2)\frac{j^p}{p!},$$
by a simple induction on $`p`$ started from $`\nu _1(j)j`$ and using
$$\underset{i=1}{\overset{j}{}}\frac{(i1)^p}{p!}\frac{j^{p+1}}{(p+1)!}.$$
Finally, as $`p!p^p\mathrm{exp}(p)`$,
$$\nu _p(j)\mathrm{exp}(Kn(p1))\mathrm{exp}(Kp(p1)/2)\mathrm{exp}(p)p^pj^p.$$
Thus, for $`j=V(n)`$ and $`p`$ larger than some constant independent of $`n`$,
$$\nu _p(V(n))n^{3p}\mathrm{exp}\left(K\left(\frac{p^2}{3}n\right)\right).$$
Therefore, for $`C_O>2`$ and $`p=C_O\sqrt{n}`$,
$`[A(V(n))B(n+C_O\sqrt{n})]`$ $``$ $`[Z_{C_O\sqrt{n}}(V(n))1]`$
$``$ $`\nu _{C_O\sqrt{n}}(V(n))`$
$``$ $`\mathrm{exp}\left(K\left({\displaystyle \frac{C_O^2}{4}}1\right)n\right),`$
for $`n`$ large enough.
Thus, by the Borel-Cantelli Lemma, with probability $`1`$, for all $`n`$ large enough,
$$A(V(n))B(n+C_O\sqrt{n}).$$
## 4 Lower bounds for the errors when $`𝒢`$ is the homogeneous tree $`𝕋_q`$
In this section we will give lower bounds for the inner and outer fluctuations in the specific case of Simple Random Walks on the homogeneous tree.
Let $`𝒢`$ be the homogeneous tree $`𝕋_q`$ of degree $`q3`$. When $`q`$ is even, $`𝕋_q`$ is the Cayley graph of the free group $`F_{q/2}`$. When $`q`$ is odd, $`𝕋_q`$ is the Cayley graph of the free product of $`F_{(q1)/2}`$ by $`_2`$. On such a graph, most of the estimates we need are known exactly (Dynkin & Maljutov and Cartier ) for the Simple Random Walk (nearest neighbours with uniform transition probability). In particular since
$$F(z)=(q1)^{|z|},$$
the hitting distance and the word distance are proportional and the ball $`B_w(n)`$ coincides with $`B(n)`$ exactly thanks to the use of the constant $`K=\mathrm{ln}(q1)`$. Thus we have
$$V(n)=\frac{q(q1)^n2}{q2}\text{ and }\mathrm{\#}B(n)=q(q1)^n.$$
Let
$$\xi _n\stackrel{\text{def.}}{=}\mathrm{min}\{k0\mathrm{s}.\mathrm{t}.S(k)B(n)\},$$
be the exit time from the ball $`B(n)`$. Then, as we deal with Simple Random Walks,
$$[S(\xi _n)=z]=\frac{1}{q(q1)^n}\delta _{zB(n)}.$$
### 4.1 Lower bound for the outer error
###### Proposition 4.1
For any constant $`C_o<(\mathrm{ln}(q1)/\mathrm{ln}q)^{1/2}`$,
$$[A(V(n))B(n+C_o\sqrt{n})=\mathrm{}\text{infinitely often}]=0$$
Proof
By Theorem 3.1, for any constant $`C_I`$ sufficiently large, if one defines the increasing sequence of events
$$\mathrm{\Omega }_{n_0}\stackrel{\text{def.}}{=}\left\{nn_0,B(nC_I\mathrm{ln}(n)+1)A(V(n))\right\}.$$
One has
$$\underset{n_0\mathrm{}}{lim}[\mathrm{\Omega }_{n_0}]=1.$$
(14)
Suppose that $`\mathrm{\Omega }_{n_0}`$ occurs and fix $`n>n_0`$. Denote $`R\stackrel{\text{def.}}{=}n1C_I\mathrm{ln}(n1)`$. Let the first $`V(n1)`$ random walks build $`A(V(n1))`$, which covers $`B(R)`$. Then we consider the next $`V(n)V(n1)`$ random walks and stop them as soon as they reach $`B(R)`$. Let $`N_x`$ be the number of random walks that are stopped in $`xB(R)`$. Letting the stopped random walks start again we can reconstruct the original Internal DLA. The same random walks construct new smaller Internal DLA’s, $`A_x(N_x)`$, that start in every $`xB(R)`$ and are built up using $`N_x`$ random walks . We state that if $`B(R)A(V(n1))`$ then
$$A(V(n))A(V(n1))\underset{xB(R)}{}A_x(N_x).$$
Thus, if there exists $`xB(R)`$ such that $`A_x(N_x)B(R+r)\mathrm{}`$, then $`A(V(n))B(R+r)\mathrm{}`$. Observe also that the random walks that start on $`B(R)`$ are independent and, thus, all the $`\left\{A_x(N_x)\right\}_{xB(R)}`$ are independent once one knows all the $`N_x`$, that is they are independent for the probability $`[|\sigma _R]`$. Here $`\sigma _R`$ is the $`\sigma `$-algebra generated by the $`S^i(k\xi _R)`$ for all $`iN\stackrel{\text{def.}}{=}V(n)V(n1)`$ and all $`k`$, that is the $`\sigma `$-algebra containing all the information on the $`N`$ random walks until they hit $`B(R)`$.
For all $`xB(R),`$ there exits a path $`𝒞_x`$ of (graph) length $`r`$ from $`x`$ to $`B(R+r)`$. Note that
$$\{A_x(N_x)B(R+r)=\mathrm{}\}\{𝒞_xA_x(N_x)\}.$$
We split the $`N_x`$ random walks that start at $`x`$ into $`\frac{N_x}{r}`$ packets of size $`r`$ ($``$ is the lower integer part), the rest being ignored. The probability that all the random walks of one packet follow exactly $`𝒞_x`$ when they start, is $`(q^r)^r`$. By definition of the Internal DLA, this event guaranty that $`𝒞_x`$ is contained in $`A_x(N_x)`$. Thus, since all packets are independent
$$[𝒞_xA_x(N_x)|\sigma _R](1q^{r^2})^{\frac{N_x}{r}}(1q^{r^2})^{\frac{N_x}{r}1}.$$
Finally, for $`n>n_0`$,
$`[\{A(n)B(R+r)=\mathrm{}\}`$ $``$ $`\mathrm{\Omega }_{n_0}]`$
$``$ $`[xB(R),A_x(N_x)B(R+r)=\mathrm{}]`$
$`=`$ $`𝔼\left[{\displaystyle \underset{xB(R)}{}}[A_x(N_x)B(R+r)=\mathrm{}|\sigma _R]\right]`$
$``$ $`𝔼\left[{\displaystyle \underset{xB(R)}{}}(1q^{r^2})^{\frac{N_x}{r}1}\right]`$
$`=`$ $`𝔼\left[(1q^{r^2})^{_{xB(R)}\frac{N_x}{r}1}\right]`$
$`=`$ $`(1q^{r^2})^{\frac{V(n)V(n1)}{r}\mathrm{\#}B(R)}`$
$``$ $`\mathrm{exp}\left(q^{r^2}\left({\displaystyle \frac{V(n)V(n1)}{r}}\mathrm{\#}B(R)\right)\right).`$
Remark that taking $`r=C_o\sqrt{n}+1+C_I\mathrm{ln}(n1)`$, with a constant $`C_o<(\mathrm{ln}(q1)/\mathrm{ln}q)^{1/2}`$ and $`C_I>\frac{1}{2}`$, we get
$$\mathrm{exp}\left(q^{r^2}\left(\frac{V(n)V(n1)}{r}\mathrm{\#}B(R)\right)\right)\mathrm{exp}(\mathrm{exp}(cn)),$$
for some constant $`c>0`$. Then, by Borel-Cantelli Lemma,
$$[\{A(V(n))B(n+C_o\sqrt{n})\mathrm{i}.\mathrm{o}.\}\mathrm{\Omega }_{n_0}]=0.$$
Letting $`n_0`$ goes to infinity leads (by (14)) to the result. $`\mathrm{}`$
### 4.2 Lower bound for the inner error
For the inner error, we prove a weaker result than the one on the outer error. Namely, we get that infinitely often, the inner error is of order $`\mathrm{ln}n`$. It assures the sharpness of the order $`\mathrm{ln}n`$ for a general upper bound but does not prevent the inner error to be sometimes smaller. Thus our result is somehow weaker than the one announced by Etheridge and Lawler.
###### Proposition 4.2
For any constant $`C_i<1/(2\mathrm{ln}(q1))`$,
$$[A(V(n))^cB(nC_i\mathrm{ln}n)\mathrm{}\text{infinitely often}]=1.$$
Proof
Let $`N\stackrel{\text{def.}}{=}V(n)`$ and $`R\stackrel{\text{def.}}{=}n1C_i\mathrm{ln}n`$. Then the event $`\{A(N)^cB(R+1)\mathrm{}\}`$ contains the event that there exists a point on the sphere of radius $`R`$ attained by none of $`N`$ independent random walks.
Let fix $`r<R`$, for every $`zB(R)`$, we define its class as
$$[z]\stackrel{\text{def.}}{=}B(R)B(z,2r),$$
and let $`A_R\stackrel{\text{def.}}{=}_{i=1}^N[S^i(\xi _R)]`$ be the set of all classes that are hit by the $`N`$ random walks when they first exit $`B(R)`$. Then we have
$`[A(N)^c`$ $``$ $`B(R)\mathrm{}]`$
$``$ $`[zB(R)\mathrm{s}.\mathrm{t}.iN,\tau _z^i=\mathrm{}]`$
$``$ $`[zB(R)\mathrm{s}.\mathrm{t}.iN,\tau _z^i=\mathrm{}\mathrm{and}A_RB(R)]]`$
$`=`$ $`{\displaystyle \underset{AB(R)}{}}[zB(R)\mathrm{s}.\mathrm{t}.iN,\tau _z^i=\mathrm{}|A_R=A][A_R=A].`$
To get a lower bound for the above conditional probability we will need, once the set $`A`$ is fixed, to choose a suitable point called $`z_A`$ on $`B(R)`$ that has few chances to be hit by the random walks after $`\xi _R^i`$. For that purpose, we construct the following algorithm. Suppose $`AB(R)`$ is a fixed set, then we will write a “Mouse algorithm” whose aim is to find the most secure place $`z_A`$ on $`B(R)`$ for a mouse knowing that $`\mathrm{\#}A`$ cats will start simple random walks (hunts) from each of the points in $`A`$.
The mouse starts from the origin and proceeds toward $`B(R)`$ choosing the direction of the subtree with the fewest cats. Denote $`(m_j)`$ ($`j=0`$ to $`R`$) the successive positions of the mouse, that is the steps of our algorithm:
* $`m_0=e`$;
* Choose $`m_{j+1}`$ inside $`\{ym_j:|y|=j+1`$ and $`B(y,Rj)A`$ is minimal$`\}`$.
Let $`z_A\stackrel{\text{def.}}{=}m_{R+1}`$. It is easy to check that $`z_AA`$. Recall that $`\sigma _R`$ is the $`\sigma `$-algebra generated by the $`S^i(k\xi _R)`$ for all $`iN`$ and all $`k`$, that is the $`\sigma `$-algebra containing all the information on the $`N`$ random walks before they hit $`B(R)`$. Then,
$`[zB(R)`$ $`\mathrm{s}.\mathrm{t}.`$ $`iN,\tau _z^i=\mathrm{}|A_R=A]`$
$``$ $`[iN,\tau _{z_A}^i=\mathrm{}|A_R=A]`$
$`=`$ $`[iN,\xi _R^i\tau _{z_A}^i=\mathrm{}|A_R=A]`$
$`=`$ $`𝔼\left[[iN,\xi _R^i\tau _{z_A}^i=\mathrm{}|\sigma _R]|A_R=A\right]`$
$`=`$ $`𝔼\left[{\displaystyle \underset{i=1}{\overset{N}{}}}[\xi _R^i\tau _{z_A}^i=\mathrm{}|\sigma _R]|A_R=A\right]`$
$`=`$ $`𝔼\left[{\displaystyle \underset{i=1}{\overset{N}{}}}^{S^i(\xi _R)}[\tau _{z_A}=\mathrm{}]|A_R=A\right].`$
By transience of the random walk, $`^{S^i(\xi _R)}[\tau _{z_A}=\mathrm{}]=1F(S^i(\xi _R),z_A)`$ has a uniform (in all parameter) strickly positive lower bound. Thus, there exists a constant $`c>0`$ such that
$`𝔼\left[{\displaystyle \underset{i=1}{\overset{N}{}}}^{S^i(\xi _R)}[\tau _{z_A}=\mathrm{}]|A_R=A\right]`$ $``$ $`𝔼\left[\mathrm{exp}\left(c{\displaystyle \underset{i=1}{\overset{N}{}}}F(S^i(\xi _R),z_A)\right)|A_R=A\right]`$
$``$ $`\mathrm{exp}\left(c{\displaystyle \underset{i=1}{\overset{N}{}}}𝔼\left[F(S^i(\xi _R),z_A)|A_R=A\right]\right)`$
$`=`$ $`\mathrm{exp}\left(cN𝔼\left[F(S^1(\xi _R),z_A)|A_R=A\right]\right).`$
By definition of the Mouse algorithm, we easily check that, for every $`j`$,
$$T_j\stackrel{\text{def.}}{=}B(z_A,2(R+1j))A=B(m_j,R+1j)A,$$
and that its cardinality is $`\mathrm{\#}T_j\mathrm{\#}A(q1)^j`$. Remark also that since $`A`$ must be a union of sets of the type $`B(R)B(z,2r)`$ with $`zB(R)`$ (otherwise $`[A_R=A]=0`$), $`z_AA`$ implies $`d(z_A,A)>2r`$. Then,
$`𝔼\left[F(S^1(\xi _R),z_A)|A_R=A\right]`$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{Rr}{}}}𝔼\left[\text{ 1 I}_{S^1(\xi _R)T_j\backslash T_{j+1}}F(S^1(\xi _R),z_A)|A_R=A\right]`$
$`=`$ $`{\displaystyle \underset{j=0}{\overset{Rr}{}}}𝔼\left[\text{ 1 I}_{S^1(\xi _R)T_j\backslash T_{j+1}}(q1)^{2(R+1j)}|A_R=A\right]`$
$`=`$ $`{\displaystyle \underset{j=0}{\overset{Rr}{}}}(q1)^{2(R+1j)}\left[S^1(\xi _R)T_j\backslash T_{j+1}|A_R=A\right].`$
Remark that the event $`A_R=A`$ is invariant by any bijection $`\gamma `$ of $`B(R)`$ such that $`\gamma (A)=A`$ and for any $`xB(R)`$, $`[\gamma (x)]=\gamma ([x])`$. Thus, the law of $`S^1(\xi _R)`$ under the condition $`A_R=A`$ is the uniform distribution on $`A`$. Thus,
$$\left[S^1(\xi _R)T_j\backslash T_{j+1}|A_R=A\right]=\frac{\mathrm{\#}(T_j\backslash T_{j+1})}{\mathrm{\#}A}\frac{\mathrm{\#}T_j}{\mathrm{\#}A}.$$
Therefore,
$$𝔼\left[F(S^1(\xi _R),z_A)|A_R=A\right]\underset{j=0}{\overset{Rr}{}}\frac{(q1)^{2(R+1j)}}{(q1)^j}(q1)^{Rr}.$$
Finally,
$`[zB(R)\mathrm{s}.\mathrm{t}.iN,\tau _z^i=\mathrm{}]`$ $``$ $`{\displaystyle \underset{AB(R)}{}}\mathrm{exp}\left(cN(q1)^{Rr}\right)[A_R=A]`$
$`=`$ $`\mathrm{exp}\left(cN(q1)^{Rr}\right)[A_RB(R)].`$
Remark that $`[A_RB(R)]`$ is the probability that $`N`$ balls uniformly distributed into $`M=q(q1)^{Rr}`$ boxes leave at least one empty box. By Lemme 4.3, that will be prove at the end of the section,
$$[A_RB(R)]1\mathrm{exp}\left(\frac{q(q1)^{Rr}}{2}\mathrm{exp}\left(4N\frac{(q1)^{rR}}{q}\right)\right).$$
When $`N=V(n)q(q1)^n`$, $`R=nC_i\mathrm{ln}n`$ and $`r=C_1\mathrm{ln}n`$ (the constant $`C_1`$ will be fixed below), we easily check that
$$\begin{array}{c}\mathrm{exp}\left(\frac{q(q1)^{Rr}}{2}\mathrm{exp}\left(4N\frac{(q1)^{rR}}{q}\right)\right)\hfill \\ \hfill \mathrm{exp}\left(\frac{qn^{(C_1+C_i)\mathrm{ln}(q1)}}{2}\mathrm{exp}\left(n\mathrm{ln}(q1)4n^{(C_1+C_i)\mathrm{ln}(q1)}\right)\right).\end{array}$$
So, with $`C_1=\frac{1}{2\mathrm{ln}(q1)}`$, for any constant $`C_i<\frac{1}{2\mathrm{ln}(q1)}`$,
$$[A_RB(R)]\underset{n\mathrm{}}{}1.$$
With the same constants $`C_1`$ and $`C_i`$, as $`C_1>C_i`$,
$$\mathrm{exp}\left(cN(q1)^{Rr}\right)\mathrm{exp}\left(cqn^{(C_1C_i)\mathrm{ln}(q1)}\right),$$
also tends to $`1`$ when $`n`$ goes to infinity. Finally,
$$[zB(nC_i\mathrm{ln}n)\mathrm{s}.\mathrm{t}.iV(n),\tau _z^i=\mathrm{}]\underset{n\mathrm{}}{}1.$$
Thus $`[A(V(n))^cB(nC_i\mathrm{ln}n)\mathrm{}]1`$ which completes the proof by Fatou Lemma. $`\mathrm{}`$
We now prove the estimate on the multinomial law that we used in the proof.
###### Lemma 4.3
Let put independently $`N`$ balls into $`M`$ boxes ($`M3`$) with uniform probability. Let $`_{N,M}`$ denote the probability law of that process. Then,
$$_{N,M}[\text{Each box contains at least one ball}]\mathrm{exp}\left(\frac{M}{2}\mathrm{exp}\left(4\frac{N}{M}\right)\right).$$
Proof
Let $`(b_k)`$ ($`k=1`$ to $`M`$) denote the $`M`$ boxes and $`(a_i)`$ ($`i=1`$ to $`N`$) the $`N`$ balls. Let also define
$`p_{N,M}`$ $`\stackrel{\text{def.}}{=}`$ $`_{N,M}[kM,iN\mathrm{s}.\mathrm{t}.a_ib_k];`$
$`q_{N,M}`$ $`\stackrel{\text{def.}}{=}`$ $`_{N,M}[iN\mathrm{s}.\mathrm{t}.a_ib_1];`$
$`I_1`$ $`\stackrel{\text{def.}}{=}`$ $`\{i\mathrm{s}.\mathrm{t}.a_ib_1\}.`$
Therefore,
$`p_{N,M}`$ $`=`$ $`{\displaystyle \underset{I\mathrm{}}{}}_{N,M}[\{k1iI\mathrm{s}.\mathrm{t}.a_ib_k\}I_1=I]`$
$`=`$ $`{\displaystyle \underset{I\mathrm{}}{}}_{N,M}[k1iI\mathrm{s}.\mathrm{t}.a_ib_k|I_1=I]\times _{N,M}[I_1=I]`$
$`=`$ $`{\displaystyle \underset{I\mathrm{}}{}}p_{N\mathrm{\#}I,M1}\times _{N,M}[I_1=I]`$
$``$ $`{\displaystyle \underset{I\mathrm{}}{}}p_{N,M1}\times _{N,M}[I_1=I]`$
$`=`$ $`p_{N,M1}q_{N,M}.`$
The third equality comes from the fact that conditioning by $`I_1=I`$, the $`N\mathrm{\#}I`$ remaining balls are uniformly distributed among the $`M1`$ remaining boxes. The inequality comes from the monotonicity in $`N`$ of $`p_{N,M}`$. By induction,
$$p_{N,M}\underset{j=2}{\overset{M}{}}q_{N,j}p_{N,1}.$$
As $`p_{N,1}=1`$ and $`q_{N,j}=1\left(1\frac{1}{j}\right)^N`$,
$`p_{N,M}`$ $``$ $`{\displaystyle \underset{j=(M+1)/2}{\overset{M}{}}}q_{N,j}{\displaystyle \underset{j=(M+1)/2}{\overset{M}{}}}\left(1\left(1{\displaystyle \frac{2}{M}}\right)^N\right)`$
$``$ $`\left(1\left(1{\displaystyle \frac{2}{M}}\right)^N\right)^{M/2}\left(1\mathrm{exp}\left(c{\displaystyle \frac{2N}{M}}\right)\right)^{M/2}`$
$``$ $`\mathrm{exp}\left({\displaystyle \frac{M}{2}}\mathrm{exp}\left(c{\displaystyle \frac{2N}{M}}\right)\right),`$
for some constant $`c>0`$ ($`c=2`$ for instance, since $`M3`$). $`\mathrm{}`$
Sébastien Blachère
Université Aix-Marseille 1
CMI, 39 rue Joliot-Curie
13453 Marseille Cedex, France
Sara Brofferio
Université Paris-Sud
Laboratoire de Mathématiques, bât. 425,
91405 Orsay Cedex, France. |
warning/0507/astro-ph0507561.html | ar5iv | text | # POWER SPECTRUM AND INTERMITTENCY OF Ly𝛼 TRANSMITTED FLUX OF QSO HE2347-4342
## 1 INTRODUCTION
The standard LCDM cosmogony is gaining more and more support from observations of cosmic structures, including the temperature fluctuations of cosmic background radiation, the clustering of galaxies and clusters, and the transmission flux fields of QSOs’ Ly$`\alpha `$ absorption spectrum. In the linear regime, the power spectrum of mass density perturbations predicted by the LCDM model is found to be consistent with observations on scales from few thousand to about 1 h<sup>-1</sup> Mpc (Pope et al. 2004). In the non-linear regime, N-body simulations of the standard LCDM model reveal that the mass density profile of dark matter halos is probably universal, and the cosmic mass field can be modeled as a superposition of the universal halos on various mass scales (e.g. Cooray & Sheth, 2002). This universal halo scenario has been successful in describing the second and higher order correlations of the evolved mass fields. It has also been extensively applied to model the formation and evolution of galaxies using the mass function, the universal density profile, the two-point correlation of the host halos, and the bias model of the relevant objects.
However, whether the LCDM model explains observations on sub-Mpc scales is still unclear. The universal density profile of dark matter halos given by LCDM N-body simulation is cuspy or singular, i.e., with a central density distribution given by $`\rho (r)r^\alpha `$, where $`\alpha =1`$ (Hernquist, 1990; Navarro et al. 1996; Jing 2000), or $`3/2`$ (Moore et al. 1999). However, the mass density profile given by the rotation curves of dwarf and low surface brightness galaxies generally show a soft core in their centers. The best fitted profiles are not in general as dense in the predicted cuspy center (Flores & Primack 1994; Swaters et al 2003; McGaugh et al. 2003; Zentner & Bullock 2003; Simon et al. 2004). Moreover, the observed substructures within halos are lower than predicted. These discrepancies have been used to prove that the power of density perturbations on small scales is less than the LCDM model prediction. This result has already motivated attempts to modify the LCDM model on small scales, such as the warm dark matter model, annihilating cold dark matter model (Kaplinghat et al 2000), and self-interacting dark matter model (Spergel & Steinhardt 2000).
On the other hand, the observed gravitational lensing of galaxy clusters, which yield constraints on small scale behavior of structure clustering, are found in good agreement with the LCDM model prediction (Metcalf 2004, Natarajan & Springel 2004). No reduction of small scale power is needed. Also, the $`N`$-body simulations show that no more than 70% of halos can be fitted by the standard universal spherical mass profile. There is considerable amount of variation in the density profile even among the halos that can be fitted by the universal mass profile (Jing 2000; Bullock et al. 2001; Ricotti 2003). The variation of the concentration parameter can be as large as a factor of two. Therefore, one needs more tests on the possible discrepancy between observations and prediction on small scales, especially using high quality samples.
In this paper, we study the small scale behavior of the cosmic field using the $`\mathrm{Ly}\alpha `$ transmitted flux of the QSO HE2347-4342, and model samples given by the hydrodynamic simulation. Along with the power spectrum, we also focus on the intermittent behavior of the field of QSO’s $`\mathrm{Ly}\alpha `$ transmitted flux on small scales. Roughly speaking, the intermittency of random field is characterized by strong enhancements (cuspy structures) scattered in a space with a low density background. If the cosmic mass field is given by the superposition of universal halos on various scales, all the non-Gaussian features should be from the halos. Thus, the intermittent features of QSOs’ Ly$`\alpha `$ transmitted flux on small scales would be effective to detect the cuspy behavior of cosmic mass field. This approach is not new. The intermittency of $`\mathrm{Ly}\alpha `$ transmitted flux has been studied using the absorption spectra of QSOs (Jamkhedkar et al. 2000; Jamkhedkar et al. 2003). The results have also been used to compare with simulations of standard LCDM model and warm dark matter model (Pando et al 2002; Feng et al. 2003). However, there are two reasons we want to revisit this topic. First, the data of HE2347-4342 $`\mathrm{Ly}\alpha `$ transmitted flux have higher resolution and S/N ratio than the Keck data used in previous works (Pando et al 2002; Jamkhedkar et al. 2003). Second, the newly developed hybrid cosmological hydrodynamic based on Weighted Essentially non-Oscillatory scheme (WENO) is especially effective in capturing singular and complex structures with a higher order spatial accuracy (Feng et al 2004). $`\mathrm{Ly}\alpha `$ transmitted flux has been extensively studied to calculate the power spectrum of mass perturbation on small scales (Croft et al. 2002; Viel et al 2004; McDonald et al. 2004).
The paper is organized as follows. In §2, we address the statistics of intermittency. Section 3 describes the data of HE2347-4342 and §4 presents the samples given by the WIGEON simulation. Sections 5 and §6 show the analysis and comparison of the power spectra and intermittent properties of observed data and simulation samples, respectively. The discussion and conclusion is presented in §7.
## 2 INTERMITTENT MASS FIELD
### 2.1 Cuspy halos and intermittency
Cuspiness of the mass density distribution can effectively be described up by density difference $`\mathrm{\Delta }\rho _r(𝐱)|\rho (𝐱+𝐫)\rho (𝐱)|`$, where $`r=|𝐫|`$. For a field given by superposition of cuspy halos, the field is regular at most locations x, i.e. $`\mathrm{\Delta }\rho _r(𝐱)0`$, when $`r0`$. On the other hand, cusps yield singular behavior,i.e., $`\mathrm{\Delta }\rho _r(𝐱)\mathrm{}`$, when $`r0`$. That is, the mass field consists of high spikes randomly and widely scattered in space, with a low field value between the spikes. Such spiky field is generally intermittent.
For a statistically homogeneous and isotropic random field, the probability distribution function (PDF) of density difference $`\mathrm{\Delta }\rho _r(𝐱)`$ for a given scale $`r`$ has to be independent of $`𝐱`$. The existence of singular density profiles means that events with large density difference $`\mathrm{\Delta }\rho _r`$ on small $`r`$ are more frequent than compared with a Gaussian field. Therefore, the PDF of $`\mathrm{\Delta }\rho _r`$ has to be long tailed, given by the events with extremely large density difference $`\mathrm{\Delta }\rho _r`$. The long tail would be more prominent for the PDF of smaller scales $`r`$. Thus, the long tail of $`\mathrm{\Delta }\rho _r`$ PDF of the mass field is an alternative tool to probe the cuspiness of the field.
As higher order moments are sensitive to the tail of the PDF, an effective measurement of the PDF long tail is given by the so-called structure function defined as
$$S_r^{2n}|\mathrm{\Delta }\rho _r(𝐱)|^{2n}.$$
(1)
where $`\mathrm{}`$ is the average over the ensemble of fields. As above mentioned, if the field is statistically homogeneous, $`S_r^{2n}`$ is independent of $`𝐱`$ and depends only on $`r`$. When $`n=1`$, we have $`S_r^2=|\mathrm{\Delta }\rho _r(𝐱)|^2`$, which is the mean of the square of the density fluctuations at $`r`$, and therefore, $`S_r^2`$ actually is the power spectrum of mass density fluctuations of the field.
Intermittency of a random field is defined by the divergence of the following ratio
$$\frac{S_r^{2n}}{[S_r^2]^n}\left(\frac{r}{L}\right)^\zeta ,$$
(2)
where $`L`$ is the size of the sample, and $`\zeta `$ is called intermittent exponent. Generally, $`\zeta `$ is $`n`$\- and $`r`$-dependent. The ratio in eq. (2) is the $`2n^{\mathrm{th}}`$ moment, $`S_r^{2n}`$, normalized by the power $`S_r^2`$. As $`S_r^2`$ measures the “width” (variance) of the PDF of $`\mathrm{\Delta }\rho _r(x)`$, and $`S_r^{2n}`$ is sensitive to the tail of the PDF, the ratio in eq. (2) measures the fraction of events in the long tail on the scale $`r`$. If the exponent $`\zeta `$ is zero or negative, the field is regular, i.e. smooth on smaller scales. If $`\zeta `$ is positive, the ratio diverges as $`r0`$, and the field is rough on small scales. In this case, the field is called to be intermittent (Gärtner & Molchanov, 1990; Zel’dovich, Ruzmaikin, & Sokoloff, 1990). Since $`\mathrm{\Delta }\rho _r(𝐱)0`$ for regular field, the $`r0`$ asymptotic behavior of $`S_r^{2n}/[S_r^2]^n`$ is dominated by cuspy structures. Intermittent exponent $`\zeta `$ measures cuspiness of the field. If the cuspy behavior is given by $`\rho (r)r^\alpha `$ with a constant $`\alpha `$, the exponent $`\zeta `$ should also be $`r`$-independent.
For a Gaussian field $`\rho (𝐱)`$, the PDF of the density difference $`\mathrm{\Delta }\rho _r`$ is also Gaussian. We have then
$$\frac{S_r^{2n}}{[S_r^2]^n}=(2n1)!!.$$
(3)
This ratio is independent of scale $`r`$, and therefore, the intermittent exponent $`\zeta =0`$.
### 2.2 Intermittent statistics with DWT variables
The basic statistical variable in eqs. (2), (3) is the density difference $`\rho (𝐱+𝐫)\rho (𝐱)`$, which contains the information of position $`𝐱`$ and spatial scale $`r`$. Therefore, it is convenient to use the statistical variables given by the discrete wavelet transform (DWT) decomposition of density field. For an 1-D field, the quantity $`\rho (x+r)\rho (x)`$ in terms of DWT variables is given by the wavelet function coefficient (WFC) defined as
$$\stackrel{~}{ϵ}_{j,l}=\psi _{j,l}(x)\rho (x)𝑑x,$$
(4)
where $`\psi _{j,l}(x)`$ is the discrete wavelet basis function, $`j`$ denotes the scale $`L/2^j`$, and $`l`$ the spatial range $`lL/2^j`$ to $`(l+1)L/2^j`$ (Daubechies, 1992; Fang & Thews 1998). The WFC, $`\stackrel{~}{ϵ}_{j,l}`$, is the density fluctuation (or difference) on scale $`L/2^j`$ at position $`l`$.
The structure functions in eq. (1) can then be re-written via the WFCs as (Farge at al. 1996)
$$S_j^n=|\stackrel{~}{ϵ}_{j,l}|^n.$$
(5)
The fair sample hypothesis allows one to calculate $`S_j^n`$ by using spatial averages, i.e., the average over $`l`$ so that
$$S_j^n=\frac{1}{2^j}\underset{l=0}{\overset{2^j1}{}}|\stackrel{~}{ϵ}_{j,l}|^n.$$
(6)
For $`n=2`$, we have
$$S_j^2=\frac{1}{2^j}\underset{l=0}{\overset{2^j1}{}}|\stackrel{~}{ϵ}_{j,l}|^2.$$
(7)
which is actually the power spectrum in the DWT modes $`P_jS_j^2`$ (Pando & Fang 1998; Fang & Feng 2000). For a Gaussian field, the Fourier power spectrum $`P(n)`$ is related to its DWT power spectrum $`P_j`$ by
$$P(n)=\frac{1}{L}\underset{j=0}{\overset{\mathrm{}}{}}P_j\left|\widehat{\psi }\left(\frac{n}{2^j}\right)\right|^2,$$
(8)
or
$$P_j=\frac{1}{2^j}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}|\widehat{\psi }(n/2^j)|^2P(n),$$
(9)
where $`\widehat{\psi }`$ is the Fourier transform of the wavelet function. This implies that the DWT power spectrum $`P_j`$ is the banded Fourier power of the flux fluctuations, and the band $`j`$ corresponds to the wavenumber around $`k=2\pi n/L2\pi 2^j/L`$.
Since $`r=L/2^j`$, the intermittent exponent $`\zeta `$ defined by eq. (2) can be calculated by
$$\frac{S_j^{2n}}{[S_j^2]^n}2^{j\zeta }.$$
(10)
Generally, $`\zeta `$ depends on $`n`$ and $`j`$. The statistics of the power spectrum and intermittency are based entirely on the DWT variables. It would be easy to make a uniform comparison between model predictions and observations with statistics on second and higher orders.
## 3 DATA OF HE2347-4342
The data used in our analysis is the transmitted flux of the $`\mathrm{Ly}\alpha `$ absorption spectrum of QSO HE2347-4342 ($`z=2.885`$, $`V=16.1`$). The optical echelle spectra were obtained at the ESO VLT UVES on 2001 November 23-24. The details on HE2347-4342 optical spectra have been described in Zheng et al 2004. The VLT data cover the wavelength range between 3600 and 4800 Å, which corresponds to the entire $`\mathrm{Ly}\alpha `$ wavelength range studied with FUSE from $`z=2.0`$ to $`2.9`$. Using IRAF tasks designed for echelle data, a normalized spectrum was obtained. The spectrum has $`24000`$ points with resolution $`\delta \lambda 0.05`$ Å. The data are given in the form of pixels with wavelength $`\lambda `$, flux $`F`$ and noise $`\sigma `$. In terms of the local velocity the resolution is $`dv3.5`$ km s<sup>-1</sup>. The S/N ratio of the spectrum is about $`110`$ per 0.1 Å bin at 4700 Å, and about $`46`$ at 3850 Å. This is respectively $``$ $`2.5`$ and $`10`$ times of the data of Keck echelle spectrum used in Feng et al 2003.
For our purpose, the useful wavelength region is from $`\mathrm{Ly}\beta `$ absorption to the $`\mathrm{Ly}\alpha `$ emission, excluding a region close to the quasar to avoid proximity effects. Below 3984 Å $`\mathrm{Ly}\beta `$ absorption starts to appear. Therefore, we take the range from $`3986.014395.600`$ Å, corresponding to redshift from $`2.278`$ to $`2.615`$. In this wavelength range, the mean transmission $`e^\tau `$ is $`0.796`$. This redshift range contains about 2<sup>13</sup> pixels. The size of a cell on the DWT scale $`j`$ corresponds to $`N=2^{13j}`$ pixels. The distance between $`N`$ pixels in the units of the local velocity scale is given by $`\delta v=2c[1\mathrm{exp}(Ndv/2c)]`$ km s<sup>-1</sup>, corresponding to comoving scale $`D=(\delta v/H_0)[\mathrm{\Omega }_m(1+z_m)^3+\mathrm{\Omega }_\mathrm{\Lambda }]^{1/2}`$.
Metal lines are a major cause of contamination in the spectra. But the doppler width of metal lines are generally narrow with $`\delta v20\mathrm{km}\mathrm{s}^1`$. In this paper, we restrict our analysis only to scales $`\delta v30\mathrm{km}\mathrm{s}^1`$ where contamination due to metal lines is low (Hu et al 1995; Boskenberg et al 2003; Kim, et al. 2004).
## 4 HYDRODYNAMIC SIMULATION
The simulation uses the newly developed hybrid cosmological hydrodynamic codes based on the Weighted Essentially Non-Oscillatory (WENO) scheme (Harten et al. 1987; Liu et al. 1994; Jiang & Shu 1996; Shu 1998; Fedkiw et al. 2003; Shu 2003). We will name this code as WIGEON, Weno for Intergalactic medium and Galaxy Evolution and formatiON. For details of the numerical method and tests, we refer to Feng et al.(2004). The simulation sample we analyze here is actually the same as those used in our previous papers on the statistical study of temperature, entropy, baryonic fraction and velocity fields of intergalactic medium (He et al, 2004; He et al. 2005; Kim et al. 2005). It was performed in a cubic box of side length 12 h<sup>-1</sup> Mpc with a 192<sup>3</sup> grid and an equal number of dark matter particles. The cosmogony is the standard LCDM model specified by the density parameter $`\mathrm{\Omega }_m=0.3`$, the baryon density $`\mathrm{\Omega }_b=0.047`$, the cosmological constant $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, the Hubble constant $`h=0.7`$, the shape factor $`\mathrm{\Gamma }=\mathrm{\Omega }_mh\mathrm{exp}[\mathrm{\Omega }_b(1+\sqrt{2h}/\mathrm{\Omega }_m)]=0.166`$, and the rms mass fluctuations in spheres of $`8h^1`$Mpc, $`\sigma _8=0.9`$. The ratio of specific heats is $`\gamma =5/3`$. Since the shock heating of cosmic gas is significant (He et al 2004), the resolution of the simulation should be less than the thickness of the shock, which is of the order of the dissipation length, i.e., the Jeans diffusion $`0.10.3`$ h<sup>-1</sup> Mpc for redshifts $`z<4`$ (Bi et al. 2003). The size of the grid is $`12/192=33/2^8=0.063`$ h<sup>-1</sup> Mpc. Therefore, the resolution of our simulation is sufficient to capture shocks.
Atomic processes including ionization, radiative cooling and heating are modelled as in Cen (1996) in a primeval plasma of hydrogen and helium of composition ($`X=0.76`$, $`Y=0.24`$). The uniform UV-background of ionizing photons is assumed to have a power-law spectrum of the form $`J(\nu )=J_{21}\times 10^{21}(\nu /\nu _{HI})^\alpha `$ergs<sup>-1</sup>cm<sup>-2</sup>sr<sup>-1</sup>Hz<sup>-1</sup>, with $`\alpha =1`$, where the photo ionizing flux is normalized by parameter $`J_{21}`$ at the Lyman limit frequency $`\nu _{HI}`$, and is suddenly switched on at $`z6`$ to heat the gas and re-ionize the universe.
One-dimensional fields are extracted along randomly selecting lines of sight in the simulation box. The density, temperature and velocity of the neutral gas fraction on grids are Gaussian smoothed using FFT techniques which form the fundamental data set. The one-dimensional grid containing the physical quantities is further interpolated by a cubic spline. Using this one-dimensional grid, the optical depth $`\tau `$ is then obtained by integrating in real space and we include the effect of the peculiar velocity and convolve with Voigt thermal broadening. To have a fair comparison with observed spectra, $`\tau `$ was Gaussian smoothed to match with the spectral resolutions of observation. The transmitted flux $`F=\mathrm{exp}(\tau )`$ is normalized such that the mean flux decrement in the spectra match with observations.
Each mock spectrum is sampled on a $`2^{10}`$ grid with the same spectral resolution as the observation. As the corresponding comoving scale for $`2^{10}`$ pixels is larger than the simulation box size, we replicate the sample periodically. To achieve the greatest statistical independence, we randomly change the direction of line of sight while crossing the boundary of the simulation box. By the way, 1000 mock spectra are generated.
## 5 POWER SPECTRUM OF THE TRANSMITTED FLUX OF HE2347-4342
### 5.1 Treatment of unwanted modes
In order to calculate the DWT power spectrum \[eq. (7)\] of the transmitted flux of HE2347-4342, we should properly treat unwanted data, including the pixels without data, contamination of metal lines etc. Although the $`S/N`$ is high on an average, it is as low as about 1 for some pixels, such as pixels with negative flux. We must reduce the uncertainty given by low S/N pixels. In the DWT analysis, the conventional technique of reducing these uncertainties is given by the algorithm of DWT denoising by thresholding (Donoho 1995) or conditional-counting, (Jamkhedkar et al. 2003) as follows
1. Calculate the SFCs of both transmission $`F(x)`$ and noise $`\sigma (x)`$, i.e.
$$ϵ_{jl}^F=F(x)\varphi _{jl}(x)𝑑x,ϵ_{jl}^N=\sigma (x)\varphi _{jl}(x)𝑑x.$$
(11)
2. Identify an unwanted mode $`(j,l)`$ using the threshold condition
$$\left|\frac{ϵ_{jl}^F}{ϵ_{jl}^N}\right|<f$$
(12)
where $`f`$ is a constant. This condition flags all modes with S/N less than $`f`$. We can also flag modes dominated by metal lines.
3. Since all the statistical quantities in the DWT representation are based on an average over the modes $`(j,l)`$, we will skip all the flagged modes while computing these averages, i.e. the average is over the un-flagged modes $`N(f)`$ only.
With this method, no rejoining and smoothing of the data are needed. The condition in eq. (12) is applied on each scale $`j`$, and therefore the unwanted modes are flagged on a scale-by-scale basis. Generally, for scale $`j`$, $`N(f)2^j`$. If the size of an unwanted data segment is $`R`$, condition in eq. (12) only flags modes $`(j,l)`$ on scales less than or comparable to $`R`$. We also flag two modes around each unwanted mode to reduce any boundary effects of the chunks. With the conditional-counting method, we can still calculate the power spectrum by the estimators of eq. (7), but the average is not over all modes $`l`$, but over the un-flagged modes only.
### 5.2 Power spectrum of the DWT modes
We calculate the power spectrum of the transmitted flux fluctuations, $`\mathrm{\Delta }F=F(\lambda )F(\lambda )`$, of HE2347-4342. To consider the correction of the noise on eq. (7), the power spectrum of transmitted flux is given by (Pando & Fang 1998b; Fang & Feng 2000; Jamkhedkar, Bi, & Fang, 2001)
$$P_j=\frac{1}{2^j}\underset{l=0}{\overset{N(f)}{}}(\stackrel{~}{ϵ}_{jl}^F)^2\frac{1}{2^j}\underset{l=0}{\overset{N(f)}{}}(\stackrel{~}{ϵ}_{jl}^n)^2.$$
(13)
The first term on the r.h.s. of eq. (13) is the same as eq. (7), in which the wavelet coefficients (WFC) $`\stackrel{~}{ϵ}_{jl}^F`$ are given by
$$\stackrel{~}{ϵ}_{jl}^F=F(x)\psi _{jl}(x)𝑑x.$$
(14)
The second term on the r.h.s. of eq.(13) is due to the noise field $`\sigma (\lambda )`$ and is calculated by
$$(\stackrel{~}{ϵ}_{jl}^n)^2=\sigma ^2(x)\psi _{jl}^2(x)𝑑x.$$
(15)
Figure 1 plots the results of the DWT power spectrum, in which the parameter $`f`$ is taken to be 1, 2, 3 and 5. At the first glance, the conditional-counting of eq. (12) would seem to preferentially drop modes in the low transmission regions, and power spectrum eq. (13) should be $`f`$-dependent. However, Figure 1 shows that the power spectrum $`P_j`$ is independent of $`f`$ on entire the scale range considered for $`f=1`$ to 5. This can be seen from eq. (13), which shows that the contribution to the power $`P_j`$ given by mode $`(j,l)`$ is $`(\stackrel{~}{ϵ}_{jl}^F)^2(\stackrel{~}{ϵ}_{jl}^n)^2`$. The noise substraction term $`(\stackrel{~}{ϵ}_{jl}^n)^2`$ guarantees that the contribution of modes with small ratio $`S/N`$ to $`P_j`$ is always small or negligible. For instance, the modes with negative flux, i.e., the modes with flux having the same order of magnitude as noise, the two terms $`(\stackrel{~}{ϵ}_{jl}^F)^2`$ and $`(\stackrel{~}{ϵ}_{jl}^n)^2`$ statistically cancel each other. Thus, all the dropped modes have a very small or negligible contribution to $`P_j`$ regardless the parameter $`f`$. Denoising by thresholding or conditional-counting is reliable.
Figure 2 compares the DWT power spectra measured in the mock samples and the observed data. We take the same parameter $`f=3`$ for both real data and mock samples. The error bars of $`P_j`$ are the maximum and minimum range of $`P_j`$ from bootstrap re-sampling. Since the PDF of $`\stackrel{~}{ϵ}_{jl}^F`$ is highly non-Gaussian (§5.1), a reasonable estimation of the errors for the average over the ensemble of $`\stackrel{~}{ϵ}_{jl}^F`$ is given by bootstrap re-sampling (Jamkhedkar et al 2003). That is, for observed sample, the bootstrap re-sampling is done on the set of $`N(f)`$ data for each scale $`j`$, and for the simulation samples the bootstrap re-sampling is on the set of $`N_{sim}(f)`$ data from the 1000 simulation samples. The error bar is calculated thus. For a data set with $`N`$ points, $`N`$ realizations or data sets are created by drawing points from the original set with replacement. The average is calculated over the $`N`$ realizations. The power spectrum is calculated on the scales from $`\delta v=224`$ to 28 km s<sup>-1</sup>, corresponding to comoving scale $`0.630.079`$ h<sup>-1</sup> Mpc in the LCDM model. The error bars for the real data are given by the maximum and minimum of bootstrap re-sampling. Figure 2 shows that simulation samples basically are in agreement with observations on the scales considered. Particularly, no discrepancy has been found even on the smallest scale $`\delta v=28\mathrm{km}\mathrm{s}^1`$ or length scale $`D=0.079`$ h<sup>-1</sup> Mpc.
## 6 INTERMITTENT BEHAVIOR
### 6.1 PDF of $`\stackrel{~}{ϵ}_{jl}^F`$
As first step to describe the intermittent behavior of the $`\mathrm{Ly}\alpha `$ transmitted flux, we show in Figure 3 the PDFs of the normalized WFCs of flux field $`\stackrel{~}{ϵ}_{jl}^F/(\stackrel{~}{ϵ}_{jl}^F)^2^{1/2}`$. These distributions are compared with a Gaussian distribution with zero mean and unit standard deviation. On the scales $`224`$ km s<sup>-1</sup> the PDFs of the normalized WFCs do not strongly deviate from Gaussian fields. But on scales $`<`$ 224 km s<sup>-1</sup>, the central peak and long-tail of the PDFs show that the field is highly non-Gaussian. On the scale 56 km s<sup>-1</sup>, the long tail extends to $`\stackrel{~}{ϵ}_{jl}^F/(\stackrel{~}{ϵ}_{jl}^F)^2^{1/2}5`$. That is, the power of some long tail events $`(\stackrel{~}{ϵ}_{jl}^F)^2`$ can be larger than the mean power $`(\stackrel{~}{ϵ}_{jl}^F)^2`$ by a factor of 20-30. This is in agreement with the result based on Keck data (Jamkhedkar et al 2003).
We see that for $`f=1,3`$ and $`5`$, the PDFs given by $`f=1,3`$ and 5 are essentially the same, especially, the central peaks of the PDFs are insensitive to the parameter $`f`$. This indicates that the statistical result does not depend on data at pixels with low number of $`S/N`$. This point is important. For instance, a saturated absorption region on scale $`j_s`$ may yield $`\stackrel{~}{ϵ}_{jl}^F0`$ on smaller scales or $`j>j_s`$, i.e. pixels with a low value of $`S/N`$. However, we cannot draw information of clustering of cosmic matter from that region. We also can not say whether this region underwent a strong nonlinear evolution. Therefore, the $`f`$-independence of the PDFs of Figure 3 provides an valuable measurement of the non-Gaussian behavior, irrespective of whether the region is saturated or not.
### 6.2 Structure functions
We now calculate the structure functions eqs. (2) or (7) for the transmitted flux fluctuations of HE2347-4342. For real data, the results are illustrated in Figure 4. In calculating the high order moment $`S_j^{2n}`$, we did not subtract the noise term in eq. (14), because as noise is considered Gaussian, its higher order moments are small (see §6.4 below). The error bars are found by bootstrap re-sampling. We see that $`\mathrm{log}_2[S_j^{2n}/(S_j^2)^n]`$ with $`n=2`$ has errors even smaller than the power spectrum of Figure 1. This is because the uncertainty in the power spectrum is caused by rare and improbable long tail events in the fluctuations, i.e. the tail of the PDF shown in Figure 3. The large uncertainty of the PDF tail leads to the large uncertainty of the power spectrum. On the other hand, the structure function is the ratio between $`S_j^{2n}`$ and $`(S_j^2)^n`$ and it reduces the effect of individual high spikes (tail events). Therefore, the structure functions are an effective and stable tool for high order statistics. Similar to the power spectra of Figures 1 and 2, and the PDFs of Figure 3, the structure functions of Figure 4 are basically independent of $`f`$. Therefore, the scale- and $`n`$-dependencies of the structure function exist regardless of the saturated absorption, and give a measurement of the non-Gaussian clustering of the baryonic gas.
The value of $`\mathrm{log}_2[S_j^{2n}/(S_j^2)^2]`$ generally is bigger for smaller length scales. For a given $`n`$, the $`j`$-dependence of $`\mathrm{log}_2[S_j^{2n}/(S_j^2)^n]`$ can approximately be fitted by eqs. (2) and (10) with a positive exponent $`\zeta `$. Therefore, the field of the transmitted flux fluctuations is highly intermittent. This shows again that the cuspy feature of cosmic clustering can be seen in high mass density areas, like massive halos, as well as in low mass density areas, like the clouds of Ly$`\alpha `$ absorption clouds (He et al. 2004; Pando et al. 2004; He et al. 2005; Kim et al. 2005). Actually the success of the semi-analytical lognormal model (Bi & Davidsen, 1997) in explaining the $`\mathrm{Ly}\alpha `$ forest has already indicated that the mass field of the cosmic baryon gas is probably intermittent, because a lognormal field is intermittent.
The structure functions measured for the mock samples are shown in Figure 5, in which the error bars are also given by bootstrap re-sampling of the 1000 samples. It shows once again that the statistical uncertainty of the structure function is reliable for a high order statistical test. Figure 5 also shows that the structure function is $`f`$-independent. Figure 6 gives a comparison between the mock samples and the real data. We see that for all scales from 224 to 28 km s<sup>-1</sup>, and all order $`n`$, the structure functions of the mock samples are consistent with real data within their error bars. The consistence is very good on the smallest scale 28 km s<sup>-1</sup>.
### 6.3 Intermittent exponent
Following the definition of intermittent exponent $`\zeta `$ eq. (2) or eq. (10), we can calculate $`\zeta `$ in the scales range $`j_1`$ to $`j_2`$ by
$$\zeta _n=\frac{1}{j_1j_2}\mathrm{log}_2\left[\frac{S_{j_1}^{2n}(S_{j_2}^2)^n}{S_{j_2}^{2n}(S_{j_1}^2)^n}\right].$$
(16)
The result is listed in Table 1. The error of $`\zeta _n`$ is estimated by $`\sigma _{\zeta _n}=(1/|j_1j_2|)\sqrt{\sigma _1^2+\sigma _2^2}`$, where $`\sigma _1`$ and $`\sigma _2`$ are the errors of $`\mathrm{log}_2[S_{j_1}^{2n}/(S_{j_1}^2)^n]`$, and $`\mathrm{log}_2[S_{j_2}^{2n}/(S_{j_2}^2)^n]`$, respectively. The scale $`j`$ represents a local velocity $`\delta v=2^{13j}\times 3.5`$ km s<sup>-1</sup>. As expected, Table 1 shows that the values of intermittent exponents of real data and mock samples are consistent with each other on all orders and scales considered.
In Table 2, we list the intermittent exponent of mock samples for different scales ranges. It clearly shows that the intermittent exponent is scale-dependent. We also tried to fit the $`j`$-dependence of $`\mathrm{log}_2[S_j^{2n}/(S_j^2)^2]`$ from $`j=7`$ to 10 by
$$\mathrm{log}_2[S_j^{2n}/(S_j^2)^n]=A+j\zeta _n$$
(17)
with assumption of $`\zeta _n`$ to be constant ($`j`$-independent). We found that the goodness-of-fit, $`Q`$, (Press et al 1992) of the data to eq.(18) with a constant $`\zeta _n`$ is always $`0.1`$. That means that the assumption that $`\zeta _n`$ is independent of scale does not hold. Therefore, $`\zeta _n`$ most likely is scale-dependent. This result is inconsistent with the cuspy center profiles $`\rho (r)r^\alpha `$ with a constant index $`\alpha `$, which predicts a constant $`\zeta _n`$ on small scales. Moreover, the $`j`$-dependence of $`\zeta _n`$ shown in Table 2 is not monotonic. This makes it more difficult to fit with the standard universal profiles. Therefore, the LCDM model seems to predict a different scale-dependence of the intermittent exponent for the Ly$`\alpha `$ absorption clouds in contrast to the cuspy behavior of the universal profile of dark matter halos.
### 6.4 $`n`$-dependence of structure function
We now turn to the $`n`$-dependence of the structure functions. Figures 7 and 8 are, respectively, $`\mathrm{log}_2[S_j^{2n}/(S_j^2)^n]`$ vs. $`n`$ for the real and mock samples of HE2347-4342 on scales $`\delta v=22428\mathrm{km}\mathrm{s}^1`$. For a Gaussian field, the $`n`$-dependence of $`\mathrm{log}_2(S_j^{2n}/(S_j^2)^n)`$ is given by eq. (3), i.e. $`\mathrm{log}_2(2n1)!!`$, which is also plotted in Figures 7 and 8. The curves of $`\mathrm{log}_2[S_j^{2n}/(S_j^2)^n]`$ vs. $`n`$ for both real and mock samples are much higher than Gaussian field on scales $`56\mathrm{and}28\mathrm{km}\mathrm{s}^1`$, but not much different from Gaussian field on scale of 224 km s<sup>-1</sup>. Therefore, it is reasonable to ignore the noise term in calculating the high order moment (§6.2).
More interesting is to fit the observed $`n`$-dependence for $`\mathrm{log}_2[S_j^{2n}/(S_j^2)^n]`$ with
$$\mathrm{log}_2\frac{S_j^{2n}}{(S_j^2)^n}n^\alpha (n1).$$
(18)
The motivation is to compare the fields with a lognormal field, for which $`\alpha =1`$. Figure 7 shows that the best fit of $`\alpha `$ is in the range $`\alpha =0.3`$ \- $`0.4`$ on scale $`56\mathrm{km}\mathrm{s}^1`$, and $`0.30.5`$ on $`28\mathrm{km}\mathrm{s}^1`$. That is, the value of $`\alpha `$ seems to approach to 1 when scale is small. The transmitted flux field is closer to a lognormal field on small scales. This somewhat supports the lognormal models of $`\mathrm{Ly}\alpha `$ forests. Figure 8 shows that values of $`\alpha `$ given by mock samples always lie in the range from $`0.30.4`$ and are consistent with real data.
## 7 DISCUSSIONS AND CONCLUSIONS
We have showed that Ly$`\alpha `$ transmitted flux field of HE2347-4342 is significantly intermittent, especially on small scales. We found that the power spectrum is in good agreement with the data of $`\mathrm{Ly}\alpha `$ transmitted flux of HE2347-4342. There is no evidence of any discrepancy between the LCDM model from observed intermittent features on scale as small as about $`\delta v=28`$ km s<sup>-1</sup>, and for statistical orders from 2 to 8. Accordingly, there is no need of reducing the power relative to the standard LCDM model up to length scale $`0.079`$ h<sup>-1</sup> Mpc.
Comparing the current results with our previous studies on the same topic, we found that the intermittency sensitively relies on the quality of both the observed data and simulation sample. In the first stage of our study, we studied the intermittency of Ly$`\alpha `$ transmitted flux with Keck data and model samples produced by pseudo-hydro simulations (Pando et al 2002). Although the simulation samples can fit the observed power spectrum, and are also intermittent, the intermittent exponent does not fit the real data. There is a discrepancy between the observed data and simulation sample on small scales. Physically, that is probably because the pseudo-hydro simulations assumed that (1) the baryon distribution is proportional to that of dark matter point-by-point, and (2) the gas temperature is related to the density by a power law equation of state. However, it has been shown that the relation between temperature and IGM density is multi-phased. The relation between temperature and density can approximately be described by a power-law equation. However, for a given density, the temperature actually is not single-valued, but varies from $`10^410^7`$ K (He et al 2004).
In the second stage, the model samples are produced by full hydro simulations with two assumptions mentioned above (Feng et al 2003). The result was a great improvement with respect to the first round. It shows that the intermittent behavior of the Keck data and simulation is basically consistent with each other, but a discrepancy can still be seen on scale $`\delta v=32`$ km s<sup>-1</sup>. In the current study, the observed data of HE2347-4342 probably is among the best quality for our purpose. Its intermittency is in good agreement with the LCDM model on small scales less than $`\delta v=30`$ km s<sup>-1</sup>.
The star formation and their feedback on the cosmic gas evolution are not considered in our simulation. Generally speaking, there are two types of the feedbacks: (1) photoionization heating by the UV emission of stars and AGNs, and (2) injection of hot gas and energy by stars. The photoionization heating can be properly considered, if the UV background is adjusted by fitting the simulation with the observed mean flux decrement of QSO’s $`\mathrm{Ly}\alpha `$ absorption spectrum. The effect of injecting hot gas and energy is localized in massive halos, and therefore, its effect is weak when we consider to avoid proximity effects. Therefore, the major conclusions would not be significantly affected even while considering the effect of star formation.
Intermittency is very effective to probe the details of the singular features of a random field. Our simulation samples show that the intermittent exponent of the Ly$`\alpha `$ transmitted flux field probably is scale-dependent. This result is different from the prediction of universal mass profile with cuspy center $`\rho (r)r^\alpha `$. If the index $`\alpha `$ is constant, the intermittent exponent should be scale-independent. Therefore, the scale-dependence of the intermittent exponent indicates that the distribution of baryon gas is decoupled from the underlying dark matter (e.g. He et al 2004, Kim et al 2005). The data of HE2347-4342 only is unable to test the prediction of scale-dependence of the intermittent exponent. More high quality QSO absorption spectra would be very valuable to test the $`j`$ dependence of $`\zeta _n`$ on small scales.
LLF acknowledges support from the National Science Foundation of China (NSFC). This work is supported in part by the US NSF under the grants AST-0507340. |
warning/0507/astro-ph0507658.html | ar5iv | text | # Study of DDO 68: nearest candidate for a young galaxy?
## 1 Introduction
Some of dwarf gas-rich galaxies with very low metallicity were considered as probable young galaxies since the seminal paper by Searle & Sargent (SS72 (1972)). The most metal-poor dwarf I Zw 18 (with the value of 12+$`\mathrm{log}`$(O/H)=7.17–7.21, e.g., Alloin et al. Alloin (1978), Izotov et al. Izotov99 (1999)) has for 30 years been the Candidate No.1 of the local young galaxies. Recently, Izotov & Thuan (IZw18CMD (2004)) have demonstrated from the very deep Hubble Space Telescope (HST) Colour-Magnitude Diagram (CMD) data that there are no stars in this object older than 0.5 Gyr; that is, this nearby object (D=15 Mpc) is indeed a genuine young local galaxy.
There are several other eXtremely Metal-Deficient objects (XMD, conditionally with the characteristic value Z$`<`$1/10 Z<sup>1</sup><sup>1</sup>1The solar oxygen abundance is accepted as 12+$`\mathrm{log}`$(O/H) = 8.66 according to Asplund et al. (Solar04 (2004))., or 12+$`\mathrm{log}`$(O/H)$``$7.65; e.g., Kunth & Östlin Kunth2000 (2000)), considered as good candidates to young local galaxies. First of all, these are the components of the galaxy pair SBS 0335–052 E and W, with 12+$`\mathrm{log}`$(O/H)=7.29 and 7.12, respectively (e.g., Izotov et al. 1997a , Lipovetsky et al. Lipovetsky99 (1999), Pustilnik et al. VLA (2001); 2004a , Izotov & Thuan IT05W (2005)). This unique pair at $`D`$54 Mpc is, however, 3.6 times more distant than I Zw 18.
In our discussion, we distinguish truly young galaxies (similar to I Zw 18), in which their young stellar generation is made up of pregalactic matter, from the so-called tidal young dwarfs (e.g., Duc & Mirabel Tidal (1999)), similar to Holmberg IX in the vicinity of M81. Holmberg IX is devoid of old Red Giant Branch (RGB) stars (Makarova et al. HoIX (2002)), so its star formation began recently. But its gas metal content is close to that of the massive parent galaxy ISM, from which it was recently formed, so in this context such tidal dwarfs should be treated not as young but as rejuvenated.
It is worth noting that the very low ISM metallicity is considered only as an indication of the possible unevolved status of a galaxy. In particular, from the theoretical point of view, several evolutionary scenarios can lead to such a result. Briefly, they include: a) the significant metal loss due to galactic superwinds in the shallow gravitational wells of dwarf galaxies; b) the inflow and the mixture of an unevolved intergalactic gas to an evolved galaxy; c) the very slow astration and related ISM enrichment characteristic of the very low surface brightness (LSB) galaxies, and d) the recent onset of the first star formation episode. All but one correspond to an aged galaxy. Moreover, observationally among 16 known XMD galaxies at distances less than 15 Mpc, at least nine objects, either from the CMD analysis (e.g., Sextans A and B, GR 8, Leo A) or from the red colours of the unresolved stellar population in their outer regions (e.g., UGCA 20, UGC 2684, van Zee et al. vZee\_UA20 (1996, 1997)), are recognized as old systems. Only one, I Zw 18, is now considered as a truly young object.
In searches for new XMD galaxies, we found that some of them populate the ‘void’ regions, where the distances from the void centers to luminous galaxies (L $`>`$ L, L correspond to M<sub>B</sub>=–19.6 for H=72 km s<sup>-1</sup> Mpc<sup>-1</sup>, assumed in the paper) exceed 4–5 Mpc for the smallest voids and 8–10 Mpc for the more typical voids. In particular, the XMD blue compact dwarf (BCD) galaxy HS 0822+3542 with 12+$`\mathrm{log}`$(O/H)=7.38, appeared near the centre of the nearby small Lynx-Cancer void (Kniazev et al. Kniazev00 (2000), Pustilnik et al. 2003b ). Voids delineated by luminous, massive galaxies are not absolutely empty. Some numbers of the lower-mass galaxies fill in voids (e.g., Salzer Salzer89 (1989), Pustilnik et al. PULTG (1995), Lindner et al. void\_gal (1996), Popescu et al. Popescu97 (1997), Grogin & Geller Grogin00 (2000) among others), probably forming the substructures of the filament types (Lindner et al. void\_gal (1996), Gottlöber et al. Gott03 (2003), Fairall et al. Fairall04 (2004)). The very low-density environment in voids and the scaled-down mass spectrum of the preformed DM halos, coupled with the effect of reionization, are all expected to cause formation of a somewhat different galaxy population from any in more typical, denser environments (e.g., Peebles Peebles01 (2001), Gottlöber et al. Gott03 (2003)). A significantly reduced probability of galaxy interactions in voids, which is suggested as an important factor of star formation history (e.g., Rojas et al. SDSS\_void (2004)), can cause the differences in the rate of chemical enrichment.
Some indication of the less evolved status of low-mass galaxies in voids was presented, e.g., by Huchtmeier et al. (void\_HI (1997)) and Pustilnik et al. (HI\_void (2002)), based on H i data. The higher star formation (SF) activity of void galaxies in respect to ‘wall’ galaxies, detected by Grogin & Geller (Grogin00 (2000)) and Rojas et al. (SDSS\_void (2004)), also indicates their larger available gas fuel reservoir. However, the metallicity issue of the void galaxy population has not yet been properly addressed, mainly due to observational difficulties. Only the study of absorption lines in Ly$`\alpha `$ clouds with low column densities in various environments led to the conclusion that such objects in voids show a significantly lower content of heavy elements (Lu et al. Lu98 (1998)).
Since the Lynx-Cancer void is one of the nearest (D$`{}_{\mathrm{center}}{}^{}`$11 Mpc), many of the underluminous galaxies falling in this region are sufficiently bright, so that in contrast to the situation in more distant voids, even inconspicuous H ii regions in more or less typical late-type dwarfs are available for determining chemical abundance in ‘routine’ programs. Therefore, it was tempting to search in this void for some other XMD objects, including likely young galaxy candidates. To that end we performed the spectrophotometry of ten dwarfs in this region, and the data will be presented separately elsewhere. In this paper we report the discovery of an extremely low oxygen abundance in dIm/BCD galaxy DDO 68 (UGC 5340=VV 542, J2000 coordinates 09:56:45.7 +28:49:35), situated at the periphery of this void at a distance of D$``$6.5 Mpc. We also present the results of $`V,R`$, and H$`\alpha `$ imaging, indicate its very blue colours, and discuss DDO 68 properties that make it a probable candidate for a young galaxy.
## 2 Observations and reduction
All observations were conducted with the SCORPIO multi-mode instrument (Afanasiev & Moiseev SCORPIO (2005)) installed in the prime focus of the SAO 6 m telescope (BTA), during 2 runs – in November 9–11, 2004 and January 7–13, 2005. For the long-slit spectral observations the grism VPH550g was used with the 2K$`\times `$2K CCD detector EEV 42-40, with the exposed region of 2048$`\times `$600 px. This gave the range 3500–7500 Å with $``$2.0 Å pixel<sup>-1</sup> and FWHM $``$12 Å along the dispersion. The scale along the slit was 0$`\stackrel{}{.}`$18 pixel<sup>-1</sup> and total extent of $``$2′. The spectra were obtained for several slit positions, crossing the bright knots, marked by their numbers on the $`V`$-band image of the galaxy obtained also in this program (Fig. 1). Exposure times varied between 15 to 60 min. When longer than 15 min, they were broken into subexposures of 15-min duration. The objects spectra were complemented before or after by the reference spectra of He–Ne–Ar lamp for the wavelength calibration. Bias and flat-field images were also acquired to perform the standard reduction of 2D spectra. Spectral standard star Feige 34 (Bohlin Bohlin96 (1996)) was observed during the night for the flux calibration.
The primary task of these observations was to get high S-to-N spectra of several H ii regions in order to detect the faint line \[O iii\]$`\lambda `$4363 and to directly measure their electron temperatures T<sub>e</sub> and oxygen abundances. Seeing was $``$0$`\stackrel{}{.}`$8 in the November 2004 run and $``$1$`\stackrel{}{.}`$4–1$`\stackrel{}{.}`$7 in the January 2005 run. We also used SCORPIO in the imaging mode (2048$`\times `$2048 px, binned 2$`\times `$2, the field of view of $``$6′$`\times `$6′) to study DDO 68 morphology, structure and colours. The standard Johnson-Cousins $`V,R_c`$ and the middle-band filters for H$`\alpha `$-line images SED665 (central wavelength = 6622 Å, FWHM=191 Å) and SED607 (central wavelength = 6063 Å, FWHM=167 Å) were used. On the night of January 12, 2005, we obtained 30-minute exposures of the galaxy in $`V`$ and $`R`$ bands (all broken into 10-min subexposures), and 15-min exposures in both SED665 and SED607 filters. The night was photometric, with the seeing between 1$`\stackrel{}{.}`$6 and 1$`\stackrel{}{.}`$8. Bias and flat-field images were acquired to perform the standard reduction. For the broad-band calibration, we observed the standard stars during the night in the field of QSO FBQ 0951+263 (Nakos et al. Nakos03 (2003)), while for the calibration of H$`\alpha `$ images we observed the spectrophotometric standard Feige 34. The coefficients of the resulting photometric system have the overall uncertainties of 0$`\stackrel{m}{.}`$046 in $`V`$ and 0$`\stackrel{m}{.}`$056 in $`R`$. The standard pipeline with the use of IRAF<sup>2</sup><sup>2</sup>2IRAF: the Image Reduction and Analysis Facility is distributed by the National Optical Astronomy Observatory, which is operated by the Association of Universities for Research in Astronomy (AURA) under cooperative agreement with the National Science Foundation (NSF). and MIDAS<sup>3</sup><sup>3</sup>3MIDAS is an acronym for the European Southern Observatory package – Munich Image Data Analysis System. was applied for the reduction of both long-slit spectra and images, which included the next steps.
Cosmic ray hits were removed from all spectra and direct images in MIDAS. Using IRAF packages from CCDRED for spectra and images, we subtracted bias, and performed flat-field correction. After that, spectra were wavelength-calibrated. Night sky background was subtracted from all spectra. Then, using the data on spectrophotometry standard stars, all spectra were transformed to absolute fluxes. One-dimensional spectra were extracted by summing up, without weighting, various numbers of rows along the slit depending on exact region of interest. The next step of the reduction of direct images consisted of shifting all V, R, SED607, and SED665 images into the common world coordinate system. Finally, all subexposures were combined into one image by summation.
To obtain the integrated photometry, the $`V`$ and $`R`$ surface brightness profiles (SBPs) and the fitting of SBPs, the method and the programs described in detail in Kniazev et al. (2004a ) were used: (1) $`V`$ and $`R`$ images were combined and filtered with a smooth-and-clip filter (Shergin, Kniazev & Lipovetsky Sh\_Kn\_Li\_96 (1996)); (2) the galaxy was detected above the 3$`\sigma `$ noise level on the combined and filtered image and the mask-frame was created to show the location of the galaxy for the subsequent steps of the analysis; (3) all background galaxies and stars were additionally masked; (4) using the mask-frame, the integrated photometry was calculated, the SBPs were created for each filter with the circular apertures with the step of 2″, and the errors were calculated; (5) fitting of SBPs by the exponential disc was performed with weights $`w_k=\sigma _k^1`$, where $`\sigma _k`$ is the error of the flux calculated within the circular aperture for each SB level in the previous module.
In the three H ii regions (Knots 1, 2, and 3 in the Northern ring/oval) where the line \[O iii\]$`\lambda `$4363 was seen, we extracted only the regions along the slit where this line was well above the noise (from 1$`\stackrel{}{.}`$8 to 3$`\stackrel{}{.}`$8 for various knots). All emission lines were measured by applying the MIDAS programs described in detail in a recent paper by Kniazev et al. (2004b ). Briefly, they draw continuum, perform robust noise estimation, fit separate lines by a single Gaussian superimposed on the continuum-subtracted spectrum and integrate the line flux. The emission lines, blended in pairs or triplets, were fitted simultaneously as blend of two or more Gaussian features. The quoted errors of singular line intensities include the following components. The first is related to the Poisson statistics of the line photon flux. The second component is the error resulting from the creation of the underlying continuum, which gives the main contribution to the errors of faint lines. For the fluxes of the lines in blends, an additional error appears related to the goodness of fit. Last, the term related to the uncertainty of the spectral sensitivity curve gives an additional error to the relative line intensities. This term is 5% for the observations presented, hence, it gives the main contribution to errors of the relative intensities of strong lines. All these components are summed squared, and the total errors have been propagated to calculate the errors of all derived parameters.
## 3 Results
### 3.1 Line intensities and element abundances
The relative intensities of all emission lines used for abundance determination in the 3 discussed H ii regions (Knots 1, 2, and 3), as well as the derived CH$`\beta `$, EWs of Balmer absorption lines, measured flux in H$`\beta `$ emission line, and the measured heliocentric radial velocities of each knot, are given in Tables 2 and 3. The position angle (PA) of the slit and the seeing $`\theta `$ during observation are also given for all spectra. The best 1D spectra for each of these regions, as well as the spectra of Knots 5 and 7, are shown in Fig. 2. For each of Knots 1, 2, and 3, we got repeated observations and can compare individual data. For Knot 1 the relative line intensities are very consistent in the two spectra for the most of lines except \[O ii\]$`\lambda `$3727 and H$`\alpha `$. This presumably is due to the different slit position on the asymmetric H ii region and also due to a factor of two different seeings, resulting in some contribution from the outer parts of this knot. For Knot 2 the first two spectra are very consistent with each other, while the 3-d indicates some additional ‘reddening’. The similar situation is for the 2nd spectrum of Knot 3 (PA=–57$`\mathrm{°}`$, $`\theta `$=1$`\stackrel{}{.}`$4), extracted from the same 2D spectrum. This seems to indicate some additional error in the spectral sensitivity curve for this night. However, the effect of this on the resulting abundances is small.
Extinction in the observed H ii regions is low: C(H$`\beta `$) $``$0.1. A larger value of C(H$`\beta `$), derived for one spectrum of Knots 2 and 3 ($``$0.2–0.3), is probably due to the worse quality of the spectral sensitivity curve for this night. This low C(H$`\beta `$) is quite consistent with the extinction data for most very metal-poor galaxies.
Chemical abundances and physical parameters are determined in the frame of the classical two-zone model of H ii region (Stasińska Stas90 (1990)), as described, e.g., in our recent papers by Pustilnik et al. (2004b ) and by Kniazev et al. (2004b ), which in turn follow the method described by Izotov et al. (1997b ). The derived individual abundances of oxygen for Knots 1, 2, and 3, as well as their T<sub>e</sub>, are given in Table 4. Since the \[S ii\] $`\lambda `$6717/$`\lambda `$6731 line ratio is higher than 1.4 for all spectra, the density N<sub>e</sub> = 10 cm<sup>-3</sup> was accepted in all our calculations (Aller Aller (1984)).
Since the line fluxes from these knots are rather small, the S-to-N in the principal line \[O iii\]$`\lambda `$4363 was at most $``$8. Therefore the majority of resulting abundances for individual observations are of medium accuracy. However, they all are consistent each to the other within their cited uncertainties. Thus, we average them and accept the mean values for each of the three regions: 12+$`\mathrm{log}`$(O/H)=7.23$`\pm `$0.05, 7.21$`\pm `$0.05, and 7.03$`\pm `$0.08, for Knots 1, 2, and 3, respectively. These mean values for different knots are consistent each to the other within their uncertainties, so we take their mean as a measure of the heavy element abundances of DDO 68 (at least, in this northern ‘ring’). If all seven measurements of O/H in these 3 knots are considered as the independent representatives of a unique value for all three knots, their weighted mean corresponds to 12+$`\mathrm{log}`$(O/H)=7.17$`\pm `$0.03. The 1st measurement for Knot 3, with 12+$`\mathrm{log}`$(O/H)=6.95$`\pm `$0.14, looks, however, a bit outlying. If we exclude this from the general mean, then the resulting mean value of the remaining six measurements (which probably is more conservative) corresponds to 12+$`\mathrm{log}`$(O/H)=7.21$`\pm `$0.03. We take this O/H as a characteristic value of DDO 68 in further discussion.
For Knot 1 the estimates of Ne, S, and Ar abundances were also derived. They correspond to $`\mathrm{log}`$(Ne/O)=–0.70$`\pm `$0.10, $`\mathrm{log}`$(S/O)=–1.21$`\pm `$0.16, and $`\mathrm{log}`$(Ar/O)=–2.20$`\pm `$0.07. Within rather large uncertainties, they are consistent with the abundance ratios for the most metal-deficient blue compact galaxies (BCGs) presented by Izotov & Thuan (IT99 (1999)).
Having the EWs of emission line H$`\beta `$, we can estimate the ages of starbursts in these three H ii regions apparently placed along a more or less regular oval. The average values of EW(H$`\beta `$) for Knots 1–3 are given in Table 5. For Knots 4, 6, and 7, we got spectra of relatively low S-to-N, so we do not discuss them in detail. However, their EWs(H$`\beta `$) are well measured and can also be used to date the starbursts at their respective locations.
We use the latest version (v5) of the models from Starburst99 (Leitherer et al. S99 (1999), Vazquez & Leitherer S99v5 (2005)) with the lowest metallicity given their ($`z`$=0.001) to estimate ages of instantaneous starbursts from the observed EWs(H$`\beta `$). For the Salpeter IMF with M<sub>low</sub> and M<sub>up</sub> of 0.1 M and 100 M, respectively, the derived ages of all emission knots are also given in Table 5. They range from 3 to 7 Myr. An independent check of the ages of all knots can be performed using the EW(H$`\alpha `$) from the obtained H$`\alpha `$ images. Both age estimates are very consistent to each other.
Knot 5 (see Fig. 1) appeared to be a compact young stellar cluster with well-seen Balmer series in absorption, besides H$`\alpha `$, on which some emission is superimposed. The radial velocity of this cluster derived on these lines (V<sub>hel</sub>=396$`\pm `$90 km s<sup>-1</sup>) is quite consistent with those for the H ii regions in the ‘ring’ and with the systemic velocity of DDO 68. The measured EWs of absorption lines H$`\beta `$, H$`\gamma `$, H$`\delta `$ are in the range of 4.4–5.8 Å (see Table 6). The value of EW(emis. H$`\alpha `$) after correction for the underlying absorption is 9.5$`\pm `$0.5 Å. The EW of absorption H$`\beta `$, in turn, after the correction for emission component, is 5.8$`\pm `$0.5 Å. EWs of H<sub>8</sub> and H<sub>9</sub> are measured with large uncertainties and are not taken for comparison. We compared all EWs of Balmer absorptions with the respective model values, presented by Gonzalez Delgado et al. (Rosa (1999)) for the cluster with metallicity $`z=0.001`$ (nearest to the observed one for DDO 68), assuming the instantaneous starburst with the Salpeter IMF with M<sub>low</sub> and M<sub>up</sub> of 1 and 80 M. The EWs of emission H$`\alpha `$ and the extrapolated one of H$`\beta `$ were compared to predictions of the latest version of Starburst99 (Vazquez & Leitherer S99v5 (2005)) for the Salpeter IMF with M<sub>up</sub> = 100 M. Both absorption and emission line EW values are consistent with the cluster age of T<sub>cluster</sub> $``$22–23 Myr.
Thus, we conclude that starbursts in 4 sites along the ‘ring’ are synchronized within the interval of 2 Myr. Only in Knot 5 did the star formation episode begin $``$18 Myr earlier than in the rest positions. One of the possible options for producing such a configuration of starbursts is the induced gas collapse behind the front of shock wave, generated by the previous strong starburst (e.g., van Dyk et al. Puche (1998), Efremov et al. Efremov02 (2002), Elmegreen et al. Elmegreen04 (2004)). However, the outlying age of Knot 5 suggests more complicated SF scenario in this region.
### 3.2 Results of imaging
The best quality image of DDO 68 obtained in these observations is the one in $`V`$-band. This is shown in Fig. 1 with the surface brightness cuts allowing the maximal extent of the galaxy to be followed. Also on this image we mark the knots for which the spectra, discussed in the previous section, were acquired. The net H$`\alpha `$ image in Fig. 1 displays the regions of current SF. We emphasise that besides the Northern ‘ring’ with diameter of $``$500 pc and rather bright emission-line knots along this ‘ring’, there is a similar structure in the Southern part of the galaxy that is connected with the Southern tail: Knot 7 in Fig. 1. As seen on the net H$`\alpha `$ image, it has the form of an almost closed ‘circle’ with a diameter of $``$270 pc. Its SE sector is significantly brighter than the rest of this ‘circle’ and consists of 3 adjacent knots. The brightest of them is quite elongated with the major axis oriented perpendicular to the arc to which it is connected. As discussed above, its starburst age is $``$4.5 Myr, which implies that the formation of the ‘ring’-like structures at both edges of DDO 68 have been synchronized by some global process.
The main issue to address based on our $`V,R,`$ and H$`\alpha `$ images is the age of the oldest visible stellar population. We built surface brightness profiles (SBPs) for $`V,R`$ images applying aperture photometry in circular apertures centred at the point near the geometrical centre of the main DDO 68 body (with J2000 coordinates of 09:56:45.8 +28:49:27). They are shown in Fig. 3.
The nebular emission in DDO 68 is not as strong and widespread as in SBS 0335–052 and I Zw 18. However, it somewhat affects the observed colours on the significant part of the colour profile in Fig. 3. The contribution of various features superimposed on the smooth distribution of the assumed underlying light is seen clearly on both SBPs (top panel) and the colour radial profile (bottom panel). The strongest features are seen clearly in the H$`\alpha `$ image (Fig. 1). The nebular emission shows the most contrast in the two ring-like regions at the N and S edges of the galaxy. They appear on both the SBP and the colour radial distribution as clear peaks at $`R`$40, 55, 80, and 100″. Their contribution within the main body ($`RR_{\mathrm{opt}}`$ = 31$`\stackrel{}{.}`$5) and adjacent regions is less pronounced, but is still significant till the radial distance of $`R`$34″.
We fitted the $`V,R`$ SBPs by the ‘minimal’ exponential disks, using for this only the ranges in radial coordinates where the contribution of the superimposed emission is minimal/absent, namely $`20\mathrm{}<R<40\mathrm{}`$ and $`60\mathrm{}<R<70\mathrm{}`$. The derived scalelengths of the disks in $`V`$ and $`R`$-bands do not differ significantly within their uncertainties. Their mean value is $`<`$$`\alpha `$$`>`$=15$`\stackrel{}{.}`$3$`\pm `$0$`\stackrel{}{.}`$3, corresponding to the linear value of 482$`\pm `$10 pc. We therefore accept the hypothesis that the underlying light is described by the unique exponential disk. The central surface brightnesses of the mean fitted disk are $`\mu _V^0`$ = 22.56 mag arcsec<sup>-2</sup> and $`\mu _R^0`$ = 22.34 mag arcsec<sup>-2</sup> with a total error of $``$0$`\stackrel{m}{.}`$06. The $`(VR)^0`$ colour radial profile along with the derived disk colour (solid line) are shown in the bottom panel of Fig. 3.
Total magnitudes of DDO 68 derived by the integration within a polygon mask including all galaxy light, with the removal of foreground stars and background galaxies, are as follows: $`V`$=14.46$`\pm `$0.05, $`R`$=14.20$`\pm `$0.05, with the resulting $`(VR)_{\mathrm{total}}`$=0.26$`\pm `$0.07. The relatively large errors appear from the uncertainties of the coefficients of the photometrical system, built on observations of the calibration field near QSO FBQ 0951+263. The total flux of DDO 68 in the H$`\alpha `$-line, derived within the same polygon mask, is (1.73$`\pm `$0.09)$`\times `$10<sup>-13</sup> erg cm<sup>-2</sup> s<sup>-1</sup>, which in turn corresponds to a total H$`\alpha `$ luminosity of 0.88$`\times `$10<sup>39</sup> erg s<sup>-1</sup>. We did not transfer this value to the traditional star formation rate (SFR), since for young starbursts the SFR derived through the commonly used formulas is highly uncertain (see Weilbacher & Fritze-v.Alvensleben (WvA (2001)). However, we compare this and other global parameters with those in I Zw 18 and SBS 0335–052 in the discussion.
Since the ‘minimal’ disk describes the widespread underlying stellar population over the whole galaxy volume, its colours can be considered as indicating the oldest stellar population in the galaxy. This colour $`(VR)_{\mathrm{disk}}`$=0$`\stackrel{m}{.}`$22, with the rms error of 0$`\stackrel{m}{.}`$07 due to the photometric system uncertainties. A small correction of $``$0$`\stackrel{m}{.}`$012 for the Galaxy extinction should also be applied before comparison with the evolutionary tracks of PEGASE.2 models (Fioc & Rocca-Volmerange PEGASE (1997, 1999)). We took the models for the metallicity of $`z`$=0.0004, the nearest to that of DDO 68. The Salpeter IMF is accepted with M<sub>low</sub> and M<sub>up</sub> of 0.1 M and 120 M, respectively. Various modes of SF (instantaneous and continuous) result in different age constraints on the oldest visible stellar population.
Taken at face value, the disk colour ($`(VR)^0`$ $``$0.21) with the cited error of $`\sigma _{(VR)}=`$0$`\stackrel{m}{.}`$07 corresponds to either 0.65–1.8–4.5 Gyr-old (lower and upper limits correspond to –1$`\sigma `$ and +1$`\sigma `$) stellar clusters for the continuous SF with constant SFR or to 0.11–0.12–1.0 Gyr-old cluster for the instantaneous starburst.
### 3.3 Properties of the young star cluster in Knot 5
The young star cluster (YSC) related to Knot 5 is one of a few known with such low metallicity that can be directly studied in great detail. The others include, e.g., some stellar clusters containing WR stars in I Zw 18 (Brown et al. Brown02 (1996)) or unresolved with HST and VLT young clusters in SBS 0335–052 E (Thuan et al. TIL97 (1997), Plante & Savage Plante02 (2002)). Apart from the difference in metallicity by a factor of $``$30, it resembles the young super star cluster in the bubble complex of NGC 6946 very well (Elmegreen et al. Elmegreen00 (2000), Efremov et al. Efremov02 (2002), Larsen et al. Larsen01 (2001)), which is situated at a similar distance of $``$6 Mpc. The latter was successfully observed with the HST and Keck HIRES to resolve its stellar content and velocity dispersion and to address its evolutionary status. While for the YSC in DDO 68, such observations are feasible in the future as well, here we summarise the properties that emerge from current observational data.
We can estimate the characteristic size of the Knot 5 cluster by assuming its intrinsic intensity distribution to be Gaussian with the respective FWHM<sub>YSC</sub>. Subtracting (quadratically) the PSF FWHM (0$`\stackrel{}{.}`$8) from the observed FWHM (1$`\stackrel{}{.}`$12), we obtained the FWHM<sub>YSC</sub> intrinsic diameter of $``$0$`\stackrel{}{.}`$8. The respective linear diameter of this cluster is 25 pc. For comparison, compact young clusters in SBS 0335–052 E, unresolved on the HST images (the PSF FWHM=0$`\stackrel{}{.}`$2) have diameters less than 50 pc (Thuan et al. TIL97 (1997)). The further analogy of the DDO 68 young star cluster with the mentioned above in NGC 6946 extends to their linear size. As described by Larsen et al. (Larsen01 (2001)), this cluster has a core with a radius of 1.3 pc, surrounded by an extended halo with the power-law luminosity profile. The halo half-light radius is 13 pc, close to half of the FWHM of the light profile of Knot 5 in DDO 68. With the age of 12-15 Myr (Efremov et al. Efremov02 (2002)), the young star cluster in NGC 6946 is, however, significantly younger.
The total $`V`$=19$`\stackrel{m}{.}`$40 of the YSC in DDO 68 corresponds to the absolute magnitude M$`{}_{}{}^{0}{}_{V}{}^{}`$=–9.66. By comparing this M$`{}_{}{}^{0}{}_{V}{}^{}`$ with that from the PEGASE.2 model for instantaneous starbursts of the metallicity $`z=0.0004`$ cluster with the Salpeter IMF for the age T=22.5 Myr (M<sub>V</sub>=2.30 per 1 M), we derived the mass of the young cluster of M<sub>YSC</sub>=6.0$`\times `$10<sup>4</sup> M. This cluster is the most massive in the ‘Northern ring’. The brightest among the others, Knot 1, with $`V`$ $``$20$`\stackrel{m}{.}`$3, has a star cluster mass of 4.4$`\times `$10<sup>3</sup> M.
The measured colour of Knot 5 $`(VR)^0`$ = 0$`\stackrel{m}{.}`$11$`\pm `$0$`\stackrel{m}{.}`$01$`\pm `$0$`\stackrel{m}{.}`$07, where the second component of the error comes from the photometric system, is redder by 0$`\stackrel{m}{.}`$09 than predicted by the PEGASE.2 model for its age ($`(VR)`$ = 0$`\stackrel{m}{.}`$02). This is most probably related to some systematics in the zero-points of the constructed photometric system (see Sect. 3.4).
### 3.4 Colour correction and improved age estimates
The main uncertainties of the derived $`(VR)`$ colours of the underlying disk and of the estimate of the age of the oldest visible stellar population are related to the zero-points uncertainty of the constructed photometrical system ($`\sigma _{(\mathrm{V}\mathrm{R})}`$ = 0$`\stackrel{m}{.}`$07). Having some independent observational data for DDO 68, we can now diminish this effect. For this we used two different approaches.
First, we used the spectral information (EWs of Balmer absorption and emission lines) for Knot 5 - the young compact star cluster in the northern ‘ring’, which implies its age of $``$22.5 Myr. Then, as shown at the end of the previous section, the observed $`(VR)`$ colour is 0$`\stackrel{m}{.}`$09 redder than predicted by the PEGASE.2 models for its respective age. This implies the presence of the additive small shift of 0$`\stackrel{m}{.}`$09 in the $`(VR)`$ zero-point due to the photometric system uncertainties.
Second, there is an independent indication of the reality of the $`(VR)`$ colour shift found above. This comes from the comparison of our measured $`(VR)_{\mathrm{total}}`$ = 0.26 and the colour $`(BV)_{\mathrm{total}}`$ = 0.17 from Karachentsev et al. (NGC (2004)) in Table 7, with $`(BR)_{\mathrm{total}}`$ = 0.32$`\pm `$0.03 from Hopp & Schulte-Ladbeck (HS95 (1995)). Combining the former, our value with the cited $`(BV)_{\mathrm{total}}`$ led to $`(BR)_{\mathrm{total}}`$ = 0$`\stackrel{m}{.}`$43, while using the above correction brought this parameter to value of 0$`\stackrel{m}{.}`$34, significantly more consistent with that of Hopp & Schulte-Ladbeck.
Now we apply this correction of the $`VR`$ zero-point and the respective reduction of its uncertainty roughly to 0$`\stackrel{m}{.}`$04 from comparison with the Hopp and Schulte-Ladbeck $`(BR)_{\mathrm{total}}`$ to the colours of the underlying population. The corrected $`(VR)^0`$ colour of the underlying disk is 0$`\stackrel{m}{.}`$12$`\pm `$0$`\stackrel{m}{.}`$04. This in turn corresponds (in the frame of the same PEGASE.2 models as above) to the ages of 200 – 450 – 900 Myr for the case of continuous SF with constant SFR, or to the ages of 100 – 105 – 115 Myr for the case of instantaneous starburst. The lower and upper limits correspond, as above, to $``$1$`\sigma `$ and +1$`\sigma `$ in the $`(VR)`$ colour.
Applying the above correction of $`(VR)`$, we somewhat improve the accuracy of the derived ages of the old population. However, due to the well known degeneracy of the colour evolutionary tracks for different SF laws, the estimated ages are still uncertain by a factor of $``$(2–8), if based only on the optical colours of unresolved stars, depending on the accepted SF law. Besides, the use of the other optical-NIR colours is important to check the results suggested by the study of only $`VR`$ colour.
The most reliable age estimates can be obtained from the colour-magnitude diagrams (CMD) of resolved stars. If the galaxy is indeed young, no RGB stars will be seen in the CMD, as in the case of I Zw 18 (Izotov & Thuan IZw18CMD (2004)). For the distance modulus of DDO 68 equal to 29.0 (for $`D=6.5`$ Mpc), stars on the tip of RGB should have $`V`$=25$`\stackrel{m}{.}`$0 and $`I`$=23$`\stackrel{m}{.}`$8.
## 4 Discussion and Conclusions
### 4.1 Properties of DDO 68 in comparison to I Zw 18
The weighted mean 12+$`\mathrm{log}`$(O/H)=7.21$`\pm `$0.03 for DDO 68 is very close to that in I Zw 18. Besides, DDO 68 is close to I Zw 18 on several other parameters. In particular, its absolute blue magnitude $`M_B`$ is only 0$`\stackrel{m}{.}`$56 fainter. Thus, both objects are very close on the L–Z diagram and, together with SBS 0335–052 E and W, comprise a group of galaxies significantly deviating from the general L–Z relation for BCGs (see, e.g., Pustilnik et al. 2003a ). Such deviations are also noticed for several luminous BCGs with 12+$`\mathrm{log}`$(O/H) $``$8.0 (e.g., Kunth & Östlin Kunth2000 (2000)), which are old galaxies, so the strong deviation from the general L–Z relation is not an exclusive property of the most metal-deficient BCGs. However, the objects of this small group deviate as well on the L–Z diagram from the region where the majority of known XMD galaxies are situated (see Kniazev et al. Kniazev03 (2003)). Based on this large deviation, Kniazev et al. (Kniazev03 (2003)) suggest that the small XMD subgroup of I Zw 18 and its ‘cousins’ are truly young galaxies. A similar conclusion was formulated for the same small subgroup of XMD galaxies by Guseva et al. (SBS1415 (2003)), based on the very blue $`(VI)`$ colours of their LSB components.
The M(H i)/L<sub>B</sub> ratio for DDO 68 is also sufficiently high (1.85), and is a factor of $``$1.4 higher than for I Zw 18. Its total H i mass (2.6$`\times `$10<sup>8</sup> M) is a factor of 1.6 higher than the H i mass of I Zw 18. Its integrated colour ($`BV`$)$`{}_{}{}^{0}{}_{\mathrm{tot}}{}^{}`$=0$`\stackrel{m}{.}`$17 (Makarova & Karachentsev Makarova98 (1998)) is close to that of I Zw 18 (0$`\stackrel{m}{.}`$12). Some other properties of these two galaxies are collected in Table 7 for comparison. Direct comparison of the current SFR (based on the H$`\alpha `$-luminosity) of DDO 68 and other similar galaxies is not meaningful, since this SFR changes significantly on a short timescale. However the specific values of H$`\alpha `$-luminosity related either to the blue luminosity or to the total H i mass are of interest. For further comparison, we take F(H$`\alpha `$) of I Zw 18 from Gil de Paz et al. (Gil03 (2003)).
In DDO 68 these parameters are significantly lower than for I Zw 18, indicating a lower current production of the ionizing radiation. Namely, L(H$`\alpha `$)/L<sub>B</sub> is $``$0.005 for DDO 68 and 0.035 - for I Zw 18. Similar, L(H$`\alpha `$)/M(H i) is 0.0016 (in units $`(M/L_\mathrm{B})`$) for DDO 68 and 0.028 – for I Zw 18. Interestingly, on these parameters DDO 68 is closer to SBS 0335–052 W, where these ratios are L(H$`\alpha `$)/L<sub>B</sub> = 0.02 and L(H$`\alpha `$)/M(H i) = 0.0027 (H$`\alpha `$-flux is from Pustilnik et al. 2004a ). It was noticed that based on its properties, SBS 0335–052 W is not a typical blue compact galaxy, but instead resembles dIrr, which based on the strength of the current ‘starburst’ is also correct for DDO 68. The close similarity of principal observational properties of DDO 68 and I Zw 18 suggests that DDO 68 may also be a young galaxy, that is it contains no stellar population with ages more than 1 Gyr.
### 4.2 Optical and H i morphology and their implications
The strongly disturbed optical morphology of DDO 68 is difficult to understand since the galaxy is rather distant from any potential disturber. The luminous galaxies are at least at $``$2.0 Mpc (NGC 2683 with M<sub>B</sub>=–20.4 at D=7.7 Mpc) and $``$2.5 Mpc (NGC 2903 with M<sub>B</sub>=–21.0 at D=8.9 Mpc, Karachentsev et al. NGC (2004)). The nearest known galaxy, a dwarf UGC 5427, with $`M_B`$ = –14.5 at $``$200 kpc in projection and with very close radial velocity, is too low-mass object to affect DDO 68. The neutral gas morphology and kinematics, as revealed by H i maps with 13$`\stackrel{}{.}`$5 angular resolution by Stil & Israel (2002a ; 2002b ), are also well disturbed with several maxima of density and with several holes with linear sizes of the order of 1 kpc.
One could consider two options to understand the optical morphology of DDO 68. The first is related to the strong disturbance by some separate, massive, poorly visible (or optically invisible) object. This could be a very low surface brightness galaxy, or an intergalactic H i cloud, similar to the case of HI 1225+01 (Salzer et al. Salzer91 (1991), Chengalur et al. Chengalur95 (1995)). The second option is related to an advanced merger, when two merging objects are already not seen as separate entities. However, the morphology of the joint material clearly indicates outflows/plumes as the characteristic features created by strong interaction. The importance of the merger channel for starbursts in BCGs was emphasised, e.g., by Östlin et al. (Ostlin01 (2001)). Of course, the combined option is also possible: for example, a well-advanced merger with an intergalactic H i cloud.
To distinguish the above options, the available H i maps are very useful. We also used overlays of the H i density map on the DDO 68 $`I`$-band and H$`\alpha `$ images, along with detailed data on H i velocity dispersion in the Ph.D. thesis of Stil (Stil99 (1999)). Analysis of all these data leads to the next conclusions. First, there is no indication of the presence of any separate massive H i object in the vicinity of DDO 68. Second, there is the large asymmetry of H i density distribution relative to the optical body of the galaxy. The H i density map can be divided conditionally into two curved chains of dense matter stretching roughly from North to South, with the typical separation between the chains of 0$`\stackrel{}{.}`$5–1$`\stackrel{}{.}`$0 (1–2 kpc). The optical body, including the Northern ‘ring’ and the Southern ‘tail’, coincides roughly with the Eastern H i chain, while practically no optical counterparts are seen for the Western H i chain (at least, on the SB level, comparable with the optical light seen along the Eastern H i chain).
The large asymmetry between the distributions of H i mass and the optical light suggests that the Western part of the H i body is unrelated directly to the optical galaxy and represents some external H i cloud in the process of merging with a gas-rich massive object, in which some star formation has already occurred and is still taking place. The elevated velocity dispersion in the four densest regions (by a factor of 2 to 3 relative to the undisturbed regions, which have $`\sigma _v`$=6–9 km s<sup>-1</sup>), all them being displaced to the Eastern half of the whole H i cloud, gives additional evidence of strong gas agitation due to the ongoing merging, as suggested by Elmegreen et al. (Elmegreen93 (1993)). If this hypothesis is correct, DDO 68 can represent the well-advanced stage of merging observed in the system HI 1225+01. The metallicity of gas in the two merging components can differ substantially, so the complex process of the gas mixing can create significant spatial inhomogeneities in the ISM metal content of DDO 68.
While the WSRT DDO 68 H i map does not show the extended tidal tails characteristic of many mergers, this is not an argument against the merger case. However, the tails are not likely likely to be seen because the available maps are relatively shallow, with the lowest H i density contour of 5$`\times `$10<sup>20</sup> atom cm<sup>-2</sup>. The extended tidal tails in the well-advanced mergers like those in NGC 7252, NGC 3921, NGC 3526 (e.g., Hibbard et al. gallery (2001)), are visible at the H i column densities of (one – a few) 10<sup>19</sup> atom cm<sup>-2</sup>. Thus, in the frame of the merger hypothesis, we predict that the deeper H i mapping of DDO 68 will uncover the extended gas tidal tails.
We also mention a more exotic variant of the first option, when the strong interaction of a gas-rich dwarf took place with a completely dark galaxy, which is a massive DM halo with the baryon mass fraction well below that typical of galaxies (Trentham et al. Trentham01 (2001)). As these authors suggest, this can be the trigger of starbursts in many isolated BCGs. However, to check this hypothesis observationally, they need to have very clear-cut predictions, since at first approximation it is difficult to distinguish this case from well-advanced merger.
### 4.3 Chemical evolution and environment
In the frame of Cold Dark Matter cosmology the galaxies formed in the regions of the very low density (voids) should be of lower mass and retarded in their formation epoch in comparison to the more massive galaxies from the average and the higher density regions (e.g., Gottlöber et al. Gott03 (2003)). Besides, in the frame of the widely accepted hierarchical galaxy formation scenario, their life cycles can differ quite a lot from the more common galaxies. Namely, for a selected galaxy, the probability of interactions and mergers during the cosmological time is scaled down in voids with the galaxy density as $`\rho _{\mathrm{galaxy}}^{+1}`$. Therefore, one could expect that some (small) fraction of void galaxies can survive in their nascent state, avoiding any significant interaction.
It is probably not by chance that one of the most metal-poor galaxies, DDO 68, is found in the very low-density ‘normal’ galaxy environment. Another XMD galaxy HS 0822+3542 is situated near the centre of the same Lynx-Cancer void (Pustilnik et al. 2003b ). At least one more dwarf galaxy with very low metallicity, KISSB 23 = KUG 0937+298 (12+$`\mathrm{log}`$(O/H)=7.65, Lee et al. Lee04 (2004)) is also situated within the same void, only at $``$650 kpc from DDO 68.
The nearest neighbour of DDO 68 is UGC 5427 (Sdm, B$``$14$`\stackrel{m}{.}`$9) at the angular distance of $``$1.8°, corresponding to $``$200 kpc in projection. Its radial velocity is only 7 km s<sup>-1</sup> lower than that of DDO 68. Karachentsev et al. (NGC (2004)) accept distances for DDO 68 and UGC 5427 derived by the method of the brightest stars (accuracies of 25%), of 5.9 and 7.1 Mpc, respectively. We suggest that both galaxies belong to the same ‘filament’ and the distance in between is on the order of their projected distance. Therefore, we accept their common radial distance as a mean of their individual determinations, consistent each to other: $`D_{\mathrm{mean}}=6.5`$ Mpc.
As discussed by Pustilnik et al. (2004a ), a substantial fraction of XMD galaxies are found either at the outskirts of galaxy aggregates (e.g., SBS 0335–052 E,W and Dw 1225+0152), or in voids (e.g., HS 0822+3542, HS 2236+1344, HS 0837+4717). While the statistics and properties of void galaxy population are still not well known, it is assumed from general consideration that the significantly reduced frequency of galaxy interactions in voids provides the most favourable conditions to allow for the stablest gas protogalaxies to survive as purely gas objects. If such very stable protogalaxies do exist in voids, they can be detected either as purely gas bodies through the blind H i surveys or through the Ly-$`\alpha `$ absorption of a background UV-bright object (e.g., Manning Manning02 (2002)), or even as young galaxies, if they recently have experienced the strong gravitational disturbance due to close collision with sufficiently massive objects.
Therefore, on one hand, it is not surprising that many ‘void’ XMD galaxies studied by us (a paper in preparation) show clear evidence of interactions and merging. This can be a kind of selection effect for the superstable objects, when starbursts in unevolved galaxies occur only due to the strong external disturbance. The case of DDO 68 is additional evidence for the common mode of such starbursts in the XMD galaxies situated in the very low density environment.
On the other hand, as explained by Kunth & Sargent (KS86 (1986)), even the first starburst in such a retarded object, due to the effective self-pollution in the young H ii regions, will be observed with a heavy element content comparable to Z $``$Z/30 for the starburst ages larger than 3 Myr.
### 4.4 Conclusions
Summarising the observational data and the discussion above, we have drawn the following conclusions:
* The oxygen abundances in three H ii regions at the northern edge of the dwarf galaxy DDO 68 are close one to another with the weighted mean of 12+$`\mathrm{log}`$(O/H)=7.21$`\pm `$0.03 dex.
* DDO 68 is one of the three galaxies (I Zw 18 and SBS 0335–052 W) with the lowest ISM metallicity, and the nearest galaxy with the same low metallicity (twice closer than I Zw 18).
* The deep photometry of DDO 68 indicates that its the reddest colours outside the central bright region and the regions of current/recent SF (where no nebular emission is seen) are very blue ($`(VR)^0`$=0.12 $`\pm `$0$`\stackrel{m}{.}`$04). Their comparison with the PEGASE.2 models for the evolving stellar clusters (with metallicity of $`z=0.0004`$) implies that the oldest visible stars contributing to the light of the LSB underlying disk have ages either of 100–115 Myr. in the case of instantaneous starburst, or of 200–900 Myr, in the case of continuous star formation with the constant SFR. Thus, $`VR`$ data indicate the possible youth of DDO 68. The data on other colours, as well on CMDs, are necessary, however, to check and strengthen this conclusion.
* The young star clusters (related, in particular to Knots 1 and 5 in the ‘Northern ring’) with the masses of (0.4–6.0)$`\times `$10<sup>4</sup> M are the nearest among all known such objects with the lowest metallicity. They are probably the best known clusters for which the direct spectroscopic study of individual massive stars with Z $``$Z/30 will be possible with the future giant ground-based and space telescopes. Study of their structure and dynamics is currently feasible with the HST and the high resolution spectroscopy at large telescopes.
* The optical morphology of DDO 68 indicates recent interaction despite no disturbing galaxy is evident in its surrounding. The most probable interpretation of this galaxy’s properties is a well-advanced merging of two gas-rich objects.
* DDO 68 is situated in the region with the low density of luminous galaxies near the rim of a nearby small Lynx-Cancer void. Apart from this galaxy, two more XMD dwarfs have been discovered in this void: HS 0822+3542 and KISSB 23. Several other XMD galaxies are found in other voids. This is consistent with theoretical expectations that dwarfs in the regions with the very low density of luminous galaxies (and the very reduced density of low-mass galaxies) may form/evolve with significant delay. This opens a new direction in searching for local young galaxies.
###### Acknowledgements.
The authors are pleased to thank V. Afanasiev and A. Moiseev for the help in the organization and preparation of observations with SCORPIO, Y. Izotov for useful comments, and J. Stil for kindly providing his Ph.D. thesis. We appreciate the criticism and suggestions of the referee, D. Kunth, which helped to improve the paper. We acknowledge the partial support from Russian state program ”Astronomy”. This research made use of the NASA/IPAC Extragalactic Database (NED), which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration. Use of the Digitized Sky Survey (DSS-II) is gratefully acknowledged. |
warning/0507/hep-th0507072.html | ar5iv | text | # Comment to ”Fine structure constant in the spacetime of a cosmic string
## Abstract
In the paper Phys. Lett. B 614 (2005) 140-142, F. Nasseri shows that the values of the fine structure constant reduces due to the presence of a cosmic string. In this comment I want to point out that this conclusion is not completely correct in the sense that the result obtained is valid only in a very special situation.
In a recent paper F. Nasseri analysed the Bohr’s atom in the cosmic string spacetime. The author claims that the fine structure constant in this system reduces by a factor of order $`8.736\times 10^{17}`$, consequently very small as the author says. In order to reach this conclusion the author considers separately two distinct forces acting on the electron: $`i)`$ the induced electrostatic self-force on the electron due to the presence of a cosmic string , and $`ii)`$ the electrostatic one between the electron and the proton. So the total force will be given by the sum of these two forces. Although being not explicitly said, the author assumes that the proton is placed on the cosmic string considered as an ideal one. The point that I want to call attention is to the fact that the induced force on the electron, Eq. $`(11)`$ of , is perpendicular to the cosmic string assumed along the $`z`$axis. So the correct expression to it is:
$`\stackrel{}{F}={\displaystyle \frac{k\pi G\mu e^2}{4\rho ^2}}\widehat{\rho }.`$ (1)
$`\rho `$ being the distance between the electron and the string. Also I have included the factor $`k`$ in the numerator of the right hand side of the above equation, missed in the Eq. $`(11)`$ of <sup>1</sup><sup>1</sup>1In Ref. the author has used the mks units where $`k=1/(4\pi ϵ_0)`$..
As to the electrostatic force, it is along the radial distance between the electron and the proton. Considering the origin of the referential system on the proton, this force reads:
$`\stackrel{}{F}={\displaystyle \frac{ke^2}{r^2}}\widehat{r}.`$ (2)
So in principle both forces are not always at the same direction.
The Eq. $`(15)`$ of is correct only when electron and proton are in the plane orthogonal to the string. Consequently the correction on the fine structure constant presented in Eq. $`(17)`$ is not general.
Acknowledgment
I want to thanks the referee for pertinent comments on this Comment. This work was partially supported by CNPq. |
warning/0507/hep-ph0507067.html | ar5iv | text | # A Prediction from the Type III See-saw Mechanism
## Abstract
A simple ansatz that is well-motivated by group-theoretical considerations is proposed in the context of the type III neutrino see-saw mechanism. It results in predictions for $`m_s/m_b`$ and $`m_\tau /m_b`$ that relates these quantities to the masses and mixings of neutrinos.
Simple unified models based on $`SO(10)`$ and related groups can lead to the so-called “type III see-saw mechanism” for neutrino masses type3 . In the most general case the type III mechanism leads to a light neutrino mass matrix given by the formula $`M_\nu =(M_NH+H^TM_N^T)(u/\mathrm{\Omega })`$, where $`M_N`$ is the Dirac mass matrix of the neutrinos, $`H`$ is a dimensionless complex three-by-three matrix and $`u/\mathrm{\Omega }`$ is the ratio of a weak-scale vacuum expectation value to a GUT-scale vacuum expectation value (VEV). In a subsequent paper the type III see-saw mechanism was shown to have certain advantages for leptogenesis, in particular allowing resonant enhancement without fine-tuning the form of neutrino mass matrices type3leptogen . In the simplest case, where a minimal set of Higgs fields breaks $`BL`$, one has $`H=I`$ and the type III see-saw formula takes the simple form
$$M_\nu =(M_N+M_N^T)\frac{u}{\mathrm{\Omega }}.$$
(1)
The main problem in constructing predictive models of neutrino masses and mixings with the usual “type I” see-saw formula type1 , $`M_\nu =M_NM_R^1M_N^T`$, is to relate the Majorana mass matrix of the right-handed neutrinos $`M_R`$, with its six complex parameters, to measurable quantities. There are very special models, such as the recently much studied “minimal $`SO(10)`$ models”, where there is such a relationship minimal . (For an exhaustive list of references on “minimal $`SO(10)`$ models” see Aulakh:2005bc .) And the study of leptogenesis may tell us something about the structure of $`M_R`$ (although leptogenesis has only a single data point to work with). In general, however, the lack of information about $`M_R`$ is a problem for the predictivity of type I see-saw models. (The so-called “type II see-saw mechanism” type2 assumes the existence of $`SU(2)_L`$-triplet Higgs fields with small VEVs that couple directly to $`\nu _L\nu _L`$. About the type II mechanism we have nothing to say in this paper.)
What makes the simplest version of the type III formula, given in Eq. (1), so remarkable and appealing is that it does not involve the masses of the superheavy right-handed neutrinos at all. As a consequence, the simplest type III formula opens the possibility of constructing models of quark and lepton masses that are extremely predictive. In particular, in models based on $`SO(10)`$ or other groups that unify an entire family within a single multiplet, the Dirac mass matrix of the neutrinos $`M_N`$ is typically closely related by the grand-unification symmetries to the mass matrices (also of Dirac type, of course) of the up quarks, down quarks and charged leptons, which we will denote respectively as $`M_U`$, $`M_D`$, and $`M_L`$. It is therefore possible in many models (for examples, see ab ; bpw ; ardhs ) to predict the matrix $`M_N`$ from a knowledge of the masses and mixings of the quarks and the masses of the charged leptons. This would allow, if Eq. (1) holds, the complete prediction of the mass ratios and mixing angles of the neutrinos with no free parameters.
In this paper we will not be so ambitious. We have not found so far a full three-family model that is as predictive as that and where all the predictions (or “postdictions”) are consistent with experiment. Rather, as an illustration of the possibilities of the type III framework, we will present here a simple ansatz for the heavier two families that is well motivated by group-theoretical considerations. This ansatz leads to two interesting predictions that are consistent with present experimental data. Before presenting the ansatz, we very briefly review the type III see-saw mechanism and formula.
In models based on $`SO(10)`$, there are two ways that the right-handed neutrinos $`N_i^c`$ ($`i=1,2,3`$) can get mass, either through a renormalizable term such as $`\mathrm{𝟏𝟔}_i\mathrm{𝟏𝟔}_j\overline{\mathrm{𝟏𝟐𝟔}}_H`$, or through a higher-dimension effective operator such as $`\mathrm{𝟏𝟔}_i\mathrm{𝟏𝟔}_j\overline{\mathrm{𝟏𝟔}}_H\overline{\mathrm{𝟏𝟔}}_H/M_{GUT}`$. The former allows automatic conservation of “matter parity”, whereas the latter makes do with smaller multiplets of Higgs fields. In the latter case, the effective $`d=5`$ operator arises most simply from integrating out three or more $`SO(10)`$-singlets, which we will denote by $`\mathrm{𝟏}_a`$ or $`S_a`$, that have the couplings $`F_{ia}\mathrm{𝟏𝟔}_i\mathrm{𝟏}_a\overline{\mathrm{𝟏𝟔}}_H`$ and $`(M_S)_{ab}\mathrm{𝟏}_a\mathrm{𝟏}_b`$. If only the Standard-Model-singlet component of the $`\overline{\mathrm{𝟏𝟔}}_H`$ has a non-zero VEV, and we denote it by $`\mathrm{\Omega }M_{GUT}`$, then one has the familiar “double see-saw” mass matrix:
$$_{neutrino}=(\nu _i,N_i^c,S_a)\left(\begin{array}{ccc}0& (M_N)_{ij}& 0\\ (M_N^T)_{ij}& 0& F_{ib}\mathrm{\Omega }\\ 0& F_{aj}\mathrm{\Omega }& (M_S)_{ab}\end{array}\right)\left(\begin{array}{c}\nu _j\\ N_j^c\\ S_b\end{array}\right).$$
(2)
By integrating out the superheavy fields $`N_i^c`$ and $`S_a`$, one obtains $`M_\nu =M_NM_R^1M_N^T`$, where $`M_R=(F\mathrm{\Omega })M_S^1(F\mathrm{\Omega })^T`$. This is just the type I see-saw formula, with an effective $`M_R`$.
Now, if we assume that the $`SU(2)_L`$-doublet Higgs field contained in $`\overline{\mathrm{𝟏𝟔}}_H`$ also gets a non-zero VEV (and there is no fundamental reason why it should not), and we denote it by $`u`$, then the double see-saw mass matrix takes the form:
$$_{neutrino}=(\nu _i,N_i^c,S_a)\left(\begin{array}{ccc}0& (M_N)_{ij}& F_{ib}u\\ (M_N^T)_{ij}& 0& F_{ib}\mathrm{\Omega }\\ F_{aj}^Tu& F_{aj}^T\mathrm{\Omega }& (M_S)_{ab}\end{array}\right)\left(\begin{array}{c}\nu _j\\ N_j^c\\ S_b\end{array}\right).$$
(3)
In this case, it is easy to show that the effective mass matrix of the light neutrinos takes the form:
$$M_\nu =M_NM_R^1M_N^T(M_N+M_N^T)\frac{u}{\mathrm{\Omega }},$$
(4)
where, as before, $`M_R=(F\mathrm{\Omega })M_S^1(F\mathrm{\Omega })^T`$. The first term is the usual type I see-saw contribution, and the second term is the type III see-saw contribution. (The origin of the type III term can be simply understood as follows. One can eliminate the $`\nu S`$ and $`S\nu `$ entries in Eq. (3), i.e. the entries $`Fu`$ and $`F^Tu`$, by doing a rotation of the $`(\nu _i,N_i^c)`$ basis by an angle $`\theta \mathrm{tan}\theta =u/\mathrm{\Omega }`$. That reduces the matrix in Eq. (3) to the same form as Eq. (2), but with the zeros replaced by terms of the type III form.) Both the type I and the type III terms in Eq. (3) are formally of order $`M_W^2/M_{GUT}`$. However, since the elements of $`M_N`$ are actually small compared to $`M_W`$ because of small Yukawa couplings (except perhaps for the third family), and since $`M_N`$ comes in quadratically in the type I term but only linearly in the type III term, one might expect the type III term to dominate for generic values of the parameters. Moreover, in the limit that the elements of $`M_S`$ are small compared to the GUT scale, the type I contribution becomes small. As was pointed out in type3leptogen , that is a good limit for the purposes of enhancing leptogenesis. It is therefore plausible that one can neglect the type I term, and we shall do so.
Now let us turn to the ansatz for the various Dirac mass matrices. Suppose that these have the form (neglecting the small masses of the first family)
$$M_U=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& a\\ 0& b& 1\end{array}\right)m_U,M_D=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& c\\ 0& d& 1\end{array}\right)m_D,$$
(5)
$$M_N=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& g\\ 0& h& 1\end{array}\right)m_U,M_L=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& e\\ 0& f& 1\end{array}\right)m_D,$$
(6)
where the “texture zero” in the 22 elements can be enforced by an abelian family symmetry, either discrete or continuous. We will say more on this later. And further suppose that the entries satisfy the conditions
$$a+b=g+h,c+d=e+f.$$
(7)
The relations given in Eq. (7) are not arbitrary, but follow from group-theory if the elements of the mass matrices come from no operators except of the following simple types:
(1) $`\mathrm{𝟏𝟔}_i\mathrm{𝟏𝟔}_j\mathrm{𝟏𝟎}_H`$,
(2) $`\mathrm{𝟏𝟔}_i\mathrm{𝟏𝟔}_j\mathrm{𝟏𝟐𝟎}_H`$,
(3) $`\mathrm{𝟏𝟔}_i\mathrm{𝟏𝟔}_j\mathrm{𝟏𝟎}_H\mathrm{𝟒𝟓}_H/M_{GUT}`$.
(4) $`\mathrm{𝟏𝟔}_i\mathrm{𝟏𝟔}_j\mathrm{𝟏𝟔}_H^{}\mathrm{𝟏𝟔}_H/M_{GUT}`$,
Eq. (7) is satisfied no matter how many operators there are of any of these types. Any operator of type (1) gives $`a=g`$, $`b=h`$, $`c=e`$, and $`d=f`$, thus satisfying Eq. (7). Any operator of type (2) gives contributions that are flavor-antisymmetric (since the $`\mathrm{𝟏𝟐𝟎}`$ is in the antisymmetric product of two spinors). Consequently, it gives $`a+b=0`$, $`c+d=0`$, $`e+f=0`$, and $`g+h=0`$, thus also satisfying Eq. (7) in a trivial way.
Any operator of type (3) gives contributions of the form $`f_if_j^cv_f[\alpha Q(f)+\beta Q(f^c)]`$. Here $`Q`$ is that generator of $`SO(10)`$ to which the VEV of the adjoint Higgs field ($`\mathrm{𝟒𝟓}_H`$) is proportional; $`Q(f)`$ is the value of this charge for the fermion $`f`$ ($`=u,d,\mathrm{}^{},\nu `$); $`v_f=v_u`$ or $`v_d`$ depending on whether $`f`$ is of the weak-isospin up or down type; and the coefficients $`\alpha `$ and $`\beta `$ depend on the way the $`SO(10)`$ indices are contracted in the operator. Thus an operator of type (3) will give, for instance, $`c+d(\alpha +\beta )(Q(d)+Q(d^c))v_d`$ and $`e+f(\alpha +\beta )(Q(\mathrm{}^{})+Q(\mathrm{}^+))v_d`$. Since the terms $`d_id_j^cH_d`$ and $`\mathrm{}^{}\mathrm{}^+H_d`$ must be invariant under the charge $`Q`$, it follows that $`Q(d)+Q(d^c)=Q(H_d)=Q(\mathrm{}^{})+Q(\mathrm{}^+)`$, and so $`c+d=e+f`$, satisfying Eq. (7). In the same way it is easily seen that $`a+b=g+h`$.
Finally, consider an operator of type (4). One of the spinor Higgs fields (say the unprimed one) gets a superlarge VEV that breaks $`SO(10)`$ down to $`SU(5)`$. The effective operator that results is then of the form $`(\alpha \mathrm{𝟏𝟎}_i\overline{\mathrm{𝟓}}_j+\beta \overline{\mathrm{𝟓}}_i\mathrm{𝟏𝟎}_j)\overline{\mathrm{𝟓}}_H`$, where the coefficients depend on the contraction of $`SO(10)`$ indices in the original operator. This gives no contribution to $`a`$, $`b`$, $`g`$, and $`h`$, and gives contributions to the other paramaters of the form $`c=f`$ and $`d=e`$ (note the transposition between $`M_D`$ and $`M_L`$). Again, such contributions satisfy Eq. (7).
Simple low-dimension operators that could give contributions not satisfying Eq. (7) are $`\mathrm{𝟏𝟔}_i\mathrm{𝟏𝟔}_j\overline{\mathrm{𝟏𝟐𝟔}}_H`$ (if the $`SU(5)`$ $`\overline{\mathrm{𝟒𝟓}}`$ contained in the $`\overline{\mathrm{𝟏𝟐𝟔}}_H`$ got a non-zero VEV), and $`\mathrm{𝟏𝟔}_i\mathrm{𝟏𝟔}_j\overline{\mathrm{𝟏𝟔}}_H^{}\overline{\mathrm{𝟏𝟔}}_H/M_{GUT}`$.
One might ask why we do not include the effects of operators of even higher dimension, such as $`\mathrm{𝟏𝟔}_i\mathrm{𝟏𝟔}_j\mathrm{𝟏𝟎}_H\mathrm{𝟒𝟓}_H^n/M_{GUT}^n`$, which are not obviously smaller than the dimension-five operators that we included in our analysis, and which would not satisfy Eq. (7) in general. Such operators ought indeed to be present. However there are reasons that one might expect them to be small, as we now explain. Consider the operator $`\mathrm{𝟏𝟔}_2\mathrm{𝟏𝟔}_3\mathrm{𝟏𝟎}_H\mathrm{𝟒𝟓}_H/M_{GUT}`$, which will contribute to the 23 and 32 elements in our illustrative model. Since these elements are somewhat small compared to the 33 elements, either the effective Yukawa couplings in this term are small or the ratio $`\mathrm{𝟒𝟓}_H/M_{GUT}`$ must be, or both. This operator can arise from integrating out a pair of multiplets $`\mathrm{𝟏𝟔}^{}+\overline{\mathrm{𝟏𝟔}}^{}`$ that have GUT-scale mass, as follows. Suppose the terms $`a\mathrm{𝟏𝟔}_3\mathrm{𝟏𝟔}^{}\mathrm{𝟏𝟎}_H`$, $`b\mathrm{𝟏𝟔}_2\overline{\mathrm{𝟏𝟔}}^{}\mathrm{𝟒𝟓}_H`$, and $`M\mathrm{𝟏𝟔}^{}\overline{\mathrm{𝟏𝟔}}^{}`$. Integrating out the primed fields gives an effective operator $`ab\mathrm{𝟏𝟔}_2\mathrm{𝟏𝟔}_3\mathrm{𝟏𝟎}_H\mathrm{𝟒𝟓}_HM^1[1+|b\mathrm{𝟒𝟓}_H/M)^2|^{1/2}`$. If $`b`$ or $`\mathrm{𝟒𝟓}_H/M`$ are small, then the higher order operators are highly suppressed. This is not to say that operators of higher dimension must always be unimportant, but it is a plausible assumption easily implemented that they can be neglected.
To return to the texture zero in the 22 elements, it could be enforced, for example, by a $`U(1)`$ family symmetry under which the $`\mathrm{𝟏𝟔}_3`$, $`\mathrm{𝟏𝟎}_H`$ and $`\mathrm{𝟏𝟔}_H^{}`$ are neutral; the $`\mathrm{𝟏𝟔}_2`$ has charge $`+1`$; and the $`\mathrm{𝟒𝟓}_H`$, $`\mathrm{𝟏𝟔}_H`$, and $`\mathrm{𝟏𝟐𝟎}_H`$ have charge $`1`$.
Given the simplest type III form (Eq. (1)), and the ansatz of Eqs. (5), (6), and (7), one has
$$M_\nu =\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& (a+b)\\ 0& (a+b)& 2\end{array}\right)\frac{um_U}{\mathrm{\Omega }}.$$
(8)
$$M_L=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& e\\ 0& (c+de)& 1\end{array}\right)m_D,$$
(9)
with the quark matrices $`M_U`$ and $`M_D`$ given by Eq. (5). Consequently, to the extent that we can ignore the first family, the five parameters $`a`$, $`b`$, $`c`$, $`d`$, and $`e`$ determine the following mass ratios and mixings of the second and third families: $`m_c/m_t`$, $`m_s/m_b`$, $`V_{cb}`$, $`m_\mu /m_\tau `$, $`m_2/m_3`$ (the neutrino mass ratio), $`U_{\mu 3}\mathrm{sin}\theta _{atm}`$, and $`m_\tau /m_b`$. Therefore there are two predictions. (We assume all the parameters are real.)
What we have done is use the values of the five quantities $`m_c/m_t`$, $`m_\mu /m_\tau `$, $`V_{cb}`$, $`m_2/m_3`$, and $`\theta _{atm}`$ to solve for the five parameters $`a`$, $`b`$, $`c`$, $`d`$, and $`e`$. Then we have used the resulting values of those parameters to “predict” the values of $`m_s/m_b`$ and $`m_\tau /m_b`$ at the GUT scale. For the first three inputs ($`m_c/m_t`$, $`m_\mu /m_\tau `$, and $`V_{cb}`$), which are fairly well known, we have taken the central experimental values and run them up to the GUT scale, assuming low energy supersymmetry. The running depends significantly on the value of $`\mathrm{tan}\beta `$, and so we make a predictions for a particular set of values of $`\mathrm{tan}\beta `$ that span the interesting range: 2, 3, 10, 25, 40, and 57. The other two inputs ($`m_2/m_3`$ and $`\theta _{atm}`$) come from neutrino oscillation experiments (see the reviews neutrinos ; Fogli:2005cq ) and have rather large error bars. (For example, $`\theta _{atm}=45^{}\pm 6^{}`$ at $`M_Z`$.) We have assumed hierarchical spectrum for neutrino masses with $`m_1<0.007`$ eV. Under this assumption the RGE evolved values of $`m_2/m_3`$ and $`\theta _{atm}`$ at the GUT scale remain within 3% of their low-scale values even for large $`\mathrm{tan}\beta `$. (For relevant renormalization group equations see Babu:1993qv ; Antusch:2003kp .) Hence, we drop their running and allow these two inputs to vary within the experimentally allowed range and plot our predictions for $`m_s/m_b`$ and $`m_\tau /m_b`$ as a function of them in Fig. 1.
We take the experimental values of the quarks from Ref. quarkmasses , except for $`m_s`$ for which we use the results of lattice calculations as given in Ref. lattice and double the error as suggested in Ref. Kim:2004ki . The values of the CKM angles and the charged lepton masses are taken from PDG 2004 pdg2004 . In presenting our results for $`m_s/m_b`$ and $`m_\tau /m_b`$ in Fig. 1, we give the percentage by which the predicted GUT values differ from the RGE-evolved experimental central values.
In doing the renormalization group running we assume that all the sparticles have mass of 1 TeV. From $`M_Z`$ to 1 TeV, the running is done at one loop, assuming the Standard Model with two Higgs doublets. From 1 TeV to the GUT scale (taken to be $`2\times 10^{16}`$ GeV) we do a two-loop running assuming the MSSM. The gauge coupling constants are taken from PDG 2004 pdg2004 . We present one example of the RGE evolution in Table 1.
It should be noted that, even with the assumption that we are making that the parameters $`a`$, $`b`$, $`c`$, $`d`$, and $`e`$ are real, there are discrete ambiguities of the relative signs of these parameters. (The overall sign does not matter.) The choice that gives by far the best fits is $`(a,b,c,d,e)=\pm (,+,+,,)`$. A typical set of values is $`a0.00455`$, $`b+0.9`$, $`c+0.04`$, $`d0.45`$, $`e0.55`$, and $`fc+de+0.14`$.
Note that the value of $`a`$ is very small. It is this that accounts for the smallness of $`m_c/m_t`$. One way that $`a`$ might be small naturally (i.e. without fine-tuning) using only the set of operators that satisfy Eq. (7) is by means of an operator of the form $`\mathrm{𝟏𝟔}_2\mathrm{𝟏𝟔}_3\mathrm{𝟏𝟎}_H\mathrm{𝟒𝟓}_H/M_{GUT}`$, where $`\mathrm{𝟒𝟓}_HQ=I_{3R}+ϵ(BL)`$, $`ϵ1`$, where $`I_{3R}`$ and $`BL`$ are the familiar $`SO(10)`$ generators ($`I_{3R}`$ the diagonal generator of $`SU(2)_R`$ in the Pati-Salam subgroup, and $`BL`$ the baryon minus lepton number), and where the fields are contracted in such a way that this generator $`Q`$ acts on the field $`\mathrm{𝟏𝟔}_2`$. (This would happen, for instance if the effective operator came from integrating out a $`\mathrm{𝟏𝟔}^{}+\overline{\mathrm{𝟏𝟔}}^{}`$ having the couplings $`\mathrm{𝟏𝟔}_2\overline{\mathrm{𝟏𝟔}}^{}\mathrm{𝟒𝟓}_H`$, $`\mathrm{𝟏𝟔}_3\mathrm{𝟏𝟔}^{}\mathrm{𝟏𝟎}_H`$, and $`M\mathrm{𝟏𝟔}^{}\overline{\mathrm{𝟏𝟔}}^{}`$.) This operator would give off-diagonal mass terms for the up quarks proportional to $`Q(u)u_3^cu_2+Q(u^c)u_2^cu_3`$. Since, $`I_{3R}(u)=0`$, this would give $`a/b=O(ϵ)`$.
The values of $`m_s/m_b`$ that we predict are satisfyingly close to the experimental (lattice) results. A couple of things should be noted in this regard. First, it was long thought that the Georgi-Jarlskog gj relation $`(m_s/m_b)_{GUT}=\frac{1}{3}(m_\mu /m_\tau )_{GUT}`$ gave a good fit to the data in SUSY GUT models. However, the recent lattice calculations have given results for $`m_s`$ that are typically only about 0.6 times the typical values that had been obtained by previous methods. Because of that, many models which were constructed in the past to give the Georgi-Jarlskog result, would be off from the current central experimental/lattice results for $`m_s/m_b`$ by about $`+60\%`$. That compares to the values we are getting, which agree with the current central value of $`m_s/m_b`$ for some of the allowed $`(m_2/m_3)(\theta _{atm})`$ parameter space, and are within $`20\%`$ for a large part of that space.
A second point is that inclusion of the first family is likely to push up the predicted value of $`m_s/m_b`$ by about $`5\%`$. The reason is that empirically the relation for the Cabbibo angle $`\theta _C\sqrt{m_d/m_s}`$ is known to work very well cabbibo . As is well-known, this formula arises naturally if the 11 element of the down quark mass matrix vanishes and the 12 and 21 elements are approximately equal fritzsch . But then diagonalizing the 12 block of the down quark mass matrix will push up the value of the 22 element by a factor of $`(1+|m_d/m_s|)`$.
In any event, we see that further improvement in the measurement of the $`\theta _{atm}`$, $`\delta m_{atm}^2`$, $`\delta m_{sol}^2`$, and the lattice results for $`m_s`$, together with an eventual determination of $`\mathrm{tan}\beta `$ will allow our simple ansatz, given in Eq. (7), to be tested. |
warning/0507/astro-ph0507211.html | ar5iv | text | # A hot story with many tales: the qWR star HD~45166 Based on observations made at the 1.5m ESO telescope at La Silla, Chile, and at Laboratório Nacional de Astrofísica/CNPq, Brazil
## 1 Introduction
HD~45166 has been observed since 1922, being an enigma that up to now has not been deciphered. Anger (ange (1933)) described observations done between 1922 and 1933, showing that the emission spectrum is highly variable. The He ii 4686 Å line is always present, with variable intensity, while the lines of N iii 4640 Å and N iv 4058 Å appear and disappear completely. The star has been observed ever since, but no photometric or spectroscopic periodic variation has been discovered, in spite of the suspicions that it is, in fact, a binary system.
The spectral classification assigned to this object has been varying along the years. Anger (ange (1933)) proposed that the object is a Wolf-Rayet of the type WN. Neubauer & Aller (neub (1948)) classified it as a W7n. Morgan et al. (morg (1955)) classified it as Bpe while Hiltner & Schild (hilt2 (1968)) recovered the WR classification. Hiltner (hilt1 (1956)) described the object as having high excitation emission lines superposed upon an approximately O9 spectral type. Heap & Aller (heap (1978)), in a frequently cited (see for instance Willis & Stickland will1 (1983) – WS83) but never published paper, classified the star as B8 V + qWR, that is, a binary with a B8 V component and a ”quasi” Wolf-Rayet one. WS83 proposed a classification of B8 V + SdO – a main sequence star and a hot sub-dwarf. The Sixth Catalog of Galactic Wolf- Rayet Stars (van der Hucht et al. huch1 (1981)) recovers for HD~45166 the classification proposed by Heap & Aller (heap (1978)) (B8 V + qWR) and puts it, together with V~Sge, in the category of low mass WR, composed by only these two objects. This category does not exist in the Seventh Catalog of Galactic WR Stars (van der Hucht huch2 (2001)), since neither HD~45166 nor V~Sge are included in this catalog.
The star differs from a classic Population I WR in the sense that its emission lines are very narrow (typically full width at half maximum (FWHM) = 300 km s<sup>-1</sup> for the He ii lines). Besides, the object presents both characteristics of WN and WC simultaneously and seems to have a significant abundance of Hydrogen judging from the He ii Pickering decrement (Heap & Aller heap (1978); WS83). The object has a much smaller luminosity than a population I WR.
Van Blerkom (bler (1978)) made an analysis of the Hydrogen and Helium lines with the hypothesis that it is a population I WR object. He concluded that the WR component has a radius of 1 R and has a small size envelope that expands with a velocity of 150 km s<sup>-1</sup>, which results in number densities of He ii of about $`10^{11}`$ cm<sup>-3</sup> and mimics the environment of a WR envelope. He found that the wind density is $`N(H)=1.6\times 10^{12}`$ cm<sup>-3</sup> and $`N(He)=3.0\times 10^{11}`$ cm<sup>-3</sup>, and the mass loss rate, $`4.5\times 10^8`$ M year<sup>-1</sup>. On the other hand, WS83 obtained, on the basis of ultraviolet observations with the IUE satellite, the following parameters for the WR component: M$`{}_{\mathrm{V}}{}^{}=0.21`$; T$`{}_{\mathrm{eff}}{}^{}=\mathrm{60\hspace{0.17em}000}`$ K; $`\mathrm{log}`$(L/L)=3.84; R=0.77 R and M<sub>1</sub>=0.5 M. They obtained, also, a distance of 1 208 pc from spectroscopic parallax of the secondary star and a reddening of $`E(BV)`$=0.15 derived from the 2200 Å interstellar extinction band.
Willis et al. (will2 (1989)) – WHSH89 – based on a sequence of high resolution IUE spectra, found an extreme constancy of the UV continuum, in contrast to the variability of the line spectrum. On long time scales they found variability in the intensities of the C iv, N iv, N v and He ii emission lines. On short time scales, however, they found variability exclusively in the C iv 1559 Å doublet absorption lines. These short time scale variations ($`<1`$ day) are seen as discrete absorption components (DACs) characterized by two main features, with mean velocities of $`950`$ km s<sup>-1</sup> and $`750`$ km s<sup>-1</sup>. DACs are usually found in stars with P Cygni profiles. They can also be seen in luminous OB stars as well as in some SdO stars (see Brown et al. brown (2004) for a recent review on DACs). In HD~45166, both components are observed to migrate in velocity with an acceleration of $`140`$ cm s<sup>-2</sup>, a recurrence time of 1.60 $`\pm `$ 0.15 days and average lifetime of approximately 3 days. WHSH89 considered that these DACs in HD~45166 could result from structural changes in the wind, produced by radiative instabilities.
Until now photometric variability has not been detected in HD~45166. Ross (ross (1961)) made $`UBV`$ observations on 24 nights and did not find any variation. WHSH89 performed 36 hours of UV observations with the IUE satellite and did not find any variation with amplitude larger than 0.02 magnitudes, that is, within the accuracy limit of the data. The photometric measurements published in the literature are summarized in the Table 1.
There are several properties in common between HD 45166 and the V Sge stars: Van der Hucht et al. (huch1 (1981)) already compared the system to V~Sge. With the parameters proposed by WS83 its mass ratio is inverted as in the V Sge stars. Moreover, its emission line properties are quite similar. Therefore, we included HD 45166 in our ongoing search program for V Sge stars.
The V Sge stars are a group of 4 stars defined by Steiner & Diaz (stei98 (1998)): V~Sge (Herbig et al. herb (1965); Diaz diaz99 (1999)), WX~Cen (Oliveira & Steiner olivwx (2004); Diaz & Steiner diaz95 (1995)), V617~Sgr (Steiner et al. stei99 (1999); Cieslinski et al. cies (1999)) and DI~Cru (Oliveira et al. olivdi (2004); Veen et al. veen (2002)). They are characterized by the presence of strong emission lines of O vi and N v. Besides, He ii 4686 Å is at least two times more intense than H$`\beta `$. The V Sge stars are very similar to the Close Binary Supersoft X-rays Sources – CBSS – common in Magellanic Clouds, but not that frequent in the Galaxy. The CBSS are interpreted as binary systems with a white dwarf that presents hydrostatic Hydrogen nuclear burning on its surface. This burning is due to the high mass transfer rate which is a consequence of the inverted mass ratio (see Kahabka & van den Heuvel kah (1997), for a review and references).
The similarities between the spectra of HD~45166 and of the V Sge stars are the simultaneous presence of WN and WC characteristics, strong lines of Hydrogen, as well as the ratios and widths of the emission lines. The main differences are that HD~45166 presents a spectrum of lower ionization: in spite of having N v, it doesn’t present O vi. Besides, the star also presents He i emission lines, differently from the V Sge or CBSS objects. On the other hand, unlike HD~45166, no V Sge star or CBSS has the spectrum of the secondary star published up to now.
In the next section we will describe our observations. In Sects. 3 and 4 we describe the optical spectrum of HD 45166 and the spectral classification of the secondary star. In Sect. 5 we derive the orbital period while in Sect. 6, two additional periods are proposed. In Sect. 7, the masses and orbital inclination are derived. In Sects. 8 and 9 we present a discussion and conclusions. In an accompanying paper (Paper II) we study the structure of the wind and discuss the nature of the system.
## 2 Observations
We made spectroscopic observations of HD~45166 from 1998 to 2004 (Table 2), using the Coudé spectrograph at the 1.6 m telescope of Laboratório Nacional de Astrofísica (LNA) in Itajubá, Brazil, and the Fiber-fed Extended Range Optical Spectrograph (FEROS) (Kaufer et al. kauf (1999)) at the 1.52 m telescope of the European Southern Observatory (ESO) in La Silla, Chile.
At the Coudé spectrograph we employed 600 l mm<sup>-1</sup> and 1800 l mm<sup>-1</sup> gratings, resulting in spectral resolution of 0.7 and 0.2 Å FWHM and reaching S/N of about 20 to 30 at the continuum. We used a retro-illuminated Site ($`1024\times 1024`$) CCD detector with 24 micrometers resolution elements. Several exposures of bias and flatfield were obtained to correct for the sensitivity of the CCD. Measurements of dark current were not necessary. The slit of the spectrograph was adjusted to 250 $`\mu `$m (about 1.1 arcsec). Observations of Thorium lamps were made for the purpose of wavelength calibration. The data reduction was performed with the standard procedures, using IRAF <sup>1</sup><sup>1</sup>1IRAF is distributed by the National Optical Astronomy Observatories, which are operated by the Association of Universities for Research in Astronomy, Inc., under cooperative agreement with the National Science Foundation. routines. Typical calibration $`rms`$ residuals were of 2 mÅ for the 1800 l mm<sup>-1</sup> grating observations.
The FEROS spectrograph, on the other hand, uses a bench mounted echelle grating with reception fibers in the Cassegrain focus. Its measured spectral resolving power is R = 48 000. A completely automatic online reduction system was available (Stahl et al. stahl (1999)) and was adopted by us. We obtained a total of 40 spectra with integration times of 15 minutes and readout time of about 7 minutes. The signal-to-noise ratio at the continuum for each individual spectrum was typically S/N=80 at the central wavelength region (4500 through 7000 Å), decreasing to about S/N=35 at the blue and red edges of the spectrum. The wavelength calibration $`rms`$ residuals were of 7 mÅ. We cut the spectra in slices of 250 Å to normalize each slice by interactively fitting a low order Legendre polynomial to the continuum.
## 3 The optical spectrum
The average of the 40 spectra obtained with the FEROS spectrograph is presented in Fig. 1 in logarithmic scale, in order to show a maximum amount of details to the reader. It presents a large set of lines in absorption as well as in emission (please note the large number of telluric absorption lines). Quantitative characteristics of the lines are presented in Tables 3, 4 and 5. The spectrum in absorption presents lines of H, He i, C ii, N i, O i, Mg ii, Si ii and Fe ii. Weak lines in absorption of C i, Mg i, Al ii, Si i, P ii, S ii, Cr ii, Fe i, and Ni ii are also present. The Fe ii lines present an average heliocentric velocity of $`V_r=+6.3\pm 0.2`$ km s<sup>-1</sup>, where the uncertainty is the standard deviation divided by the square root of the number of measurements. This $`0.2`$ km s<sup>-1</sup> uncertainty is 30 times better than the spectral resolution. The lines of Si ii have $`V_r=+7.3\pm 0.6`$ km s<sup>-1</sup> and the lines of O i, $`V_r=+5\pm 1`$ km s<sup>-1</sup>. He i lines present $`V_r=+12\pm 3`$ km s<sup>-1</sup>. The discrepancy between the radial velocities of He i lines and the other lines can be explained by the emission that can be seen in the blue wings of some of the He i lines as in the case of He i 4713 Å.
A complete list of absorption lines is given in Table 3 for all of the absorption lines with $`W_\lambda >0.02`$ Å. Besides the measured wavelength we supply, also, the rest wavelength (in parenthesis), the radial velocity in the heliocentric system, the equivalent width, the FWHM and the identification of the species <sup>2</sup><sup>2</sup>2The line identifications were done with the help of the following online databases: the ”NIST Atomic Spectra Database” (http://www.physics.nist.gov/cgi-bin/AtData/main\_asd), the ”Atomic Line List” (http://www.pa.uky.edu/~peter/atomic/) and the ”Atomic Molecular and Optical Database Systems” (http://amods.kaeri.re.kr/spect/SPECT.html).. It is worth to notice that isolated lines have radial velocities and widths that are comparable amongst themselves. Doublets, triplets and blended lines are easily identified by the discrepancies in the values of these measurements. Selected weak absorption lines are also shown in Table 4, where we present only species which have not been listed in Table 3.
The emission spectrum is also quite rich, presenting lines of H, He i/ii, C iii/iv, N iii/iv/v, O iii and Si iv. In addition, it shows lines that possibly can be identified with O ii and Si iii, as well as some unidentified lines. The H lines are blended with lines of He ii, and this affects the measurement of their widths and also their intensities. In Table 5 we show a list as complete as possible of all of the emission lines. Besides the measured and rest wavelengths, we supplied the radial velocity in the heliocentric system, the equivalent width, FWHM and the identification of the species.
An important aspect of the emission line spectrum is the diversity of the widths and line profiles. The FWHM varies from 70 km s<sup>-1</sup> for the weakest lines up to 370 km s<sup>-1</sup> for the most intense ones. The H and He lines are systematically broader than the CNO lines and have a profile that can be described approximately by a profile of Voigt/Gauss, while the CNO lines are quite well adjusted by a Lorentz profile (see Fig. 2).
The emission line spectra observed in 1998, 1999, 2002 and 2004 are qualitatively quite similar. However, the equivalent width of the He ii 4686 Å line decreased from $`21.0\pm 0.5`$ Å in 1998 and $`23.3\pm 0.4`$ Å in 1999 to $`15.4\pm 0.4`$ Å in 2002, and increased again to $`22.3\pm 0.5`$ Å in 2004. Still, the FWHM remained constant along these observations, conserving the value of 380 $`\pm `$ 10 km s<sup>-1</sup>.
## 4 The spectral classification of the secondary star
The absorption spectrum of the secondary star was detected for the first time by Hiltner (hilt1 (1956)) who described the spectrum of HD 45166 as of ”high excitation superposed on a O9 type star”. Heap & Aller (heap (1978)) classified the secondary as B8 V and this classification has been adopted in the literature until today. With the high resolution and the high signal-to-noise ratio that characterize our data, it seems to be adequate to return to the subject and to verify which, in fact, is the best classification for the star. Many of the conclusions, such as the values of the masses, depend critically on this parameter.
In order to re-evaluate the spectral class of the secondary component of the binary system, we compared the ratio of the observed Mg i 5183.6 Å/Mg ii 4481.3 Å equivalent widths (W<sub>λ</sub>(Mg i)/W<sub>λ</sub>(Mg ii) = $`0.0464\pm 0.0038`$) to the ones given by the Kurucz (kurucz (1994)) models as provided by the VALD (Vienna Atomic Line Database) service (Kupka et al. kupka (1999)). We also compared this measured ratio to the ones found at the high resolution and high S/N ratio spectra of several B type stars from the UVES Paranal Observatory Project (Bagnulo et al. bagnu (2003)). From both comparisons we derived a B7$`\pm `$1 V spectral type. This classification is similar, within the errors, to that (B8 V) obtained by Heap & Aller (heap (1978)). According to Cox (cox (1999)), the properties of a B7V star are T$`{}_{\mathrm{eff}}{}^{}=\mathrm{13\; 500}`$K, M = 4.8 M, R = 3.4 R and $`M_B=0.7`$.
## 5 The orbital period
Heap & Aller heap (1978) determined that HD 45166 is a binary system composed of a ”quasi Wolf-Rayet” star and a B8 V star. With the observations described in the previous section, we may now verify whether this star has or not a detectable orbital period. If our data present radial velocity variations on both the emission and absorption features than we may be dealing with a double spectroscopic binary. However, previous observations discovered no radial velocity variability in excess of 10 km s<sup>-1</sup> (WHSH89). Therefore it is very important that any measurement shall be made with maximum accuracy.
Photospheric absorption lines of a star are excellent tracers of its movement and potential indicators of an orbital period. To reach the necessary accuracy, we measured the radial velocities of the O i 7772/4/5 and 8446 Å absorption lines from the FEROS spectra and calculated their weighted average, taking as weight the equivalent widths, respectively W$`{}_{\lambda }{}^{}=0.115`$ Å (7772 Å), W$`{}_{\lambda }{}^{}=0.133`$ Å (7774 Å), W$`{}_{\lambda }{}^{}=0.120`$ Å (7775 Å) and W$`{}_{\lambda }{}^{}=0.163`$ Å (8446 Å). To proceed, we corrected the O i average radial velocities for possible uncertainties in the wavelength calibration. This correction was based on the averaged residuals of the nearby 7670.26 Å and 7671.31 Å telluric lines observed wavelengths, as compared to their rest wavelengths. The use of telluric absorption lines as stationary comparison source was analyzed by Griffin & Griffin griff (1973). In our case, to estimate the uncertainties resulting from this correction, one should note that the radial velocities of these two telluric lines differ one from the other by 0.2 km s<sup>-1</sup> for a same exposure and by 0.5 km s<sup>-1</sup> for different observations from the same night. Therefore, the uncertainty in the measurements of the O i lines position relative to the nearby telluric features may also be 0.2 km s<sup>-1</sup>, in spite of the observed spectral line resolution in this region being 6.2 km s<sup>-1</sup>. The relative dispersion of the radial velocities of O I in Fig. 4 (top diagram) suggests that this is quite reasonable. One could expect that measurements of absorption lines other than those O i would have uncertainties of 0.5 km s<sup>-1</sup>, because of the absence of such correction by close telluric features.
The CLEAN (Roberts et al. roberts (1987)) and PDM (Phase Dispersion Minimization, Stellingwerf stellin (1978)) routines for period search applied to the O i average radial velocities indicate a period of P = 1.596 $`\pm `$ 0.003 days (Fig. 3). The associated radial velocity curve (Fig. 4, top) displays an amplitude of K$`{}_{2}{}^{}=2.4\pm `$ 0.2 km s<sup>-1</sup> and an eccentricity of e = 0.18 $`\pm `$ 0.08. These parameters were obtained fitting the data with the Russell-Wilsing (Binnendijk binn (1960)) method.
The mass function is
$$f(M_1)=\frac{M_1^3sin^3i}{(M_1+M_2)^2}=(1.04\times 10^7)(1e^2)^{3/2}K_2^3P(M_{\mathrm{}})$$
(1)
where K<sub>2</sub> is measured in km s<sup>-1</sup> and P in days. With the values derived above we obtain a mass function of $`f(M_1)=2.2\times 10^6`$ M.
## 6 The 5 and 15 hours periods
Since the emission spectrum (associated to the primary component of the binary system) is quite rich, we measured radial displacements of these lines to verify if the orbital period determined above is also seen in this component. In a first attempt (Oliveira & Steiner oliv (2002)) analyzing only the emission lines from the data obtained at LNA in 1999, we found a candidate period of 0.357 days.
To measure the radial velocities of the emission lines from the FEROS spectra we used the C iv 5801/11 Å lines, because these features are narrower than the H and He lines and they are very well fitted by Lorentzian profiles (Fig. 2, top diagram). Besides, they are located only 100 Å away from the Na i 5890/96 Å interstellar absorption lines (Fig. 5), which can be used as fiducial marks to correct for radial velocity inaccuracies. This is the same procedure we applied to the O i absorption lines above, when the fiducial references were telluric features. Therefore, the C iv radial velocity measurements were performed by simultaneously fitting Lorentz profiles and obtaining the simple average of both radial velocities. Then we subtracted, from this average, the residuals of the radial velocities of the Na i interstellar lines relative to their median value.
The resultant radial velocity curve shows a period of 0.205 days (5 hours) (Fig. 6, top). The fact that two spectroscopic periods exist is a complicating factor because, with the small amplitudes involved, it is very difficult to determine them. Another factor that complicates the determination of the periods is that their values coincide with unit fractions of a day. We have confused (Oliveira & Steiner oliv (2002)) the period of 0.205 days ($`1/5`$ of a day) with 0.357 days ($`1/3`$ of a day). In the periodogram obtained from the LNA data, the most visible period is, in fact, $`1/3`$ of a day. One should mention that the orbital period is $`3/2`$ of a day.
At this time we should wonder which, after all, is the orbital period as we have now identified two periods. We believe that the 1.596 days period is orbital for two strong reasons. First, it was derived using the photospheric absorption lines from the secondary star. Second, the 5 hours period is so short that, if it were orbital, the secondary star would not fit in its Roche lobe.
If the 1.596 $`\pm `$ 0.003 days period is the orbital one, we should find this period also in the primary star, that is, we should search for this period in the emission lines. To do this, we subtracted the average radial velocity curve, folded with the 5 hours period, from the raw data of C iv radial velocity. The result is plotted as a function of the 1.596 days period (Fig. 4, bottom) and shows a curve symmetric to the absorption radial velocity curve with the same period, but with an amplitude slightly larger, that is, K$`{}_{1}{}^{}=2.7\pm `$ 0.2 km s<sup>-1</sup>. The radial velocity curve derived from the C iv emission lines is in anti-phase with the O i absorption lines curve. This should be expected if the emission and absorption lines are produced on distinct stellar components of the binary system and if the sample period is the orbital one. To investigate whether the C iv radial velocity variations could be attributed to superposed photospheric absorption features from the secondary component, we searched for C iv 5801/11 Å absorption lines in high resolution spectra of B stars. Our FEROS spectra of HD 103401 (B5 V), HD 104432 and HD 102465 (both B9 V), taken with the same instrumental configuration as the spectra of HD 45166, do not show any features above the noise that could be associated to C iv 5801/11 Å lines. This fact, together with the anti-phase behavior of the radial velocity curves, lead to the interpretation of the C iv radial velocity variations as produced by the orbital motion of the primary star.
The mass function of the secondary component is
$$f(M_2)=3.1\times 10^6M_{\mathrm{}}$$
(2)
After we subtract the 1.596 days and the 5 hours modulations, we applied the PDM routine on the residuals and we identified an additional period of 0.64 days (15 hours, Fig. 6, bottom). Using these techniques, no other candidate period appeared with amplitude greater than 0.5 km s<sup>-1</sup>. In Paper II we will return to the non-orbital periods and their interpretation.
## 7 Masses and inclination
The orbital period was detected in both stars (associated to the absorption and also emission lines), so this system is a double spectroscopic binary. With the radial velocity amplitudes from both stars, we obtain
$$q=\frac{M_2}{M_1}=\frac{K_1}{K_2}=1.13\pm 0.11$$
(3)
As a main sequence B7 V star, the mass of the secondary is M$`{}_{2}{}^{}=4.8\pm 0.5`$ M; therefore M<sub>1</sub> = 4.2 $`\pm `$ 0.7 M. The mass of the primary star was estimated by WS83. They assume that the primary star obeys the wind relations of typical OB stars. From an approximate wind model and observed P-Cygni profile parameters from IUE, they estimate a mass of 0.5 M. In that case the mass was not inferred from dynamical parameters such as radial velocity curves.
The inclination angle of the system, obtained from the mass function, is i = 0.77$`\mathrm{°}`$ $`\pm `$ 0.09$`\mathrm{°}`$. This is probably one of the smallest orbital inclination angles known for a binary system.
The major semi-axes of the orbits, (a<sub>1</sub>,a<sub>2</sub>), are defined as
$$a_{1,2}sini=(1.98\times 10^2)(1e^2)^{1/2}K_{1,2}PR_{\mathrm{}}$$
(4)
and, in the case of HD~45166, are a$`{}_{1}{}^{}=6.3\pm `$ 0.6 R and a$`{}_{2}{}^{}=5.5\pm `$ 0.5 R.
Would it be the case to ask whether the secondary star fills its Roche lobe? As the eccentricity of the orbit is non zero, the Roche lobes increase and decrease as a function of the orbital phase. At periastron, the Roche lobes have minimum dimensions and, at that phase, the effective radius of the secondary’s lobe is given by (Hilditch hildi (2001))
$$R_{r2}(min)=\frac{4.21P^{2/3}(1e)(K_1+K_2)(0.38+0.20\mathrm{log}(M_2/M_1))M_1}{K_2(M_1+M_2)^{2/3}}R_{\mathrm{}}$$
(5)
For the values determined for HD~45166, we have $`R_{r2}(min)`$ = 3.7 $`\pm `$ 0.8 $`R_{\mathrm{}}`$. Considering the uncertainties, this value is very close to the secondary’s radius (see Sect. 4) and, therefore, at periastron the secondary star may fill its Roche lobe.
The maximum radius for the secondary’s Roche lobe is (Hilditch hildi (2001))
$$R_{r2}(max)=\frac{4.21P^{2/3}(1+e)(K_1+K_2)(0.38+0.20\mathrm{log}(M_2/M_1))M_1}{K_2(M_1+M_2)^{2/3}}R_{\mathrm{}}$$
(6)
that is, 5.4 $`\pm `$ 1.1 R. In Table 6 we list the basic parameters for the system.
## 8 Discussion
### 8.1 Gravitational redshift
The General Theory of Relativity predicts the existence of gravitational redshift given by (for $`z=v/c1`$, where $`z`$ is the redshift, $`v`$ is velocity and $`c`$ is the velocity of light)
$$z=\frac{V_{\mathrm{gr}}}{c}=\frac{GM}{Rc^2}=\frac{2.11\times 10^6M/M_{\mathrm{}}}{R/R_{\mathrm{}}}$$
(7)
At the surface of the secondary star this corresponds to $`V_{\mathrm{gr}}`$ = +0.9 km s<sup>-1</sup> and for the primary star, $`V_{\mathrm{gr}}`$ = +2.0 km s<sup>-1</sup> (where we considered $`R_1=1.3R_{\mathrm{}}`$, see Paper II). As these values are larger than the precision of our measurements, they need to be taken into account in the data analysis. If the average velocity of the secondary absorption spectrum is +6.2 $`\pm `$ 0.3 km s<sup>-1</sup>, then the expected average photospheric velocity in our spectra of the primary star should be of +7.3 km s<sup>-1</sup>. The average heliocentric velocity of the secondary star, corrected for the gravitational redshift, is $`<v(HC)>`$ = +5.3 $`\pm `$ 0.3 km s<sup>-1</sup>. This velocity is not the radial velocity of the center of mass, because this depends on the sampling of the observation phases. The value of $`\gamma `$, that is, of the radial velocity of the mass center, is determined from the radial velocity curve and corresponds to the value that divides the curve in two segments of same underlying areas. Its value, obtained from the secondary star’s absorption lines radial velocities and corrected for the gravitational redshift, is $`\gamma `$= 4.5 $`\pm `$ 0.2 km s<sup>-1</sup> (Table 6).
### 8.2 The distance to the system
Heap & Aller (heap (1978)), as well as WS83 determined the distance to the system based on the secondary star spectroscopic parallax. The first authors found a distance of 1.26 kpc. They also determined that the neighboring star HD~44498 has a reddening of $`E(BV)=0.14`$. WS83 determined a distance of 1.208 kpc. The color excess measured by them is $`E(BV)=0.15`$ and it is quite reliable, given that it is based on the interstellar extinction band at 2200 Å.
In order to obtain an independent estimate of the interstellar extinction, we measured the equivalent width of the diffuse interstellar band (DIB) at 5780 Å and found W<sub>λ</sub>(5780 Å)$`=100\pm 6`$ mÅ. Using the relation of the DIB equivalent width versus $`E(BV)`$ color excess (Somerville some (1988))
$$E(BV)=[W_\lambda (5780\AA )+46]/940$$
(8)
where W<sub>λ</sub> is given in mÅ, we obtain, $`E(BV)=0.155\pm 0.007`$. The estimated error does not take into account the uncertainty of Somerville’s calibration. The two methods of estimating the color excess agree quite well. One should notice, however, that the DIB at 5780 Å may be contaminated by the N iv 5776.3 Å emission line.
We measured the equivalent width of Mg ii 4481 Å from the FEROS spectra of two comparison stars and obtained for HD~103401 (B5 V), W<sub>λ</sub>(4481 Å) = 0.263 $`\pm `$ 0.005 Å and for HD~102465 (B9 V), W<sub>λ</sub>(4481 Å) = 0.377 $`\pm `$ 0.004 Å. Interpolating the spectral type one would expect, for the B7$`\pm `$1 V secondary component of HD~45166, W<sub>λ</sub>(4481 Å) = 0.320 $`\pm `$ 0.005 Å. We found 0.177 $`\pm `$ 0.005 Å, instead, due to the dilution from the companion’s continuum. This means that the secondary star contributes 55% $`\pm `$ 2% in the $`B`$ filter, and the distance to HD~45166 is, therefore, d = 1.3 $`\pm `$ 0.2 kpc.
The radial velocity of the interstellar Na i lines determined by us (see Table 3) is 22.8 km s<sup>-1</sup> in the heliocentric system. After correcting this value to the Local Standard of Rest (LSR) we obtain a radial velocity of V$`{}_{LSR}{}^{}=+7.1`$ km s<sup>-1</sup>. From the galactic rotation curve we obtain, for the coordinates of HD~45166 (l = 203.1$`\mathrm{°}`$ / b = -2.3$`\mathrm{°}`$), the expected radial velocity of 10.8 (d/kpc) km s<sup>-1</sup>. The radial velocity of the average stellar spectrum, corrected for the gravitational redshift, is V$`{}_{LSR}{}^{}=9.0`$ km s<sup>-1</sup>. We conclude, therefore, that the star has an anomalous velocity of approximately $`23.1`$ km s<sup>-1</sup> with respect to LSR.
### 8.3 DACs and the photospheric connection
The coincidence between the orbital period ($`P=1.596\pm 0.003`$ days) and the recurrence time of $`1.60\pm 0.15`$ days (WHSH89) for the Discrete Absorption Components (DACs) is striking. Why should exist an association between the recurrent appearance of slowly accelerating clouds and the orbital period? Somehow the clouds responsible for the periodic absorption seem to be ejected at every orbit. What could be the periodic event, responsible for this ejection? At a first glance one could imagine that the non-zero eccentricity may play a key role in this game. At every orbit there is a periastron event in which the secondary fills its Roche lobe and is capable of transferring matter to the primary. This periodic accretion may be physically related to the DACs but how, exactly, is not clear at this point. We will revisit this issue in Paper II.
## 9 Conclusions
In what follows we present the main conclusions of this paper:
1. The optical spectrum of HD~45166 presents a great wealth of information both in the emission as well as in the absorption spectrum. The emission spectrum has lines of H, He i/ii, C iii/iv, N iii/iv/v, O ii/iii/iv, Si iii/iv. The spectrum in absorption has lines of H, He i, C ii, N i, O i, Si ii, Mg ii and Fe ii. Weak lines of C i, Mg i, Si i, P ii, S ii, Cr ii and Fe i are also seen.
2. We classified the spectral type of the secondary star as B7 V and, therefore, it should have a mass of M<sub>2</sub> = 4.8 M and a radius of R<sub>2</sub> = 3.4 R.
3. The emission lines have great diversity of widths and profiles. The most intense lines have FWHM of 370 km s<sup>-1</sup> and the weakest lines of 70 km s<sup>-1</sup>. Lines of H and He have Voigt/Gauss profiles and are systematically broader than the lines of CNO, that have Lorentz profiles.
4. HD~45166 is shown to be a double spectroscopic binary with orbital period of 1.596 $`\pm `$ 0.003 days and eccentricity of e = 0.18 $`\pm `$ 0.08.
5. Standard techniques for period search applied to the emission lines show two additional periods, of 5 and 15 hours.
6. The amplitude of the radial velocities of the orbital period are K<sub>1</sub> = 2.7 $`\pm `$ 0.2 km s<sup>-1</sup> and K<sub>2</sub> = 2.4 $`\pm `$ 0.2 km s<sup>-1</sup>. We derived M<sub>1</sub> = 4.2 $`\pm `$ 0.7 M and i = 0.77$`\mathrm{°}`$ $`\pm `$ 0.09$`\mathrm{°}`$.
7. The secondary star’s radius may be about the size of its Roche lobe at periastron. Therefore we can consider that, at periastron, it fills or it is very close to fill its Roche lobe.
8. The expected gravitational redshift is 0.9 km s<sup>-1</sup> for the secondary star and 2.0 km s<sup>-1</sup> for the primary. These values are larger than the precision of our measurements and they need to be taken into account in the analysis of the data.
9. We estimated a color excess of $`E(BV)=0.155\pm 0.007`$ and the distance to HD~45166 as 1.3 $`\pm `$ 0.2 kpc.
10. We suggest that the discrete absorption components (DACs) observed in the ultraviolet with a periodicity similar to the orbital period may be induced by the periastron events.
###### Acknowledgements.
We would like to thank Dr. Albert Bruch for his careful reading of the manuscript and the anonymous referee for the constructive comments. We are specially grateful to Dr. Luiz Paulo Vaz for kindly fitting the measurements of the radial velocity with his program of orbit adjustment with non-zero eccentricity. |
warning/0507/math0507257.html | ar5iv | text | # Likely path to extinction in simple branching models with large initial population
## 1. Introduction and main results
In population genetics it is often important to look back at the development of populations. In this paper we consider the question of how extinctions occur, and in particular, what path a population takes on the road to extinction. Using asymptotic analysis when initial population values are large, we are able to find most likely path to extinction as well as the extinction probability in two simple branching models in discrete and continuous time. In both examples we use the large deviation principle (LDP) which is non-standard since random processes are nonnegative, and we use trajectories ending up at zero.
One of the contributions of this paper is in rigorous proofs of the LDP for processes on half space. It may appear to the reader that the LDP follows from known results in Markov chains and diffusions. This is only partly correct. The standard proof of the lower bound in the local LDP relies on the change of measure. This requires a certain point (the point where maximum in the Fenchel-Legendre transform is achieved) to be finite. In our case this point is at infinity, breaking down the standard approach. We therefore give complete proofs of LDP’s in Sections 4 (discrete time) and 5 (continuous time) following the scheme of Puhalskii . His approach states that the LDP is equivalent to exponential tightness plus local LDP, and is based on the stochastic exponential method (rather than the Laplace transform). Although we follow the scheme of Puhalskii we do not use idempotent probability and give direct proofs. Since these proofs are more technical, we placed them at the end, after results on extinction. Once the LDP is established, the problem of finding most likely path to extinction is in effect the problem of minimization of the rate function. This is typically a difficult problem due to nonlinearity. We are able to solve it by setting up the Bellman equation in discrete case, Section 2, and a dynamical control problem in continuous case, Section 3.
### 1.1. Galton-Watson process.
A prototype of a branching model in discrete time is the Galton-Watson process, described as follows.
Let $`X_n`$ denote the population size at time $`n`$, and $`\xi _{n+1}^j`$ the number of offspring of the $`j`$th individual. For each $`n=1,2,\mathrm{}`$, $`\{(\xi _n^j)_{j1}\}`$ is the sequence of independent identically distributed integer-valued random variables with the probability distribution function $`𝖯(\xi _n^j=\mathrm{})=p_{\mathrm{}},\mathrm{}=0,1,\mathrm{}.`$ The population size at time $`n+1`$ is given by
$$X_{n+1}=\underset{j=1}{\overset{X_n}{}}\xi _{n+1}^j,$$
where $`X_0=K>0`$. The state $`\{0\}`$ is absorbing, and the branching process $`(X_n)_{n0}`$ might be absorbed in $`\{0\}`$ at the extinction time
$$\tau =inf\{n:X_n=0\}.$$
If $`p_0=0`$, the population does not become extinct. However if $`p_0>0`$, it is well known (e.g. , ) that the extinction time $`\tau `$ is finite with probability one if and only if the offspring mean $`𝔪=_\mathrm{}1\mathrm{}p_{\mathrm{}}`$ does not exceed one ($`𝔪1`$). Moreover, for any $`𝔪`$, the distribution function of $`\tau `$ is computed using the offspring probability generating function $`𝔣(s)=_\mathrm{}0p_{\mathrm{}}s^{\mathrm{}},`$ $`0s1`$: for any $`N1`$,
$$𝖯(\tau N)=(𝔣_N(0))^K,$$
(1.1)
where $`𝔣_n(s)`$ is the $`n`$-th iterate of $`𝔣(s)`$, i.e. $`𝔣_n(s)=𝔣(𝔣_{n1}(s))`$ with $`𝔣_1(0)=𝔣(0)=p_0`$.
A natural question is how to find the “path to extinction” given that extinction occurred at time $`N`$, $`\tau =N.`$ The conditional distribution of the chain conditioned on extinction: for $`n=1,\mathrm{},N1`$,
$$\pi _{n|N}(i):=𝖯(X_n=i|\tau =N),i=1,2,\mathrm{}$$
gives the complete description. It can be used to find the conditional median or the traditional optimal in the mean square sense estimate $`\widehat{X}_n=_{i=1}^{\mathrm{}}i\pi _{n|N}(i)`$. Unfortunately such computations are involved, even using the Markov property of $`(X_n)`$. However, for large values of $`X_0=K`$, one path has an overwhelmingly large probability compared to the rest. Consider the normed branching process
$$x_n^K=\frac{X_n}{K}.$$
The limit in probability $`𝖯\text{-}lim_K\mathrm{}x_n^K=\widehat{x}_n`$ exists (see , ) and satisfies $`\widehat{x}_{n+1}=𝔪\widehat{x}_n,\widehat{x}_0=1.`$ The process $`\widehat{x}_n`$ is always positive, irrespective of the value of $`𝔪`$, so that, the approximation $`\widehat{x}_n`$ is inadequate for study of extinction, the fact is already mentioned in . In the approach we take, $`(x_n^K)_{nN}`$ is approximated on the set $`\{\tau N\}`$ by a deterministic sequence $`u_{}^{}:=(u_n^{})_{nN}`$ with $`u_0^{}=1`$, positive $`u_n^{}`$’s and $`u_N^{}=0`$, such that for small $`\delta >0`$ and large $`K`$,
$$𝖯\left(\underset{n=1}{\overset{N}{}}|x_n^Ku_n^{}|\delta \right)𝖯\left(\tau N\right).$$
This choice of $`u_{}^{}`$ might be warranted by the following argument. Since $`𝔣_n(0)`$ increases in $`n`$, for large $`K`$, $`(𝔣_N(0))^K`$ is considerably larger than any of $`(𝔣_n(0))^K`$ for $`n<N`$. Then, by (1.1), $`𝖯(\tau N)=𝖯(\tau =N)+𝖯(\tau N1)𝖯(\tau =N).`$ Consequently, for any $`u_{}=(u_n)_{nN}`$ with $`u_0=1`$ and $`u_n0`$,
$$𝖯\left(\underset{n=1}{\overset{N}{}}|x_n^Ku_n|\delta \right)𝖯\left(\tau N\right).$$
For large $`K`$, extinction for the process $`x_n^K`$ is a rare event, since the limit process $`\widehat{x}_n`$ is positive. Therefore, as in , we approach the problem of extinction using the large deviations theory, obtaining a new result as well as recover an asymptotic version of the well-known result (1.1) by using this theory. According to LDP, Theorem 4.1 and by analogy with the maximal likelihood estimator, the path $`(u_n^{})_{nN}`$ is said to be the most likely path to extinction of the normed population $`x_n^K`$.
Clearly, $`\tau `$ is the extinction time for both processes $`X_n`$ and $`x_n^K`$, so that, $`Ku_n^{}`$ (with large $`K`$) sets the pattern for the extinction path in the original branching process.
Figure below demonstrates likely paths to extinction for a binary splitting model with different parameters, $`p=p_0`$, illustrating the general result.
For formulating the main result, we use the log moment generating function, assuming its existence up to some $`t_0>0`$,
$$𝔤(t)=\mathrm{log}\underset{\mathrm{}0}{}e^t\mathrm{}p_{\mathrm{}},t(\mathrm{},t_0).$$
(1.2)
It is related to the moment generating function by
$$\mathrm{log}𝔣_n(0)𝔤_n(\mathrm{})\text{(Lemma }\text{2.1}\text{)}.$$
###### Theorem 1.1.
Assume $`p_0>0`$ and (1.2). Then, for any $`N1`$,
(i)
$`(u_n^{})_{nN}`$ $`=\underset{\begin{array}{c}u_0=1,u_N=0\\ u_n>0,nN1\end{array}}{argmax}\underset{\delta 0}{lim}\underset{K\mathrm{}}{lim}{\displaystyle \frac{1}{K}}\mathrm{log}𝖯\left({\displaystyle \underset{n=1}{\overset{N}{}}}|x_n^Ku_n|\delta \right)`$
with
$$u_n^{}=\underset{1in}{}𝔤^{}(𝔤_{Ni}(\mathrm{})),nN,$$
(1.3)
where $`𝔤_i(t)`$ is $`i`$-th iterate of $`𝔤(t)`$, $`𝔤_0(t)=t`$.
(ii)
$$\underset{\delta 0}{lim}\underset{K\mathrm{}}{lim}\frac{1}{K}\mathrm{log}𝖯(\underset{n=1}{\overset{N}{}}|x_n^Ku_n^{}|\delta )=\underset{K\mathrm{}}{lim}\frac{1}{K}\mathrm{log}𝖯(\tau n)).$$
### 1.2. Branching diffusion.
In continuous time, we consider the model of a branching diffusion $`X_t`$ defined by the Itô equation
$$dX_t=\alpha X_tdt+\sigma \sqrt{X_t}dB_t$$
(1.4)
with a positive initial condition $`X_0=K`$, where $`B_t`$ is a Brownian motion, $`\sigma ^2>0`$, and $`\alpha `$. Stochastic equation (1.4) possesses a strong nonnegative solution. Since the diffusion parameter degenerates, one way to see this is to construct the solution from the following approximating sequence $`(X_t^i)_{i1}`$:
$$X_t:=X_t^1I_{\{t\tau _1\}}+\underset{i1}{}X_{\tau _i}^iI_{\{\tau _i<t\tau _{i+1}\}},$$
where $`dX_t^i=\alpha X_t^idt+\sigma \sqrt{|X_t^i|i^1}dB_t,`$ $`X_0^i=K`$, and $`\tau _i=inf\{X_t^ii^1\}`$ the increasing sequence of stopping times $`(\tau _i)_{i1}`$ relative to the filtration generated by Brownian motion $`(B_t)`$ (see also Theorem 13.1, ). The strong uniqueness of (1.4) follows from Yamada-Watanabe’s theorem (see, e.g. Rogers and Williams, p. 265 ) since its drift and diffusion parameters are Lipschitz and Hölder (with coefficient $`\frac{1}{2}`$) continuous respectively.
Obviously,
$$\tau =inf\{t:X_t=0\}=\underset{i\mathrm{}}{lim}\tau _i.$$
We analyze the normed process $`x_t^K=\frac{X_t}{K}`$. Due to (1.4), $`x_t^K`$ solves the Itô equation
$$dx_t^K=\alpha x_t^Kdt+\frac{\sigma }{\sqrt{K}}\sqrt{x_t^K}dB_t,$$
(1.5)
with $`x_0^K=1`$. It can be readily shown that $`𝖯\text{-}lim_K\mathrm{}x_t^K=\widehat{x}_t`$ exists and solves $`{\displaystyle \frac{d\widehat{x}_t}{dt}}=\alpha \widehat{x}_t,\widehat{x}_0=1.`$ However, $`\widehat{x}_t`$ is always positive and is far from to be estimated path to extinction. As in the discrete time, in order to evaluate path to extinction for $`(x_t^K)_{tT}`$ for fixed $`T>0`$, we approximate $`(x_t^K)_{tT}`$ on the set $`\{\tau T\}`$ by a deterministic function $`(u_t^{})_{tT}`$ with $`u_0^{}=1,u_T^{}=0`$ and $`u_t^{}>0`$, such that for a small $`\delta >0`$ and large $`K`$,
$$𝖯\left(\underset{tT}{sup}|x_t^Ku_t^{}|\delta \right)𝖯\left(\underset{tT}{sup}|x_t^Ku_t|\delta \right)$$
for any $`(u_t)_{tT}`$ from the set $`\{u_0=1,(u_t>0)_{t<T},u_T=0\}`$.
Unfortunately, the useful formula of (1.1) type is not known to us in this case. Here we obtain its asymptotic version as $`K\mathrm{}`$, see (ii) below.
###### Theorem 1.2.
For any $`T>0`$,
1. $$(u_t^{})_{tT}=\underset{\begin{array}{c}u_0=1,u_T=0\\ u_t>0,t<T\end{array}}{argmax}\underset{\delta 0}{lim}\underset{K\mathrm{}}{lim}\frac{1}{K}\mathrm{log}𝖯\left(\underset{tT}{sup}|x_t^Ku_t|\delta \right)$$
is given by
$$u_t^{}=\{\begin{array}{cc}e^{\alpha t}\left(1\frac{1e^{\alpha t}}{1e^{\alpha T}}\right)^2,\hfill & \alpha 0\hfill \\ \left(1\frac{t}{T}\right)^2,\hfill & \alpha =0.\hfill \end{array}$$
(1.6)
2. $$\underset{K\mathrm{}}{lim}\frac{1}{K}\mathrm{log}𝖯(\tau T)=\underset{\delta 0}{lim}\underset{K\mathrm{}}{lim}\frac{1}{K}\mathrm{log}𝖯\left(\underset{tT}{sup}|x_t^Ku_t^{}|\delta \right)$$
$$=\{\begin{array}{cc}\frac{1}{\sigma ^2}\frac{\alpha }{1e^{\alpha T}},\hfill & \alpha 0\hfill \\ \frac{1}{\sigma ^2T},\hfill & \alpha =0\hfill \end{array}$$
###### Corollary 1.
1. $`u_{}^{}`$ has the remarkable property: it is the same for subcritical and supercritical case: $`u_t^{}(\alpha )u_t^{}(\alpha )`$.
2. For large $`K`$, the probability of extinction in $`[0,T]`$ is given by
$$𝖯(\tau T)\mathrm{exp}\left(\frac{K}{\sigma ^2}\frac{\alpha }{1e^{\alpha T}}\right).$$
In particular, for $`\alpha =0`$, $`𝖯(\tau T)e^{\frac{K}{2\sigma ^2T}}.`$
## 2. Proof of Theorem 1.1
We begin with
###### Lemma 2.1.
For any $`n1`$, $`𝔤_n(\mathrm{})=\mathrm{log}𝔣_n(0).`$
###### Proof.
The result follows by induction from the identity $`𝔤_n(\mathrm{log}t)\mathrm{log}𝔣_n(t)`$ for $`t(0,t_0).`$ Write
$$𝔤(\mathrm{log}(t))=\mathrm{log}\underset{\mathrm{}0}{}e^{\mathrm{}\mathrm{log}(t)}p_{\mathrm{}}=\mathrm{log}\underset{\mathrm{}0}{}e^{\mathrm{log}(t^{\mathrm{}})}p_{\mathrm{}}=\mathrm{log}\underset{\mathrm{}0}{}t^{\mathrm{}}p_{\mathrm{}}=\mathrm{log}𝔣(t).$$
If $`𝔤_{n1}(\mathrm{log}t)\mathrm{log}𝔣_{n1}(t)`$, then
$$𝔤_n(\mathrm{log}t)=𝔤(𝔤_{n1}(t)=𝔤(\mathrm{log}(𝔣_{n1}(t))=\mathrm{log}𝔣(𝔣_{n1}(t))=\mathrm{log}(𝔣_n(t)).$$
The proof of Theorem 1.1 is done in a number of steps.
(1) Recall that $`𝔤(t)`$ is convex function with $`𝔤(0)=0`$, $`𝔤(\mathrm{})=\mathrm{log}(p_0)`$ and $`𝔤^{}(t)>0,t>\mathrm{}`$ while $`𝔤^{}(\mathrm{})=lim_t\mathrm{}𝔤^{}(t)=0`$.
(2) By the local LDP (see, Theorem 4.1), for $`u_0=1,u_N=0`$ and other positive $`u_n`$’s, it holds
$$\underset{\delta 0}{lim}\underset{K\mathrm{}}{lim}\frac{1}{K}\mathrm{log}𝖯\left(\underset{nN}{}|x_n^Ku_n|\delta \right)=\underset{nN}{}I(u_n,u_{n1}).$$
(3) In order to find $`(u_n^{})_{nN}`$ such that for $`u_0=1,u_n>0,u_N=0`$
$$\underset{in}{}I(u_i,u_{i1})\underset{in}{}I(u_i^{},u_{i1}^{}),$$
(2.1)
we apply the Dynamic Programming.
Since $`u_N=0`$,
$$I(u_N,u_{N1})=\underset{t(\mathrm{},t_0)}{sup}(u_{N1}𝔤(t))=u_{N1}𝔤(\mathrm{})=:B_n(u_{N1})$$
(2.2)
is the boundary condition for the Bellman equation
$$B_n(u_{n1})=\underset{u>0}{inf}\left[B_{n+1}(u)+I(u,u_{n1})\right],1nN1.$$
(2.3)
For $`n=N1`$, we have
$$B_{N1}(u_{N2})=\underset{u>0}{inf}\left[u𝔤(\mathrm{})+\underset{t(\mathrm{},t_0)}{sup}\{tuu_{N2}𝔤(t)\}\right].$$
(2.4)
(2.4) provides the inequality,
$$B_{N1}(u_{N2})\underset{u>0}{inf}\left[u𝔤(\mathrm{})+tuu_{N2}𝔤(t)\right],t(\mathrm{},t_0)$$
which, with $`t=𝔤(\mathrm{})`$, is transformed into
$$B_{N1}(u_{N2})u_{N2}𝔤_2(\mathrm{}).$$
(2.5)
We show that the above inequality is equality. For $`u,u_{N2}>0`$, “$`sup_t`$” in (2.4) is attained at the point $`t^{}=t^{}(u,u_{N2})`$, so that, for any $`u>0`$,
$$B_{N1}(u_{N2})u\left[t^{}(u,u_{N2})𝔤(\mathrm{})\right]u_{N2}𝔤(t^{}(u,u_{N2})).$$
We choose $`u=u_{N1}^{}`$ such that $`t^{}(u_{N1}^{},u_{N2})=𝔤(\mathrm{})`$. This is possible since
$`𝔤(\mathrm{})=\mathrm{log}p_0,t^{}(0,u_{N2})=\mathrm{},`$
$`𝔤^{}(\mathrm{})=0,t^{}(𝔪,u_{n2})=0,𝔤^{}(0)=𝔪,`$
so that, the existence of $`u_{N1}^{}`$ follows from continuity, in $`u`$, of $`t^{}(u,u_{N2})`$.
The choice of $`u_{N1}^{}`$ gives the inequality
$$B_{N1}(u_{N2})u_{N2}𝔤(t^{}(u_{N1}^{},u_{N2}))=𝔤\left(𝔤(\mathrm{})\right)=𝔤_2(\mathrm{}).$$
Consequently, the opposite inequality for (2.5) holds true and, therefore,
$$B_{N1}(u_{N2})=u_{N2}𝔤_2(\mathrm{}).$$
It is obvious too that for any $`u_{N2}>0`$,
$$u_{N1}^{}=u_{N2}𝔤^{}(t^{}(u_{N1}^{},u_{N2}))=u_{N2}𝔤^{}(𝔤(\mathrm{})).$$
Further, by induction, we find the following pairs:
$`u_{N1}^{}`$ $`=𝔤^{}\left(𝔤(\mathrm{})\right)u_{N2}^{}`$
$`B_{N1}(u_{N2}^{})`$ $`=𝔤_2(\mathrm{})u_{N2}^{}`$
$`\mathrm{}`$
$`u_{N2}^{}`$ $`=𝔤^{}\left(𝔤_2(\mathrm{})\right)u_{N3}^{}`$
$`B_{N2}(u_{N3}^{})`$ $`=𝔤_3(\mathrm{})u_{N3}^{}`$
$`\mathrm{}`$
$`u_1^{}`$ $`=𝔤^{}\left(𝔤_{n1}(\mathrm{})\right)u_0`$
$`B_1(u_0)`$ $`=𝔤_n(\mathrm{})u_0(u_0=1).`$
With chosen $`(u_n^{})_{1nN1}`$, the Bellman equation (2.3) is transformed into the backward recurrent equation
$$B_n(u_{n1}^{})=B_{n+1}(u_n^{})+I(u_n^{},u_{n1}^{}),1nN1$$
with boundary condition $`u_{N1}^{}𝔤(\mathrm{})`$ (see, (2.2)).
Thus, $`B_1(1)=_{1nN}I(u_n^{},u_{n1}^{}).`$
On the other hand, the Bellman equation also yields
$$B_1(1)\underset{1nN1}{}I(u_n,u_{n1})+B_N(u_{N1})=\underset{1nN}{}I(u_n,u_{n1})$$
what proves (2.1).
(4) We recall that $`_{nN}I(u_n^{},u_{n1}^{})=𝔤_n(\mathrm{})`$, that is, by Lemma 2.1 and (1.1),
$$\underset{1nN}{}I(u_n^{},u_{n1}^{})=\mathrm{log}𝔣_N(0)=\frac{1}{K}\mathrm{log}𝖯(\tau N),K>0.$$
(5) Thus, (1)-(3) imply the statement (i); formula (1.3) follows from recurrence $`u_n^{}=𝔤^{}\left(𝔤_2(\mathrm{})\right)u_{n1}^{}`$, $`u_0^{}=1`$.
Finally (ii) follows from (4). ∎
## 3. Proof of Theorem 1.2
We apply the LDP Theorem 5.1. By the local LDP, with $`u_0=1`$, $`u_t>0`$ and $`u_T=0`$, we have
$$\underset{\delta 0}{lim}\underset{K\mathrm{}}{lim}\frac{1}{K}\mathrm{log}𝖯\left(\underset{tT}{sup}\left|x_t^Ku_t\right|\delta \right)=J_T(u),$$
where $`J_T(u)=\{\begin{array}{cc}\frac{1}{2\sigma ^2}_0^T\frac{(\dot{u}_tu_t)^2}{u_t}I_{\{u_t>0\}}𝑑t,\hfill & u_0=1,du_t=\dot{u}_tdt\hfill \\ \mathrm{},\hfill & \text{otherwise}.\hfill \end{array}`$
Therefore (i) is reduced to minimization of $`J_T(u)`$ in a class of absolutely continuous test functions $`u_t`$ with $`u_0=1`$, $`u_t>0`$ and $`u_T=0`$.
Set $`w_t={\displaystyle \frac{\dot{u}_tu_t}{\sqrt{u_t}}}`$, $`t[0,T)`$ and notice that the minimization of $`J_T(u_{})`$ is equivalent to the following control problem with the controlled process $`u_t`$, solving a differential equation
$$\dot{u}_t=\alpha u_t+\sqrt{u_t}w_t,t[0,T)$$
subject to $`u_0=1`$. The control action $`w_t`$ belongs to a class of measurable functions with $`_0^Tw_t^2𝑑t<\mathrm{}`$ bringing $`u_t`$ to zero at the time $`T`$. The control action $`w_t^{}`$ from this class is optimal if for any $`w_t`$,
$$_0^T(w_t^{})^2𝑑t_0^Tw_t^2𝑑t.$$
If $`w_t^{}`$ exists, then the controlled process $`u_t^{}`$ related to $`w_t^{}`$ minimizes $`J_T(u_{})`$ in the required class of continuous functions $`u_{}=(u_t)_{tT}`$.
In order to find $`w_t^{}`$, it is convenient to deal with (recall $`u_t0`$) $`v_t=\sqrt{u_t}`$ since $`v_t`$ solves the linear differential equation $`\dot{v}_t=\frac{\alpha }{2}v_t+\frac{1}{2}w_t,v_0=1.`$ If $`w_t^{}`$ exists, then $`w_t^{}`$ brings $`v_t`$ to zero at the time $`T`$, that is, $`0=v_T=e^{\frac{\alpha }{2}T}+_0^Te^{\frac{\alpha }{2}(Tt)}w_t^{}𝑑t`$ or, equivalently,
$$1=\frac{1}{2}_0^Te^{\frac{\alpha }{2}t}w_t^{}𝑑t.$$
(3.1)
Hence, by the Cauchy-Schwarz inequality $`1\frac{1}{2}_0^Te^{t\alpha }𝑑t_0^T(w_t^{})^2𝑑t`$, that is, the following lower bound holds: $`_0^T(w_t^{})^2𝑑t{\displaystyle \frac{2\alpha }{1e^{\alpha T}}}.`$ This lower bound is valid for any $`w_t`$ providing (3.1) , so that, the condition
$$_0^T(w_t^{})^2𝑑t=\frac{2\alpha }{1e^{\alpha T}}$$
is valid for $`w_t^{}=ce^{t\frac{\alpha }{2}}`$ for any constant $`c`$, bring $`w_t^{}=c^{}e^{t\frac{\alpha }{2}}`$ with $`c^{}`$ solving
$$1=_0^Te^{t\frac{\alpha }{2}}w_t^{}𝑑t=c^{}_0^Te^{t\alpha }𝑑t.$$
Hence,
$$c^{}=\{\begin{array}{cc}\frac{2\alpha }{1e^{T\alpha }},\hfill & \alpha 0\hfill \\ \frac{2}{T},\hfill & \alpha =0\hfill \end{array}\text{and}w_t^{}=\{\begin{array}{cc}\frac{2\alpha e^{t\frac{\alpha }{2}}}{1e^{T\alpha }},\hfill & \alpha 0\hfill \\ \frac{2}{T},\hfill & \alpha =0\hfill \end{array}$$
$$_0^T(w_t^{})^2𝑑t=\{\begin{array}{cc}\frac{2\alpha }{1e^{\alpha T}},\hfill & \alpha 0\hfill \\ \frac{2}{T},\hfill & \alpha =0.\hfill \end{array}$$
Finally, we find that
$`v_t^{}`$ $`=e^{t\frac{\alpha }{2}}{\displaystyle \frac{\alpha }{1e^{T\alpha }}}{\displaystyle _0^t}e^{(ts)\frac{\alpha }{2}}e^{s\frac{\alpha }{2}}𝑑s`$
$`=e^{t\frac{\alpha }{2}}\left[1{\displaystyle \frac{1e^{t\alpha }}{1e^{T\alpha }}}\right]=e^{t\frac{\alpha }{2}}\left({\displaystyle \frac{e^{t\alpha }e^{T\alpha }}{1e^{T\alpha }}}\right)`$
and, since $`u_t^{}=(v_t^{})^2`$, we obtain (1.6) and the proof of (i) is complete.
(ii) By (i),
$$J_T(u^{})=\frac{1}{\sigma ^2}\frac{\alpha }{1e^{\alpha T}}$$
(3.2)
We show that
$$\underset{K\mathrm{}}{lim}\frac{1}{K}\mathrm{log}𝖯\left(\tau T\right)=J_T(u^{}).$$
To this end, use the fact that $`\{\tau T\}=\{(\omega ,t):tT,x_t^K(\omega )=0\}`$. For notational convenience denote $`𝔄:=\{\tau T\}`$. Set $`𝔄^{\mathrm{𝖼𝗅}}`$ and $`𝔄^{\mathrm{𝗂𝗇𝗍}}`$ the closure and interior of $`𝔄`$. Then, by the LDP, we have
$`\overline{\mathrm{lim}}_K\mathrm{}{\displaystyle \frac{1}{K}}\mathrm{log}𝖯\left(𝔄^{\mathrm{𝖼𝗅}}\right)\underset{u:\{\begin{array}{c}u_s>0,s<t;\\ u_t=0\\ tT\end{array}}{inf}J_t(u)=\underset{tT}{inf}J_t(u^{})`$
$`\underset{¯}{\mathrm{lim}}_K\mathrm{}{\displaystyle \frac{1}{K}}\mathrm{log}𝖯\left(𝔄^{\mathrm{𝗂𝗇𝗍}}\right)\underset{u:\{\begin{array}{c}u_s>0,s<t;\\ u_t=0\\ tT\end{array}}{inf}J_t(u)=\underset{tT}{inf}J_t(u^{}).`$
Since $`\underset{¯}{\mathrm{lim}}_K\mathrm{}=\overline{\mathrm{lim}}_K\mathrm{}`$ implies the existence of $`lim_K\mathrm{}`$, it remains to show that $`inf_{tT}J_t(u^{})=J_T(u^{})`$.
Notice that (3.2) is valid with $`T`$ replaced by any $`t<T`$ with $`u_{}^{}`$ replaced by the corresponding $`u_{}^{,t}=\{u_0^{,t}=1;u_s^{,t}>0,s<t;u_t^{s,t}=0\}`$. In other words, for any $`t`$,
$$J_t(u^{,t})=\frac{1}{\sigma ^2}\frac{\alpha }{1e^{\alpha t}},$$
and $`J_t(u^{,t})`$ increases to $`J_T(u_{}^{})`$ with $`tT`$. ∎
## 4. LDP in Discrete Time
Let $`m=inf\{nN:u_n=0\}`$ and $`m=\mathrm{}`$ if all $`(u_n)_{nN}`$ are positive.
$$I(y,x)=\underset{t(\mathrm{},t_0)}{sup}[tyx𝔤(t)].$$
###### Theorem 4.1.
Assume (1.2). For any $`N1`$, the family $`\{(x_n^K)_{nN}\}_K\mathrm{}`$ obeys the LDP in $`_+^N`$, supplied by the Euclidian metric $`\varrho _N`$, with the speed $`\frac{1}{K}`$ and the rate function
$$J_N(u_{})=\{\begin{array}{ccc}\underset{n=1}{\overset{m1}{}}I(u_n,u_{n1})u_{m1}\mathrm{log}(p_0),\hfill & \begin{array}{c}u_0=1\\ u_n=0,n>m\end{array}\hfill & \\ \underset{n=1}{\overset{N}{}}I(u_n,u_{n1}),\hfill & \begin{array}{c}u_0=1\\ u_n>0,nN\end{array}.\hfill & \\ \mathrm{},\hfill & \begin{array}{c}n:u_n=0,u_{n+1}>0\\ \text{or}u_01\end{array}\hfill & \end{array}$$
###### Remark 1.
LDP for branching processes have been considered in the literature, see, for example, , , . However, they were concerned with the sequence $`\frac{X_n}{X_{n1}}`$, as $`n\mathrm{}`$, whereas here we consider the LDP for $`\frac{X_n}{X_0}`$ processes indexed by the large initial value.
###### Remark 2.
The nonnegativity of $`x_n^K`$ provides some difficulty for verification of LDP at the “point of extinction” where the test function becomes zero. For set $`𝕊`$ of test functions that keep away from zero the statement of the theorem is implied by a result in Klebaner and Zeitouni, and other known results that can be adapted to our setting (see, e.g. Kifer, , Puhalskii, , Klebaner and Liptser, , etc.). But $`\{\tau N\}𝕊`$ , and for the sake of completeness and accuracy we give the complete proof below, with a new proof of the lower bound in the local LDP.
### 4.1. Proof of Theorem 4.1
We follow standard (necessary and sufficient) conditions for proving the LDP by showing the exponential tightness:
$$\underset{C\mathrm{}}{lim}\overline{\mathrm{lim}}_K\mathrm{}\frac{1}{K}\mathrm{log}𝖯\left(\mathrm{\Omega }𝒦_C\right)=\mathrm{}$$
with compacts $`𝒦_C=\{\mathrm{max}_{1nN}x_NC\}`$, $`C\mathrm{}`$, and the local LDP:
$$\underset{\delta 0}{lim}\underset{K\mathrm{}}{lim}\frac{1}{K}\mathrm{log}𝖯\left(\varrho _N(x_{}^K,u_{})\delta \right)=J_N(u_.).$$
Notice that (1.2) implies the existence of a stochastic exponential, with $`t_nKt_0`$,
$$_{(t_1,\mathrm{},t_N)}^K(x_1^K,\mathrm{},x_{N1}^K)=\underset{n=1}{\overset{N}{}}𝖤\left(e^{t_nx_n^K}|_{n1}\right),$$
where $`(_n)_{n0}`$ is the filtration, with $`_0=\{\mathrm{},\mathrm{\Omega }\}`$, generated by $`(x_n^K)_{n1}`$.
Set
$$𝔷_n=e^{_{in}t_{\mathrm{}}x_i^K\mathrm{log}_{(t_1,\mathrm{},t_n)}^K(x_1^K,\mathrm{},x_{n1}^K)}.$$
(4.1)
The random process $`(𝔷_n,_n)_{nN}`$ is the (positive) martingale,
$$𝖤𝔷_N=1.$$
(4.2)
#### 4.1.1. Exponential tightness
Since $`\mathrm{max}_{1nN}x_i^K_{1nN}x_n^K`$, it is enough to show
$$\underset{C\mathrm{}}{lim}\overline{\mathrm{lim}}_K\mathrm{}\frac{1}{K}\mathrm{log}𝖯\left(\underset{1iN}{}x_i^KC\right)=\mathrm{}.$$
Set $`t^{}=argmax_{t(\mathrm{},t_0)}[t𝔤(t)]`$. Since $`𝔤(0)=0`$, we have that $`t^{}(0,t_0)`$ and $`𝔤(t^{})<t^{}`$. We choose $`t_nt^{}K(<Kt_0)`$, and introduce $`𝔄=\left\{_{1in}x_i^KC\right\}`$. With chosen $`t_n`$, we have $`𝖤𝔷_N=1`$ and, therefore, $`𝖤I_𝔄𝔷_N1`$. Taking into account this inequality and (4.1), write
$`1`$ $`𝖤I_𝔄e^{_{\{1nN\}}t^{}x_n^K\mathrm{log}_{(t^{},\mathrm{},t^{})}^K(x_1^K,\mathrm{},x_{N1}^K)}`$
$`=𝖤I_𝔄e^{Kt^{}_{\{1nN\}}x_n^KK𝔤(t^{})_{\{1nN\}}x_{n1}^K}`$
$`𝖤I_𝔄e^{K_{\{1nN\}}[t^{}𝔤(t^{})]x_n^KK|𝔤(t^{})|}`$
$`𝖤I_𝔄e^{KC[t^{}𝔤(t^{})]}=e^{KC[t^{}𝔤(t^{})]K|𝔤(t^{})|}𝖯\left(𝔄\right).`$
Therefore, $`\frac{1}{K}\mathrm{log}𝖯\left(𝔄\right)\underset{>0}{\underset{}{[t^{}𝔤(t^{})]}}C+|𝔤(t^{})|\underset{C\mathrm{}}{\overset{}{}}\mathrm{}.`$
#### 4.1.2. Local LDP. Upper bound
We may restrict ourselves by the test function $`u_{}=\{\underset{>0}{\underset{}{u_1,\mathrm{},u_{N1}}},\underset{=0}{\underset{}{u_N}}\}`$ and show that
$$\overline{\mathrm{lim}}_{\delta 0}\overline{\mathrm{lim}}_K\mathrm{}\frac{1}{K}\mathrm{log}𝖯\left(\varrho _N(x_{}^K,u_{})\delta \right)J_N(u_{}).$$
(4.3)
For the test function with all positive $`u_n`$’s and $`u_0=1`$ the proof of (4.3) is similar. For test function with $`u_n=0,u_{n+1}>0`$ or $`u_01`$, (4.3) is obvious. For others test functions the verification of (4.3) is reduced to the above-mentioned ones.
Let now $`𝔄=\left\{\rho _N(x_{}^K,u_{})\delta \right\}`$. By (4.2), we have
$$1𝖤I_𝔄𝔷_N=𝖤I_𝔄e^{_{\{1nN\}}[t_nx_i^KKx_{n1}^K𝔤(\frac{t_n}{K})]}.$$
(4.4)
Set $`t_n^{}=argmax_{t(\mathrm{},t_0)}[tu_nu_{n1}𝔤(t)]`$, $`nN1`$, and $`t_N^{}=l(l>0)`$, and take $`t_n=Kt_n^{}`$, then we derive from (4.4)
$`1`$ $`𝖤I_𝔄e^{K_{\{1nN\}}[t_n^{}u_nu_{n1}𝔤(t_n^{})]K_{1nN1}(t_n^{}+|𝔤(t_n^{})|)\delta }`$
$`=𝖤I_𝔄e^{K[_{\{1nN1\}}I(u_n,u_{n1})u_{N1}𝔤(l)]K_{1nN1}(|t_n^{}|+|𝔤(t_n^{})|)\delta }`$
$`=𝖤I_𝔄e^{K[J_{N1}(u_{})u_{N1}𝔤(l)]K_{1nN1}(|t_n^{}|+|𝔤(t_n^{})|)\delta }.`$
Hence, taking into account that $`lim_l\mathrm{}𝔤(l)=\mathrm{log}(p_0)`$, we obtain
$`{\displaystyle \frac{1}{K}}\mathrm{log}𝖯\left(𝔄\right)`$ $`[J_{N1}(u_{})+u_{N1}𝔤(l)]+{\displaystyle \underset{1iN1}{}}(|t_i^{}|+|𝔤(t_i^{})|)\delta `$
$`\underset{\delta 0}{\overset{}{}}[J_{N1}(u_{})+u_{N1}𝔤(l)]\underset{l\mathrm{}}{\overset{}{}}J_N(u_{}).`$
### 4.2. Local LDP. Lower bound
Obviously for $`u_{}`$ with $`J_N(u_{})=\mathrm{}`$, it is nothing to verify. Further as in the upper bound verification, we may restrict ourselves by the test function $`u_{}=\{\underset{>0}{\underset{}{u_1,\mathrm{},u_{N1}}},\underset{=0}{\underset{}{u_N}}\}`$ with $`𝖯(\xi _1^1=0)=p_0>0`$ and show that
$$\underset{¯}{\mathrm{lim}}_{\delta 0}\underset{¯}{\mathrm{lim}}_K\mathrm{}\frac{1}{K}\mathrm{log}𝖯\left(\varrho _N(x_{}^K,u_{})\delta \right)J_N(u_{}).$$
Write
$`\{\varrho _N(x_{}^k,u_{})\delta \}`$ $`=\left\{\varrho _{N1}(x_{}^k,u_{})+x_N^K\delta \right\}`$
$`\left\{\varrho _{N1}(x_{}^k,u_{})0.5\delta ,x_N^K0.5\delta \right\}`$
$`\left\{\varrho _{N1}(x_{}^k,u_{})0.5\delta ,x_N^K=0\right\}`$
$`\left\{\varrho _{N1}(x_{}^k,u_{})0.5\delta ,{\displaystyle \frac{1}{K}}{\displaystyle \underset{j=1}{\overset{Kx_{N1}^K}{}}}\xi _N^j=0\right\}`$
$`\left\{\varrho _{N1}(x_{}^k,u_{})0.5\delta ,{\displaystyle \frac{1}{K}}{\displaystyle \underset{j=1}{\overset{K(u_{N1}+\delta )}{}}}\xi _N^j=0\right\}`$
$`=\left\{\varrho _{N1}(x_{}^k,u_{})0.5\delta ,{\displaystyle \underset{j=1}{\overset{K(u_{N1}+\delta )}{}}}\xi _N^j=0\right\}.`$
The sets $`𝔄_1=\left\{\varrho _{N1}(x_{}^k,u_{})0.5\delta \right\}`$ and $`𝔄_2=\left\{_{j=1}^{K(u_{N1}+\delta )}\xi _N^j=0\right\}`$ are independent, so that,
$$𝖯\left(\varrho _N(x_{}^k,u_{})\delta \right)𝖯\left(\varrho _{N1}(x_{}^k,u_{})0.5\delta \right)𝖯^{K(u_{N1}+\delta )}\left(\xi _1^1=0\right).$$
Consequently,
$$\underset{¯}{\mathrm{lim}}_{\delta 0}\underset{¯}{\mathrm{lim}}_K\mathrm{}\frac{1}{K}\mathrm{log}𝖯\left(\varrho _N(x_{}^k,u_{})\delta \right)$$
$$\underset{¯}{\mathrm{lim}}_{\delta 0}\underset{¯}{\mathrm{lim}}_K\mathrm{}\frac{1}{K}\mathrm{log}𝖯\left(\varrho _{N1}(x_{}^k,u_{})0.5\delta \right)+u_{N1}\mathrm{log}𝖯(\xi _1^1=0).$$
If
$$\underset{¯}{\mathrm{lim}}_{\delta 0}\underset{¯}{\mathrm{lim}}_K\mathrm{}\frac{1}{K}\mathrm{log}𝖯\left(\varrho _{N1}(x_{}^K,u_{})\delta \right)J_{N1}(u_{}),$$
(4.5)
provided that $`u_n>0`$, $`nN1`$, the required lower bound holds true.
Thus, it is left to verify the validity of (4.5).
Set $`\mathrm{\Lambda }_{N1}(x_{}^K)=𝔷_{N1},`$ that is,
$$\mathrm{\Lambda }_{N1}(x_{}^K)=e^{_{n=1}^{N1}K\left[t_n^{}x_n^Kx_{n1}^K𝔤(t_n^{})\right]},𝖤\mathrm{\Lambda }_{N1}(x_{}^K)=1.$$
We introduce the probability measure $`𝖰_{N1}^K`$ with $`d𝖰_{N1}^K=\mathrm{\Lambda }_{N1}(x_{}^K)d𝖯.`$ Since $`\mathrm{\Lambda }_{n1}(x_{}^K)>0`$, $`𝖯`$-a.s., we also have $`d𝖯=\mathrm{\Lambda }_{n1}^1(x_{}^K)d𝖰_{n1}^K.`$
In particular, for $`𝔄=\left\{\varrho _{N1}(x_{}^K,u_{})\delta \right\},`$
$$𝖯(𝔄)=_𝔄\mathrm{\Lambda }_{N1}^1(x_{}^K)𝑑𝖰_{N1}^K.$$
So, the following lower bound, on the set $`𝔄`$, is valid:
$`\mathrm{\Lambda }_{N1}^1(x_{}^K)`$ $`e^{KJ_{N1}(u_{})K\delta \mathrm{max}_{nN1}(|t_n^{}|+|𝔤(t_n^{})|)}`$
$`e^{KJ_{N1}(u_{})K\delta \mathrm{max}_{nN1}(|t_n^{}|+|𝔤(t_n^{})|)}`$
or, equivalently,
$$\frac{1}{K}\mathrm{log}𝖯(𝔄)J_{N1}(u_{})\delta \underset{nN1}{\mathrm{max}}(|t_n^{}|+|𝔤(t_n^{})|)+\frac{1}{K}\mathrm{log}𝖰_{N1}^K(𝔄).$$
The latter inequality implies (4.5) if
$$\underset{K\mathrm{}}{lim}\frac{1}{K}\mathrm{log}𝖰_{N1}^K(𝔄)=0.$$
(4.6)
A simple condition, providing (4.6), is $`lim_K\mathrm{}𝖰_{N1}^K(𝔄)=1`$ or, equivalently,
$$\underset{K\mathrm{}}{lim}𝖰_{N1}^K\left(\varrho _{N1}(x_{}^K,u_{})>\delta \right)=0.$$
(4.7)
We verify (4.7) by showing<sup>1</sup><sup>1</sup>1$`𝖤_{N1}^K`$ denotes the expectation with respect to $`𝖰_{N1}^K`$
$$𝖤_{N1}^K\varrho _{N1}^2(x^K,u_{})=\frac{u_{N1}}{K}\underset{n=1}{\overset{N1}{}}\frac{u_{n1}}{u_n^2}𝔤^{\prime \prime }(t_n^{}).$$
(4.8)
Notice that the positiveness of $`(u_n)_{nN1}`$ provides a boundedness for the right hand side of (4.8)‘ and, in turn by Chebyshev’s inequality, the validity of (4.7).
In order to establish (4.8), we apply the identity relative to $`t_n^{}`$:
$$1=𝖤\left(\frac{\mathrm{\Lambda }_n(x_{}^K)}{\mathrm{\Lambda }_{n1}(x_{}^K)}|_{n1}\right)=𝖤e^{K\left[t_n^{}x_n^Kx_{n1}^K𝔤(t_n^{})\right]}.$$
(4.9)
Differentiating twice (4.9) in $`t_n^{}`$, we find that
$`0`$ $`=𝖤\left([x_n^Kx_{n1}^K𝔤^{}(t_n^{})]{\displaystyle \frac{\mathrm{\Lambda }_i(x_{}^K)}{\mathrm{\Lambda }_{n1}(x_{}^K)}}|_{n1}\right)`$ (4.10)
$`0`$ $`=𝖤\left(\left\{K[x_n^Kx_{n1}^K𝔤^{}(t_n^{})]^2x_{n1}^K𝔤^{\prime \prime }(t_n^{})\right\}{\displaystyle \frac{\mathrm{\Lambda }_n(x_{}^K)}{\mathrm{\Lambda }_{n1}(x_{}^K)}}|_{n1}\right).`$
By the Bayes formula, e.g. , : for any integrable random variable $`\alpha `$,
$$𝖤_{N1}^K(\alpha |_{n1})=𝖤\left(\alpha \frac{\mathrm{\Lambda }_n(x_{}^K)}{\mathrm{\Lambda }_{n1}(x_{}^K)}|_{n1}\right).$$
By taking $`\alpha =x_n^K`$ and $`\alpha =[x_n^Kx_{n1}^K𝔤^{}(t_i^{})]^2`$, we derive with the help of (4.10) that
$$𝖤_{N1}^K(x_n^K|_{n1})=x_{n1}^K𝔤^{}(t_n^{})$$
(4.11)
$$𝖤_{N1}^K\left([x_n^Kx_{n1}^K𝔤^{}(t_n^{})]^2|_{n1}\right)=x_{n1}^K\frac{𝔤^{\prime \prime }(t_n^{})}{K}$$
(4.12)
Since $`u_n,u_{n1}`$ are positive, we have $`𝔤^{}(t_n^{})={\displaystyle \frac{u_n}{u_{n1}}}`$. Hence and by (4.11), we obtain that $`𝖤_{N1}^Kx_n^K=\frac{u_n}{u_{i1}}𝖤_{N1}^Kx_{n1}^K.`$ Consequently, iterating the above recursion and taking into account $`u_0=1`$, we find that
$$𝖤_{N1}^Kx_n^K=u_n.$$
Further, with the help of (4.12) we find a recursion
$$𝖤_{N1}^K(x_n^K)^2=\left(\frac{u_n}{u_{n1}}\right)^2𝖤_{N1}^K(x_{n1}^K)^2+u_{n1}\frac{𝔤^{\prime \prime }(t_n^{})}{K}.$$
By using $`𝖤_{N1}^K(x_n^Ku_n)^2=𝖤_{N1}^K(x_n^K)^2u_n^2`$ and and $`u_n^2=\left(\frac{u_n}{u_{n1}}\right)^2u_{n1}^2`$, we establish a recursion for $`\mathrm{}_n=𝖤_{N1}^K(x_n^Ku_n)^2`$:
$$\mathrm{}_n=\left(\frac{u_n}{u_{n1}}\right)^2\mathrm{}_{n1}+u_{n1}\frac{𝔤^{\prime \prime }(t_n^{})}{K}$$
supplied by $`\mathrm{}_0=0`$. Then, $`\frac{\mathrm{}_0}{u_0^2}=0`$ and
$$\frac{\mathrm{}_n}{u_n^2}=\frac{\mathrm{}_{n1}}{u_{n1}^2}+\frac{u_{n1}}{u_n^2}\frac{𝔤^{\prime \prime }(t_n^{})}{K},\mathrm{}_{N1}=\frac{u_{N1}}{K}\underset{n=1}{\overset{N1}{}}\frac{u_{n1}}{u_n^2}𝔤^{\prime \prime }(t_n^{}).$$
It is left to recall that $`\mathrm{}_{N1}=𝖤_{N1}^K\varrho _{N1}^2(x^K,u_{}).`$
## 5. LDP in Continuous Time
We introduce the filtration $`(_t^B)_{t0}`$ generated by Brownian motion $`B_t`$, with the general conditions. All random processes considered in this section are adapted to this filtration. Henceforth, by agreement,
$$\frac{0}{0}=0.$$
###### Theorem 5.1.
For any $`T>0`$, the family $`\{(x_t^K)_{tT}\}_K\mathrm{}`$ obeys the LDP in $`_{[0,T]}(_+)`$, supplied by the uniform metric $`\varrho _T`$, with the speed $`\frac{1}{K}`$ and the rate function
$$J_T(u_{})=\{\begin{array}{ccc}\frac{1}{2\sigma ^2}_0^T\frac{(\dot{u}_t\alpha u_t)^2}{u_t}𝑑t,\hfill & u_0=1,du_t=\dot{u}_tdt,\hfill & \\ \mathrm{},\hfill & \text{otherwise}.\hfill & \end{array}$$
###### Remark 3.
Since $`u_t0`$, Freidlin-Wentzell’s rate function, , $`\frac{1}{2\sigma ^2}_0^T\frac{(\dot{u}_t\alpha u_t)^2}{u_t}𝑑t`$ is not compatible with $`u_t=0`$. Our branching diffusion model is a very particular case of a model studied by Puhalskii in . To apply the LDP analysis from to the family $`\{(x_t^K)_{tT}\}_K\mathrm{}`$, one has to “disentangle” many details of the proof to make it compatible with our case. Finally, in Donati-Martin et all, , the LDP analysis deals with a rate function of the following type $`_0^T\frac{(\dot{u}_t\rho )^2}{u_t}𝑑t`$ for $`u_t0`$ related to a family of diffusion type processes without extinction. A reader interested in details of the direct proof can find them below.
###### Proof.
It suffices to verify:
(i) $`C`$-exponential tightness (see ),
$$\underset{C\mathrm{}}{lim}\overline{\mathrm{lim}}_K\mathrm{}\frac{1}{K}\mathrm{log}𝖯\left(\underset{tT}{sup}x_t^KC\right)=\mathrm{},$$
(5.1)
$$\underset{\mathrm{\Delta }0}{lim}\overline{\mathrm{lim}}_K\mathrm{}\underset{\gamma T}{sup}\frac{1}{K}\mathrm{log}𝖯\left(\underset{t\mathrm{\Delta }}{sup}|x_{\gamma +t}^Kx_\gamma ^K|\eta \right)=\mathrm{},\eta >0,$$
(5.2)
where $`\gamma `$ is stopping time relative to $`(_t^B)_{t0}`$,
(ii) the Local LDP,
$$\underset{\delta 0}{lim}\underset{K\mathrm{}}{lim}\frac{1}{K}\mathrm{log}𝖯\left(\underset{tT}{sup}\left|x_t^Ku_t\right|\delta \right)=J_T(u_{}).$$
(i)-Verification. The Itô equation (1.5) is equivalent to the integral equation $`x_t^K=e^{\alpha t}\left(1+\frac{1}{\sqrt{K}}_0^te^{\alpha s}\sqrt{x_s^K}𝑑B_s\right).`$ Hence,
$$\underset{tT}{sup}x_t^K2e^{|\alpha |T}\left(1\frac{\sigma }{\sqrt{K}}\underset{tT}{sup}_0^te^{\alpha s}\sqrt{x_s^K}𝑑B_s\right),$$
(5.3)
so that, (5.1) holds true provided that
$$\underset{C\mathrm{}}{lim}\overline{\mathrm{lim}}_K\mathrm{}\frac{1}{K}\mathrm{log}𝖯\left(\underset{tT}{sup}_0^te^{\alpha s}\sqrt{x_s^K}𝑑B_s\sqrt{K}C\right)=\mathrm{}.$$
(5.4)
In order to verify (5.4), let us introduce a continuous martingale and its variation process
$$M_t=\frac{\sigma }{\sqrt{K}}_0^te^{\alpha s}\sqrt{x_s^K}𝑑B_s\text{and}M_t=\frac{\sigma ^2}{K}_0^te^{2\alpha s}x_s^K𝑑s$$
respectively and the stopping time $`\tau _C=inf\{tT:M_tC\}`$, where $`inf\{\mathrm{}\}=\mathrm{}`$ which enables us to claim that (5.4) is valid if
$$\underset{C\mathrm{}}{lim}\overline{\mathrm{lim}}_K\mathrm{}\frac{1}{K}\mathrm{log}𝖯\left(\tau _CT\right)=\mathrm{}.$$
(5.5)
We proceed with verification of (5.5). With $`\lambda >0`$, set
$$𝔷_t=e^{\lambda M_t\frac{1}{2}M_t}.$$
It is well known that the process $`(𝔷_t,_t^B)_{t0}`$ is the positive local martingale and so, the supermartingale too with $`𝖤𝔷_\theta 1`$ for any stopping time $`\theta `$ relative to $`(_t^B)`$. By choosing $`\theta =\tau _C`$, we find that $`1𝖤I_{\{\theta T\}}𝔷_\theta .`$ Then, due to a lower bound on the set $`\{\theta T\}`$: $`\mathrm{log}𝔷_\theta \lambda C\frac{\sigma ^2\lambda ^2}{2K}_0^\theta e^{2\alpha s}x_s^K𝑑s`$ and (5.3) there exists positive $`l`$ such that
$$\mathrm{log}𝔷_\theta \lambda C\frac{\sigma ^2\lambda ^2}{2K}_0^\theta e^{2\alpha s}\lambda C\frac{l\lambda ^2}{2K}(1+C).$$
Further, a choice of $`\lambda =\frac{KC}{(1+C)l}`$ implies $`𝔷_\theta e^{\frac{KC^2}{(1+C)l}}`$. Consequently,
$$\frac{1}{K}\mathrm{log}𝖯\left(\tau _CT\right)\frac{C^2}{(1+C)l}\underset{C\mathrm{}}{\overset{}{}}\mathrm{}.$$
By (5.1), the proof of (5.2) is reduced to the verification of two conditions: for any $`\eta ,C>0`$,
$$\underset{\mathrm{\Delta }0}{lim}\overline{\mathrm{lim}}_K\mathrm{}\underset{\gamma T}{sup}\frac{1}{K}\mathrm{log}𝖯\left(\underset{t\mathrm{\Delta }}{sup}_\gamma ^{\gamma +t}x_s^K𝑑s\eta ,\underset{sT}{sup}x_s^KC\right)=\mathrm{}$$
$$\underset{\mathrm{\Delta }0}{lim}\overline{\mathrm{lim}}_K\mathrm{}\underset{\gamma T}{sup}\frac{1}{K}\mathrm{log}𝖯\left(\frac{\sigma }{\sqrt{K}}\underset{t\mathrm{\Delta }}{sup}\left|_\gamma ^{\gamma +t}\sqrt{x_s^K}𝑑B_s\right|\eta ,\underset{sT}{sup}x_s^KC\right)=\mathrm{}.$$
The first is obvious while the second is equivalent to
$$\underset{\mathrm{\Delta }0}{lim}\overline{\mathrm{lim}}_K\mathrm{}\underset{\gamma T}{sup}\frac{1}{K}\mathrm{log}𝖯\left(\underset{t\mathrm{\Delta }}{sup}I_{T,C}\left|M_{\gamma +t}^KM_\gamma ^K\right|\eta \right)=\mathrm{},$$
(5.6)
where $`I_{t,C}=I_{\{sup_{st}x_s^KC\}}`$, $`tT`$.
Set $`N_t^K=M_{\gamma +t}^KM_\gamma ^K`$ and notice that $`(N_t^K,_{\gamma +t}^B)_{t0}`$ is a local martingale with the variation process $`N^K_t=\frac{\sigma ^2}{K}_\gamma ^{\gamma +t}x_s^K𝑑s`$.
Further, the use of $`I_{T,C}N_t^K=I_{T,C}_0^tI_{s,C}𝑑N_s^K`$ simplifies (5.6) to
$$\underset{\mathrm{\Delta }0}{lim}\overline{\mathrm{lim}}_K\mathrm{}\underset{\gamma T}{sup}\frac{1}{K}\mathrm{log}𝖯\left(\underset{t\mathrm{\Delta }}{sup}\left|_0^tI_{s,C}𝑑N_s^K\right|\eta \right)=\mathrm{}.$$
(5.7)
The local martingale $`N_t^{K,C}:=_0^tI_{s,C}𝑑N_s^K`$ possesses the variation process
$$N^{K,C}_t=_0^tI_{s,C}dN^K_s=\frac{\sigma ^2}{K}_0^tI_{s,C}x_s^K𝑑s,$$
that is, d$`N^{K,C}_t\frac{\sigma ^2C}{K}dt.`$
Now, we are able to verify (5.7) with the help of stochastic exponential technique. Let
$$𝔷_t(\lambda )=e^{\lambda N_t^{K,C}\frac{\lambda ^2}{2}N^{K,C}_t},\lambda .$$
Since $`𝔷_t(\lambda )`$ is a continuous local martingale and supermartingale too, for any stopping time $`\theta `$, $`𝖤𝔷_\theta (\lambda )1.`$ Let $`\theta =inf\{t\mathrm{\Delta }:N_t^{K,C}\eta \}`$. Taking into account that $`\{\theta \mathrm{\Delta }\}=\{N_\theta ^{K,C}\eta \}`$, write $`1𝖤I_{\{\theta \mathrm{\Delta }\}}𝔷_\theta (\lambda ).`$ The value $`𝔷_\theta (\lambda )`$ is evaluated below on the set $`\{\theta \mathrm{\Delta }\}`$ as follows: with $`\lambda >0`$ and $`N^{K,C}_\theta \frac{\sigma ^2C}{K}\theta \frac{\sigma ^2C}{K}\mathrm{\Delta },`$
$$𝔷_\theta (\lambda )e^{\lambda \eta \frac{\lambda ^2\sigma ^2C}{2K}\mathrm{\Delta }}.$$
Therefore, $`\mathrm{log}𝖯(\theta \mathrm{\Delta })\left[\lambda \eta \frac{\lambda ^2\sigma ^2C}{2K}\mathrm{\Delta }\right]`$ and the choice of $`\lambda =\frac{K\eta }{\sigma ^2C\mathrm{\Delta }}`$ provides
$$\frac{1}{K}\mathrm{log}𝖯(\theta \mathrm{\Delta })\frac{\eta ^2}{2\sigma ^2C\mathrm{\Delta }}\underset{\mathrm{\Delta }0}{\overset{}{}}\mathrm{}.$$
It is clear that the same result remains valid for $`\theta =inf\{t:N_t^{K,C}\eta \}`$. Combining both, we obtain (5.7).
(ii)-Verification. The upper bound. For $`u_01`$ or $`du_t\ll ̸dt`$, the proof is obvious. For $`u_0=1`$ and $`du_t=\dot{u}_tdt`$, the stochastic exponential technique is applicable. With an absolutely continuous deterministic function $`\lambda (t)`$ let us introduce a continuous martingale $`M_t`$ and its predictable variation process $`M_t`$:
$$M_t=\frac{\sigma }{\sqrt{K}}_0^t\lambda (s)\sqrt{x_s^K}𝑑B_s\text{and}M_t=\frac{\sigma ^2}{K}_0^t\lambda ^2(s)x_s^K𝑑s.$$
It is well known that the stochastic exponential $`𝔷_t=e^{M_t0.5M_t}`$ is a local martingale and a supermartingale too with $`𝖤𝔷_T1`$. The use of this property implies
$$1𝖤I_{\{sup_{tT}|x_t^Ku_t|\delta \}}𝔷_T.$$
(5.8)
The next helpful step of the proof gives a deterministic lower bound for $`𝔷_T`$ on the set $`\{sup_{tT}|x_t^Ku_t|\delta \}=:𝔄_\delta .`$ By (1.4), $`M_t=_0^t\lambda (s)(dx_s^K\alpha x_s^Kds)`$, so that,
$`\mathrm{log}𝔷_T`$ $`={\displaystyle _0^T}\lambda (s)(dx_s^K\alpha x_s^Kds){\displaystyle \frac{\sigma ^2}{2K}}{\displaystyle _0^t}\lambda ^2(s)x_s^K𝑑s`$
$`={\displaystyle _0^T}\left[\lambda (s)(\dot{u}_s\alpha u_sds){\displaystyle \frac{\sigma ^2}{2K}}\lambda ^2(s)u_s\right]𝑑s`$
$`+{\displaystyle _0^T}\lambda (s)d(x_s^Ku_s)[=\lambda _T(x_T^Ku_T){\displaystyle _0^T}(x_s^Ku_s)\dot{\lambda }_t𝑑s]`$
$`{\displaystyle _0^T}\left[\lambda (s)\alpha \left\{x_s^Ku_s\right\}+{\displaystyle \frac{\sigma ^2}{2K}}\lambda ^2(s)\left\{x_s^Ku_s\right\}\right]𝑑s.`$
Now, by taking $`\lambda (s)=K\theta (s)`$, we find a lower bound of $`𝔷_T`$ on the set $`𝔄_\delta :=\{sup_{tT}|x_t^Ku_t|\delta \}`$,
$`\mathrm{log}𝔷_T`$ $`K{\displaystyle _0^T}\left[\theta (s)(\dot{u}_s\alpha u_s){\displaystyle \frac{\sigma ^2}{2}}\theta ^2(s)u_s\right]𝑑s`$
$`\delta K\left[|\theta _T|+{\displaystyle _0^T}\left(|\dot{\theta }_s|+|\alpha \theta (s)|+{\displaystyle \frac{\sigma ^2\theta ^2(s)}{2}}\right)𝑑s\right].`$
This lower bound jointly with (5.8) implies the following upper bound: for any absolutely continuous deterministic function $`\theta (s)`$,
$$\overline{\mathrm{lim}}_{\delta 0}\overline{\mathrm{lim}}_K\mathrm{}\frac{1}{K}\mathrm{log}𝖯\left(𝔄_\delta \right)_0^T\left[\theta (s)(\dot{u}_s\alpha u_s)\frac{\sigma ^2}{2}\theta ^2(s)u_s\right]𝑑s.$$
Since $`u_s`$ is only nonnegative, it makes sense, for computational convenience, to use a corrected upper bound, with $`\epsilon >0`$
$$\overline{\mathrm{lim}}_{\delta 0}\overline{\mathrm{lim}}_K\mathrm{}\frac{1}{K}\mathrm{log}𝖯\left(𝔄_\delta \right)_0^T\left[\theta (s)(\dot{u}_s\alpha u_s)\frac{\sigma ^2}{2}\theta ^2(s)(u_s+\epsilon )\right]𝑑s.$$
(5.9)
If $`\dot{u}_t`$ is absolutely continuous function, a choice of $`\theta (s)={\displaystyle \frac{\dot{u}_s\alpha u_s}{\sigma ^2(u_s+\epsilon )}}`$ provides
$`\overline{\mathrm{lim}}_{\delta 0}\overline{\mathrm{lim}}_K\mathrm{}{\displaystyle \frac{1}{K}}\mathrm{log}𝖯\left(𝔄_\delta \right)`$ $`{\displaystyle \frac{1}{2\sigma ^2}}{\displaystyle _0^T}{\displaystyle \frac{(\dot{u}_s\alpha u_s)^2}{u_s+\epsilon }}𝑑s`$
$`{\displaystyle \frac{1}{2\sigma ^2}}{\displaystyle _0^T}{\displaystyle \frac{(\dot{u}_s\alpha u_s)^2}{u_s}}𝑑s,\epsilon 0.`$
In general case, one can choose a sequence $`\theta _n(s)`$, $`n1`$ of absolutely continuous functions such that
$$\underset{n\mathrm{}}{lim}\left[\theta _n(s)(\dot{u}_s\alpha u_s)\frac{\sigma ^2}{2}\theta _n^2(s)(u_s+\epsilon )\right]=\underset{\varphi }{sup}\left[\varphi (\dot{u}_s\alpha u_s)\frac{\sigma ^2}{2}\varphi ^2(u_s+\epsilon )\right]$$
$$=\frac{1}{2\sigma ^2}\frac{(\dot{u}_s\alpha u_s)^2}{u_s+\epsilon }.$$
Hence, for sufficiently large $`n`$, $`\left[\theta _n(s)(\dot{u}_s\alpha u_sds)\frac{\sigma ^2}{2}\theta _n^2(s)(u_s+\epsilon )\right]0.`$ Then, due to (5.9) being valid with $`\theta (s)`$ replaced by $`\theta _n(s)`$, and Fatou’s theorem, we find that
$$\overline{\mathrm{lim}}_{\delta 0}\overline{\mathrm{lim}}_K\mathrm{}\frac{1}{K}\mathrm{log}𝖯\left(𝔄_\delta \right)\underset{¯}{\mathrm{lim}}_n\mathrm{}_0^T\left[\theta _n(s)(\dot{u}_s\alpha u_s)\frac{\sigma ^2}{2}\theta _n^2(s)(u_s+\epsilon )\right]𝑑s$$
$$_0^T\underset{¯}{\mathrm{lim}}_n\mathrm{}\left[\theta _n(s)(\dot{u}_s\alpha u_s)\frac{\sigma ^2}{2}\theta _n^2(s)(u_s+\epsilon )\right]ds=\frac{1}{2\sigma ^2}_0^T\frac{(\dot{u}_s\alpha u_s)^2}{u_s+\epsilon }$$
$$\frac{1}{2\sigma ^2}_0^T\frac{(\dot{u}_s\alpha u_s)^2}{u_s},\epsilon 0.$$
(ii)-Verification. The proof of
$$\underset{¯}{\mathrm{lim}}_{\delta 0}\underset{¯}{\mathrm{lim}}_K\mathrm{}\frac{1}{K}\mathrm{log}𝖯\left(\underset{tT}{sup}|x_s^Ku_s|\delta \right)\frac{1}{2\sigma ^2}_0^T\frac{(\dot{u}_t\alpha u_t)^2}{u_t}𝑑t$$
(5.10)
is done in three steps.
1. It suffices to analyse the case $`_0^T\frac{(\dot{u}_s\alpha u_s)^2}{u_s}𝑑s<\mathrm{}`$, which enables us to consider only those test functions that remain zero after arriving at zero. In other words, we shall give the proof of (5.10) for absolutely continuous $`u_{}`$ with $`u_0=1`$ and $`(u_t>0)_{t<T}`$, $`u_T0`$.
2. Set $`\tau _C=inf\{tT:x_t^KC\},\text{where }inf\{\mathrm{}\}=\mathrm{}`$ and notice that if for any $`C>0`$
$$\underset{¯}{\mathrm{lim}}_{\delta 0}\underset{¯}{\mathrm{lim}}_K\mathrm{}\frac{1}{K}\mathrm{log}𝖯\left(\underset{tT\tau _C}{sup}|x_s^Ku_s|\delta \right)\frac{1}{2\sigma ^2}_0^T\frac{(\dot{u}_t\alpha u_t)^2}{u_t}𝑑t,$$
(5.11)
then (5.10) holds. This can be seen as follows. Since
$$𝔄_\delta \left\{\underset{tT\tau _C}{sup}|x_t^Ku_t|\delta \right\}\{\tau _C=\mathrm{}\}$$
$$=\left\{\underset{tT\tau _C}{sup}|x_t^Ku_t|\delta \right\}\left\{\underset{tT\tau _C}{sup}|x_t^Ku_t|\delta \right\}\{\tau _CT\},$$
we have $`\{\tau _CT\}𝔄_\delta \left\{sup_{tT\tau _C}|x_t^Ku_t|\delta \right\}`$, so that,
$$2\left[𝖯\left(𝔄_\delta \right)𝖯\left(\tau _CT\right)\right]𝖯\left(\underset{tT\tau _C}{sup}|x_s^Ku_s|\delta \right).$$
Hence, due to (5.11),
$$\underset{¯}{\mathrm{lim}}_{\delta 0}\underset{¯}{\mathrm{lim}}_K\mathrm{}\frac{1}{K}\mathrm{log}𝖯\left(𝔄_\delta \right)\underset{C\mathrm{}}{lim}\underset{¯}{\mathrm{lim}}_K\mathrm{}\frac{1}{K}\mathrm{log}𝖯\left(\tau _CT\right)$$
$$\frac{1}{2\sigma ^2}_0^T\frac{(\dot{u}_t\alpha u_t)^2}{u_t}𝑑t$$
and it is left to recall that $`\{\tau _CT\}=\left\{sup_{tT}x_t^KC\right\}`$ and to refer to (5.1).
3. By 1., $`_0^T\dot{u}_t^2𝑑t<\mathrm{}`$. We proceed with the verification of (5.11). Define a continuous martingale $`M_t`$ and its variation process $`M_t`$: with $`\epsilon >0`$,
$$M_t=_0^{t\tau _C}\sqrt{K}\frac{\dot{u}_s\alpha x_s^K}{\sigma \sqrt{x_s^K+\epsilon }}𝑑B_s\text{and}M_t=_0^{t\tau _C}K\frac{(\dot{u}_s\alpha x_s^K)^2}{\sigma ^2(x_s^K+\epsilon )}𝑑s.$$
By definition of $`\tau _C`$, we have $`M_T\frac{2K}{\sigma ^2\epsilon }_0^T\left(\dot{u}_t^2+\alpha ^2C^2\right)𝑑s<\mathrm{},`$ so that the stochastic exponential $`(𝔷_t,_t^B,𝖯)_{tT}`$ with $`𝔷_t=e^{M_t0.5M_t}`$ is a uniformly integrable martingale, $`𝖤𝔷_T=1`$. We use the latter property to define a new probability measure $`\overline{𝖯}`$ on $`(\mathrm{\Omega },_T^B)`$ by letting $`d\overline{𝖯}=𝔷_Td𝖯`$ and apply
$$𝖯\left(\underset{tT\tau _C}{sup}|x_t^Ku_t|\delta \right)=\underset{\{sup_{tT\tau _C}|x_t^Ku_t|\delta \}}{}𝔷_T^1𝑑\overline{𝖯}$$
for verification of (5.11). This approach heavily uses a semimartingale description of the processes $`(x_t^K,_t^B,\overline{𝖯})_{tT}`$ and $`(𝔷_t^1,_t^B,\overline{𝖯})_{tT}`$. We begin with the process $`(B_t,_t^B,\overline{𝖯})_{tT}`$. The random processes $`(B_t,_t^B,𝖯)_{tT}`$ and $`(𝔷_t,_t^B,𝖯)_{tT}`$ are continuous martingales and, in particular,
$$d𝔷_t=I_{\{\tau _Ct\}}𝔷_t\sqrt{K}\frac{\dot{u}_t\alpha x_t^K}{\sigma \sqrt{x_t^K+\epsilon }}dB_t.$$
Hence, the co-variation process for $`𝔷_t`$, $`B_t`$ is defined as:
$$𝔷,B_t=_0^{t\tau _C}𝔷_s\sqrt{K}\frac{\dot{u}_s\alpha x_s^K}{\sigma \sqrt{x_s^K+\epsilon }}𝑑s.$$
It is well known (see, e.g. Ch. 4, §5 in ) that the random process $`(\overline{B}_t,_t^B,\overline{𝖯})_{tT}`$ with
$$\overline{B}_t=B_t_0^t𝔷_s^1d𝔷,B_s=B_t_0^{t\tau _C}\sqrt{K}\frac{\dot{u}_s\alpha x_s^K}{\sigma (x_s^K+\epsilon )}𝑑s$$
is a Brownian motion. Consequently, we find that, $`\overline{𝖯}`$-a.s.,
$`x_t^K`$ $`=1+{\displaystyle _0^t}I_{\{\tau _Cs\}}\dot{u}_s𝑑s+{\displaystyle _0^t}\alpha x_s^K\left[1I_{\{\tau _Cs\}}\left({\displaystyle \frac{x_s^K}{x_s^K+\epsilon }}\right)^{0.5}\right]𝑑s`$ (5.12)
$`+{\displaystyle _0^t}{\displaystyle \frac{\sigma }{\sqrt{K}}}\sqrt{x_s^K}𝑑\overline{B}_s`$
$`\mathrm{log}z_t^1`$ $`={\displaystyle _0^{t\tau _C}}\sqrt{K}{\displaystyle \frac{\dot{u}_s\alpha x_s^K}{\sigma \sqrt{x_s^K+\epsilon }}}𝑑\overline{B}_s{\displaystyle \frac{1}{2}}{\displaystyle _0^{t\tau _C}}K{\displaystyle \frac{(\dot{u}_s\alpha x_s^K)^2}{\sigma ^2(x_s^K+\epsilon )}}𝑑s.`$
Now, we evaluate from below the value $`\frac{1}{K}\mathrm{log}𝔷_T^1`$ on the set $`\{sup_{tT\tau _V}|x_t^Ku_t|\}`$. Write
$`{\displaystyle \frac{1}{K}}\mathrm{log}𝔷_T^1`$ $`{\displaystyle \frac{1}{2\sigma ^2}}{\displaystyle _0^T}{\displaystyle \frac{(\dot{u}_s\alpha u_s)^2}{u_s}}𝑑s+h(C,\epsilon ,\delta )`$
$`{\displaystyle \frac{1}{\sqrt{K}}}\underset{tT}{sup}\left|{\displaystyle _0^{t\tau _C}}{\displaystyle \frac{\dot{u}_s\alpha x_s^K}{\sigma \sqrt{x_s^K+\epsilon }}}𝑑\overline{B}_s\right|,`$
where $`h(C,\epsilon ,\delta )\underset{\delta 0}{\overset{}{}}0`$. Therefore, (5.12) can be transformed into (here $`\eta `$ is a positive constant)
$`{\displaystyle \frac{1}{K}}\mathrm{log}𝖯\left(\underset{tT\tau _C}{sup}|x_t^Ku_t|\delta \right){\displaystyle \frac{1}{2\sigma ^2}}{\displaystyle _0^T}{\displaystyle \frac{(\dot{u}_s\alpha u_s)^2}{u_s}}𝑑s+h(C,\epsilon ,\delta )`$
$`+{\displaystyle \frac{1}{K}}\mathrm{log}\overline{𝖯}(\underset{tT\tau _C}{sup}|x_t^Ku_t|\delta ,{\displaystyle \frac{1}{\sqrt{K}}}\underset{tT}{sup}|{\displaystyle _0^{t\tau _C}}{\displaystyle \frac{\dot{u}_s\alpha x_s^K}{\sigma \sqrt{x_s^K+\epsilon }}}d\overline{B}_s\eta \left|\right).`$
This lower bound makes it possible to claim that (5.11) holds true, provided that
$$\underset{K\mathrm{}}{lim}\overline{𝖯}\left(\underset{tT}{sup}\left|_0^{t\tau _C}\frac{\dot{u}_s\alpha x_s^K}{\sigma \sqrt{x_s^K+\epsilon }}𝑑\overline{B}_s\right|>\sqrt{K}\eta \right)=0$$
(5.13)
$$\underset{\epsilon 0}{lim}\overline{\mathrm{lim}}_K\mathrm{}\overline{𝖯}\left(\underset{tT\tau _C}{sup}|x_t^Ku_t|>\delta \right)=0.$$
(5.14)
Since $`I_{\{\tau _st\}}\frac{(\dot{u}_s\alpha x_s^K)^2}{\sigma ^2(x_s^K+\epsilon )}\frac{(|\dot{u}_s|+C)^2}{\sigma ^2\epsilon }`$ the Doob inequality (here $`\overline{𝖤}`$ is the expectation relative to $`\overline{𝖤}`$)
$$\overline{𝖤}\left(\underset{tT}{sup}\right|_0^{t\tau _C}\frac{\dot{u}_s\alpha x_s^K}{\sigma \sqrt{x_s^K+\epsilon }}d\overline{B}_s>\sqrt{K}\eta \left|\right)\frac{4}{K\eta ^2}_0^T\frac{(|\dot{u}_s|+C)^2}{\sigma ^2\epsilon }ds,$$
jointly with 3., establish (5.13).
Due to the first part of (5.12), the proof of (5.14) is reduced to the verification of
$$\underset{\epsilon 0}{lim}\underset{0xC}{sup}x\left[1\left(\frac{x}{x+\epsilon }\right)^{0.5}\right]=0,$$
which is obvious, and
$$\underset{K\mathrm{}}{lim}\overline{𝖯}\left(\underset{tT}{sup}\right|_0^{t\tau _C}\sqrt{x_s^K}d\overline{B}_s>\sqrt{K}\eta \left|\right)=0,$$
which is similar to the proof of (5.13). |
warning/0507/math0507389.html | ar5iv | text | # Theorem .
## 1. Construction
We refer to \[B\] for background on Coxeter groups, their matrices and graphs.
Let $`G_1`$ be a group generated by 5 involutions $`\{r_1,\mathrm{},r_5\}=S_1`$ and relations
$`(r_ir_j)^5=1`$ if $`(i,j)=(4,5)`$;
$`(r_ir_j)^3=1`$ if $`(i,j)=(1,2)`$, $`(2,3)`$ or $`(3,4)`$;
$`(r_ir_j)^2=1`$ otherwise.
$`(G_1,S_1)`$ may be thought as a cocompact group of isometries of the hyperbolic space $`^4`$. Its fundamental polytope is a simplex with Coxeter diagram $`\mathrm{\Sigma }_1`$:
$`\mathrm{\Sigma }_1`$
Now define a group $`G_k`$. It consists of $`k`$ copies of $`G_1`$ (an $`m`$-th copy is generated by 5 involutions $`r_{5(m1)+1},\mathrm{},r_{5m}`$) with the following additional relations between generators of different copies:
$`(r_ir_j)^5=1`$ if $`ij1mod5`$;
$`(r_ir_j)^2=1`$ otherwise.
Let $`S_k`$ be the set of generating reflections $`\{r_1,\mathrm{},r_{5k}\}`$. The Coxeter graph $`\mathrm{\Sigma }_k`$ of $`(G_k,S_k)`$ consists of $`k`$ copies of Coxeter graph $`\mathrm{\Sigma }_1`$ sticking out of a complete graph on $`k`$ vertices with triple edges. The graph $`\mathrm{\Sigma }_3`$ is drawn on Fig. 1.
Since any two involutions $`r_i`$ and $`r_j`$ generate a finite group, the signature of $`\mathrm{\Sigma }_k`$ is uniquely defined. To prove the theorem, we show that $`G_k`$ is word-hyperbolic and the signature of $`\mathrm{\Sigma }_k`$ equals $`(4k,k)`$. Then Lemma 12 of the paper \[V1\] implies that the set of reflections $`S_k`$ does not preserve any quadratic form with negative inertia index less than $`k`$. This means that $`G_k`$ cannot be embedded into $`O(n,k1)`$ for any $`n`$ with all generators being conjugates of $`\{r_i\}`$.
## 2. Main tools
To prove that $`G_k`$ is word-hyperbolic, we use the following result of G. Moussong \[M\].
###### Proposition 1 (\[M\], cor. of Th. 17.1).
Let $`\mathrm{\Sigma }`$ be a Coxeter graph of a Coxeter group $`G`$ without parabolic subgraphs. Then $`G`$ is word-hyperbolic if and only if no pair of unjoined subgraphs $`\mathrm{\Sigma }^{}`$, $`\mathrm{\Sigma }^{\prime \prime }`$ of $`\mathrm{\Sigma }`$ corresponds to two infinite groups.
To compute the determinant of $`\mathrm{\Sigma }_k`$ we use the following result of E. Vinberg \[V2\].
The weight of an $`(m2)`$-tuple edge of $`\mathrm{\Sigma }`$ is a real number equal to $`\mathrm{cos}(\frac{\pi }{m})`$. For any cyclic path $`\gamma `$ of $`\mathrm{\Sigma }`$ denote by $`p(\gamma )`$ the product of the weights of all edges of $`\gamma `$.
###### Proposition 2 (\[V2\], Prop. 11).
The determinant of a Coxeter graph $`\mathrm{\Sigma }`$ is equal to the sum of products of the form
$$(1)^sp(\gamma _1)\mathrm{}p(\gamma _s),$$
where $`\{\gamma _1,\mathrm{},\gamma _s\}`$ ranges over all unordered collections (including the empty one) of pairwise disjoint oriented cyclic paths of $`\mathrm{\Sigma }`$.
## 3. Proof of the theorem
First, we prove that $`G_k`$ is word-hyperbolic. It is clear that any subgraph of $`\mathrm{\Sigma }_k`$ corresponding to infinite subgroup contains at least one vertex of the complete subgraph on $`k`$ vertices. By Prop. 1, this implies word-hyperbolicity of $`G_k`$.
Hiding all vertices of the complete subgraph in $`\mathrm{\Sigma }_k`$, we obtain a graph corresponding to the finite group $`kH_4`$. So, the positive inertia index of $`\mathrm{\Sigma }_k`$ is at least $`4k`$. To prove that the negative inertia index of $`\mathrm{\Sigma }_k`$ equals $`k`$, it is sufficient to show that the sign of $`\text{det}(\mathrm{\Sigma }_k)`$ equals $`(1)^k`$. We finish the proof by the following
###### Lemma .
$$\text{det}(\mathrm{\Sigma }_k)=\frac{(2\sqrt{5})^k}{2^{5k}}(k+1)$$
###### Proof.
Denote $`\text{det}(\mathrm{\Sigma }_k)`$ by $`d_k`$, and denote by $`d`$ the determinant of $`H_4`$. Prop. 2 implies the following recurrent formula for $`d_k`$:
$$d_{k+1}=d_kd_1+\underset{m=1}{\overset{k}{}}(1)^mm!\left(\genfrac{}{}{0pt}{}{k}{m}\right)(\mathrm{cos}(\frac{\pi }{5}))^{m+1}d^{m+1}d_{km}.$$
($``$)
Here $`\left(\genfrac{}{}{0pt}{}{k}{m}\right)`$ is the number of $`(m+1)`$-simplices with one fixed vertex in the complete graph on $`k`$ vertices, and $`m!`$ is the number of oriented Hamiltonian cycles in 1-skeleton of $`(m+1)`$-simplex. Notice now that
$$d_1=\frac{2\sqrt{5}}{16},d=\frac{73\sqrt{5}}{32},\mathrm{cos}(\frac{\pi }{5})=\frac{1\sqrt{5}}{4},d\mathrm{cos}(\frac{\pi }{5})=\frac{2\sqrt{5}}{32},$$
and assume that the lemma is true for all $`lk`$ (we put $`d_0=1`$, so the lemma is true for $`k=1`$). Then the expression ($``$) implies that
$`d_{k+1}=\frac{(2\sqrt{5})^k}{2^{5k}}(k+1)\frac{2\sqrt{5}}{16}+\underset{m=1}{\overset{k}{}}(1)^mm!\left(\genfrac{}{}{0pt}{}{k}{m}\right)\left(\frac{2\sqrt{5}}{32}\right)^{m+1}\frac{(2\sqrt{5})^{km}}{2^{5(km)}}(km+1)=`$
$`=\frac{(2\sqrt{5})^{k+1}}{2^{5(k+1)}}\left(2(k+1)+\underset{m=1}{\overset{k}{}}(1)^m\frac{k!}{(km)!}(km+1)\right)=\frac{(2\sqrt{5})^{k+1}}{2^{5(k+1)}}(k+2).`$
Independent University of Moscow
e-mail: ow felikson@mccme.ru
owowwww pasha@mccme.ru |
warning/0507/nlin0507055.html | ar5iv | text | # A 𝑈_𝑞((𝑔𝑙)̂(2|2))₁-vertex model: Creation algebras and quasi-particles I
## 1. Introduction
This paper continues the study of the $`U_q\left(\widehat{gl}(2|2)\right)`$-vertex model started in . There an integrable vertex model based on the vector representation of $`U_q\left(gl(2|2)\right)`$ and its dual is investigated. In the limit $`q0`$, the action of the corner transfer matrix Hamiltonian on the space of half-infinite configurations takes a trigonal form provided that the configurations obey a particular boundary condition. A one-to-one correspondence between the half-infinite configurations and the weight states of a reducible level-one module of $`U_q\left(\widehat{sl}(2|2)\right)/`$ with grade $`n`$ is observed for $`n3`$. Here the grade corresponds to the diagonal element of the corner transfer matrix Hamiltonian. The reducible module decomposes into one weakly integrable irreducible module and one nonintegrable irreducible module of $`U_q\left(\widehat{sl}(2|2)\right)/`$. Both are highest weight modules.
Here the investigation of half-infinite configurations is supplemented taking into account a second boundary condition. Various choices of composite vertices and $`U_q\left(gl(2|2)\right)`$-weights are considered. At $`q=0`$, a one-to-one correspondence between the half-infinite configurations and the weight states of reducible or irreducible level-one modules of $`U_q\left(\widehat{sl}(2|2)\right)/`$ is found at the grades $`0,1,2`$. This correspondence is assumed to hold true at all grades. Similar as the modules relevant to the first boundary condition, the reducible modules decompose into a weakly integrable and a nonintegrable highest-weight module. These irreducible modules may be assembled in a reducible module in various ways. In each case, one possibility is distinguished by a simple free boson realization within the scheme employed in .
The arguments developed for the algebraic analysis of integrable $`U_q\left(\widehat{gl}(N)\right)`$-vertex models (,,) suggest that the full space of states is interpreted as a sum of tensor products each combining a highest weight module with a suitable dual (and therefore lowest weight) module. The levels of both constituents of the tensor product add to zero. According to the method presented in these references, the row-to-row transfer matrix should be described in terms of type I vertex operators associated with the vector representation $`W`$ of $`U_q\left(gl(2|2)\right)`$ and its dual $`W^{}`$. Type II vertex operators are expected to create the eigenstates of the row-to-row transfer matrix. In general, only a subset of all type II vertex operators existing for the affine quantum algebra at fixed level gives rise to such eigenstates. The correct choice of vertex operators follows from the decomposition of the tensor products modelling the space of states into irreducible components. For the $`U_q\left(\widehat{gl}(N)\right)`$-vertex models this decomposition can be done at $`q=0`$ , making use of the crystal base theory , . The path space of the model is described in terms of a creation algebra whose generators can be seen as formal $`q0`$ limits of the appropriate type II vertex operators. This picture is based on the existence of the $`q0`$ limits of the creation operators acting on the true groundstate of the model. The existence of these limits has been conjectured and checked for the $`XXZ`$-model in .
In this study, a similar description is proposed for an infinite configuration space of the $`U_q\left(\widehat{gl}(2|2)\right)`$-model at $`q=0`$. The second boundary condition is imposed in both directions. In the limit $`q0`$, the diagonal elements of the corner transfer matrix Hamiltonian decouple into a contribution depending only on the $`W`$-part of the half-infinite configuration and into a second part depending only on the $`W^{}`$-part. For the boundary condition chosen here, the $`W`$-components and the $`W^{}`$-components of the infinite configurations can be described separately in terms of two creation algebras $`𝒜`$ and $`𝒜^{}`$. According to their defining relations, bases of both algebras are given by their sets of normal forms. This has been proven in for a quite general type of creation algebras including $`𝒜`$ and $`𝒜^{}`$. The infinite $`W`$-components ($`W^{}`$-components) are in one-to-one correspondence with the normal forms of $`𝒜^{}`$ ($`𝒜`$) supplemented by the unity.
To demonstrate these correspondences, two types of infinite border strips are employed. Border strips consisting of a finite number of rows or columns of finite length framed by either two infinite rows or two infinite columns are termed horizontal or vertical border strips, respectively. The semi-standard supertableaux of the horizontal (or vertical) border strips satisfying a suitable boundary condition are mapped onto the infinite $`W`$-(or $`W^{}`$-)components and vice versa. Then a weight-preserving one-to-one map between the semi-standard supertableaux of the horizontal (vertical) border strips and the set of normal forms of $`𝒜^{}`$ ($`𝒜`$) enlarged by the unity is given explicitely. A similar description can be constructed for the infinite configuration space restricted by the first boundary condition in both directions. If different boundary conditions are imposed in the left and right direction, the simple separation into the $`W`$-and $`W^{}`$-components is lost. These cases will be considered in a separate publication.
The defining relations of $`𝒜`$ and $`𝒜^{}`$ are expected to emerge from the formal $`q0`$ limits of the commutation relations satisfied by particular type II vertex operators. Suitable vertex operators should be associated with a reducible, infinite-dimensional highest weight module $`\stackrel{̊}{V}`$ in the case of $`𝒜`$ and with its dual module $`\stackrel{̊}{V}^{}`$ in the case of $`𝒜^{}`$. The module $`\stackrel{̊}{V}`$ decomposes into a one-dimensional module and an infinite-dimensional irreducible module. The latter coincides with the module $`V`$ introduced in . Based on a free boson realization for one component, the existence of such vertex operators is conjectured. The commutation relations of two vertex operators are governed by the R-matrices related to $`\stackrel{̊}{V}`$ and $`\stackrel{̊}{V}^{}`$. The free boson realizations allow to determine the normalisation of the R-matrices defined on $`\stackrel{̊}{V}\stackrel{̊}{V}`$ or $`\stackrel{̊}{V}^{}\stackrel{̊}{V}^{}`$. Evaluation of their $`q0`$ limits formally leads to the defining relations of $`𝒜`$ or $`𝒜^{}`$. It seems more difficult however to demonstrate the compatibility of the commutation relations involving both types of vertex operators with the separation into $`W`$\- and $`W^{}`$-components found at $`q=0`$. In the mixed case it is less obvious how to make use of $`q0`$ limits of single R-matrix elements. The $`q0`$ limits of the R-matrix elements on $`\stackrel{̊}{V}\stackrel{̊}{V}`$ or $`\stackrel{̊}{V}^{}\stackrel{̊}{V}^{}`$ reflect the structure of the irreducible components of these tensor products. Hence this structure underlies the defining relations of $`𝒜`$ or $`𝒜^{}`$. This motivates the search for an R-matrix action on the irreducible components of the tensor products $`\stackrel{̊}{V}^{}\stackrel{̊}{V}`$ or $`\stackrel{̊}{V}\stackrel{̊}{V}^{}`$. This action is obtained on $`\stackrel{̊}{V}^{}\stackrel{̊}{V}`$ only for a partial range of the spectral parameter. Its formal $`q0`$ limit does not distinguish between different irreducible components.
From the above findings it may be conjectured that a part of the eigenstates of the row-to-row transfer matrix is generated by the type II vertex operators associated with $`\stackrel{̊}{V}`$ and $`\stackrel{̊}{V}^{}`$. A full discussion of the mixed R-matrix elements is more conveniently given in context with the R-matrix elements encountered in the case of mixed boundary conditions and is therefore relegated to a forthcoming publication.
The main findings presented here are formulated by conjecture 1 in section 3, conjectures 2 and 3 in section 6 and by the results 1-3 in section 5.
The paper is organised as follows. For the convenience of the reader, subsection 2.1 collects notations and recalls the free boson realization of $`U_q\left(\widehat{gl}(2|2)\right)`$ at level-one used in . Subsection 2.2 gives a short account of the vertex model and the structure of the half-infinite configuration space subject to the boundary condition chosen in . The second boundary condition is considered in section 3. Since the analysis is quite analogous to the one presented in , the outline is kept short here. In section 4, the infinite configuration space is mapped onto the semi-standard supertableaux. For one boundary condition, a one-to-one correspondence between the semi-standard supertableaux and the normal forms of two creation algebras is specified in section 5. Section 6 deals with the type II vertex operators related to the creation algebras. In section 7, the mixed case is investigated. Appendices A and C contain some details relegated from sections 5 and 7. Some properties of the R-matrix associated with the infinite-dimensional module and the list of explicit expressions for the R-matrix elements are given in appendix B.
## 2. The model
### 2.1. The quantum affine superalgebra $`U_q\left(\widehat{gl}(2|2)\right)`$
The integrable vertex model investigated in is based on the quantum affine superalgebra $`U_q\left(\widehat{gl}(2|2)\right)`$. Defining relations in terms of the Chevalley or Drinfeld basis as well as references regarding the representation theory of the algebra are given in . $`U_q\left(\widehat{gl}(2|2)\right)`$ is an associative $`_2`$-graded algebra over $`[[q1]]`$ with generators $`E_n^{k,\pm },k=1,2,3,n`$ and $`\mathrm{\Psi }_{\pm n}^{l,\pm },l=1,2,3,4,n_+`$, the central element $`c`$ and the grading operator $`d`$. All simple roots are chosen odd. A $`_2`$ grading $`||`$ is defined by $`|E_n^{k,\pm }|=1`$ $`k,n`$ and $`|q^c|=|d|=|\mathrm{\Psi }_{\pm n}^{l,\pm }|=0`$ $`l,n`$. The defining relations of $`U_q\left(\widehat{gl}(2|2)\right)`$ in terms of the Drinfeld basis can be found in . $`U_q\left(\widehat{sl}(2|2)\right)`$ is the subalgebra with generators $`E_n^{k,\pm },n`$ and $`\mathrm{\Psi }_{\pm n}^{k,\pm },n_+`$ with $`k=1,2,3`$, the central element $`q^c`$ and the grading operator $`d`$. $`U_q^{}\left(\widehat{gl}(2|2)\right)`$ and $`U_q^{}\left(\widehat{sl}(2|2)\right)`$ denote the superalgebras obtained by discarding the grading operator $`d`$ from the set of generators of $`U_q\left(\widehat{gl}(2|2)\right)`$ or $`U_q\left(\widehat{sl}(2|2)\right)`$, respectively.
On any $`U_q\left(\widehat{gl}(2|2)\right)`$-module considered here, $`c`$ acts as a scalar taking the value $`0`$ or $`1`$. This value is referred to as the level of the module. Parts of the subsequent analysis involve a free boson realisation of $`U_q\left(\widehat{gl}(2|2)\right)`$ applying to the level-one case. The generating functions
(1)
$$E^{k,\pm }(z)=\underset{n}{}E_n^{k,\pm }z^{n1}\mathrm{\Psi }^{l,\pm }(z)=\underset{n0}{}\mathrm{\Psi }_{\pm n}^{l,\pm }z^n$$
and the grading operator $`d`$ can be expressed in terms of six sets $`\{\phi ^l,\phi _0^l,\phi _n^l,l=1,2,3,4;n\}`$ and $`\{\beta ^{\overline{l}},\beta _0^{\overline{l}},\beta _n^{\overline{l}},\overline{l}=1,2;n\}`$ of bosonic oscillators. These oscillators satisfy the commutation relations
(2)
$$\begin{array}{cc}\hfill [\phi _n^l,\phi _m^l^{}]& =\delta _{l,l^{}}\delta _{n+m,0}\frac{[n]^2}{n}n,m0\hfill \\ \hfill [\phi ^l,\phi _0^l^{}]& =i\delta _{l,l^{}}\hfill \end{array}$$
and
(3)
$$\begin{array}{cc}\hfill [\beta _n^{\overline{l}},\beta _m^{\overline{l}^{}}]& =n\delta _{\overline{l},\overline{l}^{}}\delta _{n+m,0}n,m0\hfill \\ \hfill [\beta ^{\overline{l}},\beta _0^{\overline{l}^{}}]& =i\delta _{\overline{l},\overline{l}^{}}\hfill \end{array}$$
where $`[n]\frac{q^nq^n}{qq^1}`$. The currents $`\mathrm{\Psi }^{l,\pm }(z)`$ are realized by
(4)
$$\mathrm{\Psi }^{4,\pm }(z)=q^{\pm (\phi _0^1i\phi _0^4)}\mathrm{exp}\left(\pm (qq^1)\underset{n>0}{}\left(\phi _{\pm n}^1i\phi _{\pm n}^4\right)z^n\right)$$
and
(5)
$$\mathrm{\Psi }^{l,\pm }(z)=q^{i^l(\phi _0^{l+1}+i\phi _0^l)}\mathrm{exp}\left((qq^1)\underset{n>0}{}i^l\left(\phi _{\pm n}^{l+1}+i\phi _{\pm n}^l\right)z^n\right)$$
for $`l=1,2,3`$. To express the remaining generating functions, it is convenient to introduce the deformed free fields
(6)
$$\phi ^{l,\pm }(z)=\phi ^li\phi _0^l\mathrm{ln}z+i\underset{n0}{}\frac{q^{\frac{1}{2}|n|}}{[n]}\phi _n^lz^n$$
and
(7)
$$\beta ^{\overline{l}}(z)=\beta ^{\overline{l}}i\beta _0^{\overline{l}}\mathrm{ln}z+i\underset{n0}{}\frac{1}{n}\beta _n^{\overline{l}}z^n$$
Then the generating functions $`E^{k\pm }(z)`$ can be written
(8)
$$E^{k,\pm }(z)=:\mathrm{exp}(i^{k+1}(\phi ^{k+1,\pm }(z)+i\phi ^{k,\pm }(z))):\mathrm{exp}(\pm i\pi \delta _{k,1}\phi _0^3)X^{k,\pm }(z)$$
with
(9)
$$\begin{array}{cc}\hfill X^{1,}(z)=X^{2,+}(z)& =\frac{1}{z(qq^1)}(:\mathrm{exp}\left(\beta ^1(q^1z)\right)::\mathrm{exp}\left(\beta ^1(qz)\right):)\hfill \\ \hfill X^{1,+}(z)=X^{2,}(z)& =:\mathrm{exp}(\beta ^1(z)):\hfill \\ \hfill X^{3,+}(z)& =\frac{1}{z(qq^1)}(:\mathrm{exp}\left(\beta ^2(q^1z)\right)::\mathrm{exp}\left(\beta ^2(qz)\right):)\hfill \\ \hfill X^{3,}(z)& =:\mathrm{exp}(\beta ^2(z)):\hfill \end{array}$$
The grading operator $`d`$ is characterised by the properties
(10)
$$w^dE^{k,\pm }(z)w^d=wE^{k,\pm }(wz)w^d\mathrm{\Psi }^{l,\pm }(z)w^d=\mathrm{\Psi }^{l,\pm }(wz)$$
With (4), (5) and (8), this implies
(11)
$$d=\frac{1}{2}\underset{l=1}{\overset{4}{}}\left(\phi _0^l\right)^2\underset{l=1}{\overset{4}{}}\underset{n>0}{}\frac{n^2}{[n]^2}\phi _n^l\phi _n^l+\frac{1}{2}\underset{\overline{l}=1,2}{}\beta _0^{\overline{l}}(\beta _0^{\overline{l}}i)+\underset{\overline{l}=1,2}{}\underset{n>0}{}\beta _n^{\overline{l}}\beta _n^{\overline{l}}$$
Expressions (4), (5), (8) and (11) satisfy the defining relations of the Drinfeld basis of $`U_q\left(\widehat{gl}(2|2)\right)`$ with $`c`$ replaced by the scalar $`1`$. Later analysis of Fock spaces associated to the bosonic oscillators requires four further fields $`\eta ^{\overline{l}}(z)`$ and $`\xi ^{\overline{l}}(z)`$ introduced by
(12)
$$\eta ^{\overline{l}}(z)=\underset{n}{}\eta _n^{\overline{l}}z^{n1}=:e^{\beta ^{\overline{l}}(z)}:\xi ^{\overline{l}}(z)=\underset{n}{}\xi _n^{\overline{l}}z^n=:e^{\beta ^{\overline{l}}(z)}:$$
with $`\overline{l}=1,2`$. The relations (3) imply $`\xi _n^{\overline{l}}\eta _m^{\overline{l}^{}}+\eta _m^{\overline{l}^{}}\xi _n^{\overline{l}}=\delta _{\overline{l},\overline{l}^{}}\delta _{n+m,0}`$ and $`\xi _n^{\overline{l}}\xi _m^{\overline{l}^{}}+\xi _m^{\overline{l}^{}}\xi _n^{\overline{l}}=\eta _n^{\overline{l}}\eta _m^{\overline{l}^{}}+\eta _m^{\overline{l}^{}}\eta _n^{\overline{l}}=0`$. Both $`\eta _0^1`$ and $`\eta _0^2`$ commute with all generators (4), (5) and (8). In section 6, the free boson realisation of the superalgebra will be employed to establish the commutation relations of various type-II vertex operators.
For various purposes, it is convenient to write $`\mathrm{\Psi }^{l,\pm }(z)`$ in terms the generators $`h_l`$ and $`H_n^l`$ with $`l=1,2,3,4`$ and $`n\backslash 0`$:
(13)
$$\mathrm{\Psi }^{l,\pm }(z)=q^{\pm h_l}\mathrm{exp}\left(\pm (qq^1)\underset{n>0}{}H_{\pm n}^lz^n\right)l=1,2,3,4$$
Equations (4) and (5) yield
(14) $`h_4`$ $`=\phi _0^1i\phi _0^4`$ $`h_l`$ $`=i^l\left(\phi _0^{l+1}+i\phi _0^l\right)`$
(15) $`H_n^4`$ $`=\phi _n^1i\phi _n^4`$ $`H_n^l`$ $`=i^l\left(\phi _n^{l+1}+i\phi _n^l\right)`$
with $`l=1,2,3,n0`$. The quantum superalgebras generated by $`E_0^{k,\pm }`$ and $`q^{\pm h_l}`$ with $`k,l=1,2,3`$ or $`k=1,2,3`$, $`l=1,2,3,4`$ are denoted by $`U_q\left(sl(2|2)\right)`$ or $`U_q\left(gl(2|2)\right)`$, respectively. The generators $`h_1+h_3`$ and $`H_n^1+H_n^3`$ with $`n0`$ constitute the commutative algebra $``$. All generators of $`U_q^{}\left(\widehat{sl}(2|2)\right)`$ commute with $``$.
In terms of the basis $`\{\tau _l\}_{1l4}`$ with the bilinear form $`(\tau _l,\tau _l^{})=(1)^{l+1}\delta _{l,l^{}}`$, the classical simple roots $`\overline{\alpha }_l`$ are written $`\overline{\alpha }_l=(1)^{l+1}(\tau _l+\tau _{l+1})`$ for $`l=1,2,3`$ and $`\overline{\alpha }_4=\tau _1\tau _4`$. The classical weights $`\overline{\mathrm{\Lambda }}_l`$ with $`l=1,2,3,4`$ are given by $`\overline{\mathrm{\Lambda }}_l=_{l^{}=1}^l\tau _l^{}\frac{1}{2}_{l^{}=1}^4\tau _l^{}`$. In addition, an affine root $`\delta `$ and the corresponding affine weight $`\mathrm{\Lambda }_0`$ with the properties $`(\mathrm{\Lambda }_0,\mathrm{\Lambda }_0)=(\delta ,\delta )=(\tau _l,\mathrm{\Lambda }_0)=(\tau _l,\delta )=0`$ and $`(\mathrm{\Lambda }_0,\delta )=1`$ are introduced. Then the set of simple roots is expressed by $`\alpha _0=\delta \overline{\alpha }_1\overline{\alpha }_2\overline{\alpha }_3`$ and $`\alpha _l=\overline{\alpha }_l`$ for $`l=1,2,3,4`$. The weight lattice is the free Abelian group $`P=_{l=0}^4\mathrm{\Lambda }_l+\delta `$ with $`\mathrm{\Lambda }_l=\overline{\mathrm{\Lambda }_l}+\mathrm{\Lambda }_0`$ for $`l=1,2,3`$ and $`\mathrm{\Lambda }_4=\overline{\mathrm{\Lambda }}_4`$. $`P`$ and its dual lattice $`P^{}=_{l=0}^4h_l+d`$ with $`h_0=ch_1h_2h_3`$ can be identified via the bilinear form $`(,)`$ by setting $`\alpha _l=h_l`$ and $`d=\mathrm{\Lambda }_0`$.
### 2.2. The vertex model
The vertex model is constructed from the four-dimensional $`U_q^{}\left(\widehat{sl}(2|2)\right)`$-module $`W`$ with basis $`\{w_k\}_{0k3}`$ and the dual module $`W^{}`$ with basis $`\{w_k^{}\}_{0k3}`$. Their $`U_q^{}\left(\widehat{sl}(2|2)\right)`$-structures are given by
(16) $`h_jw_k`$ $`=(1)^{j+1}(\delta _{j,k+1}+\delta _{j,k})w_k`$
$`\mathrm{\Psi }_{\pm n}^{j,\pm }w_k`$ $`=\pm (1)^{j+1}(qq^1)q^{\pm n(2\delta _{j,2})}\left(\delta _{j,k+1}+\delta _{j,k}\right)w_k`$ $`n>0`$
$`h_jw_k^{}`$ $`=(1)^j(\delta _{j,k+1}+\delta _{j,k})w_k^{}`$
$`\mathrm{\Psi }_{\pm n}^{j,\pm }w_k^{}`$ $`=\pm (1)^j(qq^1)q^{\pm n\delta _{j,2}}\left(\delta _{j,k+1}+\delta _{j,k}\right)w_k^{}`$ $`n>0`$
with $`j=1,2,3`$ and
(17)
$$\begin{array}{cc}\hfill E_n^{j,+}w_j& =(1)^{j+1}q^{n(2\delta _{j,2})}w_{j1}\hfill \\ \hfill E_n^{j,}w_{j1}& =q^{n(2\delta _{j,2})}w_j\hfill \\ \hfill E_n^{j,+}w_{j1}^{}& =q^{(n+2)\delta _{j,2}1}w_j^{}\hfill \\ \hfill E_n^{j,}w_j^{}& =(1)^{j+1}q^{(n2)\delta _{j,2}+1}w_{j1}^{}\hfill \end{array}$$
$`n`$. Extension of (17) and (16) to $`U_q^{}\left(\widehat{gl}(2|2)\right)`$-structures is not unique. A convenient choice is
(18) $`h_4w_k`$ $`=(\delta _{k,0}\delta _{k,3})w_k`$
$`\mathrm{\Psi }_{\pm n}^{4,\pm }w_k`$ $`=\pm (qq^1)\left(q^{\pm 3n}\delta _{k,0}q^{\pm n}\delta _{k,3}\right)w_k`$ $`n>0`$
$`h_4w_k^{}`$ $`=(\delta _{k,0}\delta _{k,3})w_k^{}`$
$`\mathrm{\Psi }_{\pm n}^{4,\pm }w_k^{}`$ $`=(qq^1)\left(q^n\delta _{k,0}q^{\pm n}\delta _{k,3}\right)w_k^{}`$ $`n>0`$
The modules $`W`$ and $`W^{}`$ are attributed alternately to the horizontal as well as the vertical lines of the lattice composing the integrable vertex model. To each of the four types of elementary vertices, a spectral parameter $`z`$ or $`(q^2w)^{\pm 1}z`$ is assigned as shown in figure 1. Boltzmann weights are assigned to the elementary vertices depending on the spectral parameter and on the configuration of basis elements $`\{w_k\}_{0k3}`$ or $`\{w_k^{}\}_{0k3}`$ on the joining links. Their values follow from the R-matrices intertwining the action of $`U_q\left(\widehat{gl}(2|2)\right)`$ on the tensor products of two evaluation modules (72) (see ). Four neighbouring vertices may be viewed as a composite vertex of type A or B as illustrated in figure 1. In the limit $`q0`$, the Boltzmann weights of a composite vertex provide a well-defined map $`(WW^{})^2(WW^{})^2`$ (type A) or $`(W^{}W)^2(W^{}W)^2`$ (type B). Due to the particular spectral inhomogeneity chosen, the maps are invertible. These properties allow to consider the Hamiltonian of a corner transfer matrix at $`q=0`$. A northwest corner transfer matrix built from composite vertices of type $`A`$ or type $`B`$ acts on the half-infinite configurations $`(\mathrm{}v_{k_4}v_{k_3}v_{k_2}v_{k_1})`$ subject to a suitable boundary condition. These configurations are specified by $`v_{2r}=w_{2r}`$, $`v_{2r1}=w_{2r1}^{}`$ for type $`A`$ and by $`v_{2r}=w_{2r}^{}`$, $`v_{2r1}=w_{2r1}`$ for type $`B`$. The space of all configurations $`(\mathrm{}v_{k_4}v_{k_3}v_{k_2}v_{k_1})`$ with $`k_r=k`$ for almost all $`r`$ may be called $`\mathrm{\Omega }_A^{(k)}`$ in the type A case and $`\mathrm{\Omega }_B^{(k)}`$ in the type B case. In the following, the values $`k=1`$ and $`k=3`$ are considered. Analysis and structure of the results for $`\mathrm{\Omega }_A^{(1)}`$ and $`\mathrm{\Omega }_B^{(1)}`$ are similar to the procedure and findings described in for the spaces $`\mathrm{\Omega }_A^{(3)}`$ and $`\mathrm{\Omega }_B^{(3)}`$. The action of the CTM Hamiltonians becomes triangular in the limit $`q0`$. Moreover, in this limit the diagonal elements decouple into a part depending only on the labels $`k_{2r}`$ and another part determined by the entries $`k_{2r1}`$ only.
Denoting the diagonal element for the configuration $`(\mathrm{}w_{k_4}w_{k_3}^{}w_{k_2}w_{k_1}^{})\mathrm{\Omega }_A^{(k)}`$ by $`h_{(\mathrm{},k_4,k_3,k_2,k_1);(\mathrm{},k_4,k_3,k_2,k_1)}`$, the two contributions are specified by
(19)
$$h_{(\mathrm{},k_4,k_3,k_2,k_1);(\mathrm{},k_4,k_3,k_2,k_1)}=\underset{r=1}{\overset{\mathrm{}}{}}r\left(x_{k_{2r+1},k_{2r1}}+y_{k_{2r+2},k_{2r}}\right)$$
with
(20)
$$x_{k_1,k_2}=y_{k_2,k_1}=\{\begin{array}{cc}0\hfill & \text{if }k_1>k_2\text{ or }k_1=k_2=1,3;\hfill \\ 1\hfill & \text{if }k_1<k_2\text{ or }k_1=k_2=0,2.\hfill \end{array}$$
The diagonal element for the configuration $`(\mathrm{}w_{k_4}^{}w_{k_3}w_{k_2}^{}w_{k_1})\mathrm{\Omega }_B^{(k)}`$ is given by
(21)
$$\underset{r=1}{\overset{\mathrm{}}{}}r\left(x_{k_{2r+2},k_{2r}}+y_{k_{2r+1},k_{2r1}}\right)$$
In view of the triangular action of the corner transfer matrix Hamiltonians, $`x_{k_1,k_2}`$ and $`y_{k_1,k_2}`$ may be called the quasi-energy functions of the vertex model.
A $`U_q\left(gl(2|2)\right)`$-weight for a configuration $`(\mathrm{}w_{k_4}w_{k_3}^{}w_{k_2}w_{k_1}^{})\mathrm{\Omega }_A^{(k)}`$ or $`(\mathrm{}w_{k_4}^{}w_{k_3}w_{k_2}^{}w_{k_1})\mathrm{\Omega }_B^{(k)}`$ follows unambiguously from a $`U_q\left(gl(2|2)\right)`$-reference weight $`(\overline{h}_1^{ref},\overline{h}_2^{ref},\overline{h}_3^{ref},\overline{h}_4^{ref})`$ introduced for the configuration $`(\mathrm{}w_kw_k^{}w_kw_k^{})`$ or $`(\mathrm{}w_k^{}w_kw_k^{}w_k)`$, respectively. For an arbitrary configuration in $`\mathrm{\Omega }_A^{(k)}`$, the $`U_q\left(gl(2|2)\right)`$-weight is given by
(22)
$$\begin{array}{c}h_l(\mathrm{}w_{k_4}w_{k_3}^{}w_{k_2}w_{k_1}^{})=\hfill \\ \hfill \left\{\overline{h}_l^{ref}+\overline{h}_l^{A,k}(\mathrm{},k_4,k_3,k_2,k_1)\right\}(\mathrm{}w_{k_4}w_{k_3}^{}w_{k_2}w_{k_1}^{})\end{array}$$
with
(23)
$$\begin{array}{cc}\hfill \overline{h}_1^{A,k}(\mathrm{},k_4,k_3,k_2,k_1)=& \overline{h}_3^{A,k}(\mathrm{},k_4,k_3,k_2,k_1)=\hfill \\ & \delta _{k,1}\underset{r=1}{\overset{\mathrm{}}{}}\left(\delta _{k_{2r},2}+\delta _{k_{2r},3}\delta _{k_{2r1},2}\delta _{k_{2r1},3}\right)\hfill \\ & +\delta _{k,3}\underset{r=1}{\overset{\mathrm{}}{}}\left(\delta _{k_{2r},0}+\delta _{k_{2r},1}\delta _{k_{2r1},0}\delta _{k_{2r1},1}\right)\hfill \\ \hfill \overline{h}_2^{A,k}(\mathrm{},k_4,k_3,k_2,k_1)=& \delta _{k,1}\underset{r=1}{\overset{\mathrm{}}{}}\left(\delta _{k_{2r},0}+\delta _{k_{2r},3}\delta _{k_{2r1},0}\delta _{k_{2r1},3}\right)\hfill \\ & \delta _{k,3}\underset{r=1}{\overset{\mathrm{}}{}}\left(\delta _{k_{2r},1}+\delta _{k_{2r},2}\delta _{k_{2r1},1}\delta _{k_{2r1},2}\right)\hfill \\ \hfill \overline{h}_4^{A,k}(\mathrm{},k_4,k_3,k_2,k_1)=& \delta _{k,1}\underset{r=1}{\overset{\mathrm{}}{}}\left(\delta _{k_{2r},0}\delta _{k_{2r},3}\delta _{k_{2r1},0}+\delta _{k_{2r1},3}\right)\hfill \\ \hfill +\delta _{k,3}\underset{r=1}{\overset{\mathrm{}}{}}& \left(2\delta _{k_{2r},0}+\delta _{k_{2r},1}+\delta _{k_{2r},2}2\delta _{k_{2r1},0}\delta _{k_{2r1},1}\delta _{k_{2r1},2}\right)\hfill \end{array}$$
For any configuration $`(\mathrm{}w_{k_4}^{}w_{k_3}w_{k_2}^{}w_{k_1})\mathrm{\Omega }_B^{(k)}`$, the weight reads
(24)
$$\begin{array}{c}h_l(\mathrm{}w_{k_4}^{}w_{k_3}w_{k_2}^{}w_{k_1})=\hfill \\ \hfill \left\{\overline{h}_l^{ref}\overline{h}_l^{A,k}(\mathrm{},k_4,k_3,k_2,k_1)\right\}(\mathrm{}w_{k_4}^{}w_{k_3}w_{k_2}^{}w_{k_1})\end{array}$$
A suitable choice of reference weights for the spaces $`\mathrm{\Omega }_A^{(k)}`$ and $`\mathrm{\Omega }_B^{(k)}`$ is
(25) $`h_l^{ref}`$ $`=\delta _{l,2}+\delta _{l,4}`$ $`\text{for}k=1,`$
$`h_l^{ref}`$ $`=0l`$ $`\text{for}k=3.`$
Another reference weight for $`\mathrm{\Omega }_A^{(k)}`$ is provided by adding the $`U_q\left(gl(2|2)\right)`$-weight of $`w_k^{}`$ to the weight (25):
(26) $`h_l^{ref}`$ $`=\delta _{l,4}+2\delta _{l,2}\delta _{l,1}`$ $`\text{for}k=1,`$
(27) $`h_l^{ref}`$ $`=\delta _{l,4}\delta _{l,3}`$ $`\text{for}k=3.`$
The weight assignments (22) with (25) and (26) will be called $`\overline{h}^{A,k}`$ and $`\overline{h}^{A,k}`$, respectively. Similarly, adding the $`U_q\left(gl(2|2)\right)`$-weight of $`w_k`$ to (25) gives a second reference weight for $`\mathrm{\Omega }_B^{(k)}`$:
(28)
$$\begin{array}{cc}\hfill h_l^{ref}& =\delta _{l,1}+\delta _{l,4}\text{for}k=1,\hfill \\ \hfill h_l^{ref}& =\delta _{l,3}\delta _{l,4}\text{for}k=3.\hfill \end{array}$$
The assignments (24) with (25) and (28) are denoted by $`\overline{h}^{B,k}`$ and $`\overline{h}^{B,k}`$. In the following section, a further reference weight
(29) $`h_l^{ref}`$ $`=s\delta _{l,1}+(1+s)\delta _{l,2}(1+2s^{})\delta _{l,4}`$ $`\text{for}k=1,`$
$`h_l^{ref}`$ $`=s\delta _{l,3}+s^{}\delta _{l,4}`$ $`\text{for}k=3`$
with arbitrary $`s^{}`$ and $`s`$ will be taken into account both for $`\mathrm{\Omega }_A^{(1)}`$ and $`\mathrm{\Omega }_B^{(1)}`$. The assignments (24) with (29) will be referred to as $`\overline{h}^{C,k}`$.
In , the half-infinite configurations of the spaces $`\mathrm{\Omega }_A^{(3)}`$ and $`\mathrm{\Omega }_B^{(3)}`$ are compared to the weight vectors of reducible level-one modules of $`U_q\left(\widehat{sl}(2|2)\right)/`$ denoted by $`\stackrel{~}{V}(\mathrm{\Lambda }_0)`$, $`\stackrel{~}{V}(\mathrm{\Lambda }_1+\mathrm{\Lambda }_4)`$ and $`\stackrel{~}{V}(2\mathrm{\Lambda }_0\mathrm{\Lambda }_3+\mathrm{\Lambda }_4)`$. The eigenvalues of the grading operator $`d`$ and the generators $`h_l`$, $`l=1,2,3,4,`$ are called the grade and the $`U_q\left(gl(2|2)\right)`$-weight of a weight vector. For a fixed weight assignment $`\overline{h}^A,\overline{h}^A,\overline{h}^B`$ or $`\overline{h}^B`$, the $`U_q\left(gl(2|2)\right)`$-weights of all configurations with the diagonal element of the CTM Hamiltonian given by $`n`$ are collected. At $`n=0,1,2,3`$, they are in one-to-one correspondence with the $`U_q\left(gl(2|2)\right)`$-weights of all vectors with grade $`n`$ found in one of the reducible level-one modules. This correspondence may be assumed to hold true at any grade. Then the character of the reducible module can be expressed in terms of the quasi-energy functions. Due to the decomposition of the diagonal elements (19), (21), the character expressions factorise into two parts, each of them depending only on one quasi-energy function. Each of the reducible module can be decomposed into two irreducible level-one modules. One them is weakly integrable , the other nonintegrable. Table 1 specifies the relevant level-one modules for the four choices of composite vertices and assignment of weights. The weakly integrable irreducible module is listed left of the nonintegrable module.
A similar analysis suggests that the weight states of the irreducible, nonintegrable level-one module $`V\left((1s)\mathrm{\Lambda }_0+s\mathrm{\Lambda }_3+s^{}\mathrm{\Lambda }_4\right)`$-module of $`U_q\left(\widehat{sl}(2|2)\right)/`$ correspond to the configurations if the reference weight $`\overline{h}^{C,3}`$ is adopted.
## 3. The second boundary condition
The configurations in $`\mathrm{\Omega }_A^{(1)}`$ and $`\mathrm{\Omega }_B^{(1)}`$ can be related to the weight vectors of level-one modules of $`U_q\left(\widehat{sl}(2|2)\right)/`$. An assignment of weights $`\overline{h}^{A,1}`$, $`\overline{h}^{A,1}`$, $`\overline{h}^{B,1}`$, $`\overline{h}^{B,1}`$ or $`\overline{h}^{C,1}`$ introduced in the previous section is fixed. Then the $`U_q\left(gl(2|2)\right)`$-weights of all configurations with a fixed value $`n`$ of the diagonal elements (19) or (21) are compared to the $`U_q\left(gl(2|2)\right)`$-weights of all vectors with the associated grade present in a suitable level-one module. The value of the associated grade may equal the diagonal element or differ from it by a constant value. Consideration of the three lowest values $`n=0,1,2`$ indicates the appropriate modules. Three reducible level-one modules of $`U_q\left(\widehat{sl}(2|2)\right)/`$ denoted by $`\stackrel{̊}{V}(\mathrm{\Lambda }_0)`$, $`\stackrel{̊}{V}(2\mathrm{\Lambda }_0\mathrm{\Lambda }_3+\mathrm{\Lambda }_4)`$ and $`\stackrel{̊}{V}(\mathrm{\Lambda }_1+\mathrm{\Lambda }_4)`$ account for the first four choices of composite vertices and $`U_q\left(gl(2|2)\right)`$-weights. Similar as the reducible modules related to $`\mathrm{\Omega }_A^{(3)}`$ and $`\mathrm{\Omega }_B^{(3)}`$, each of them decomposes into an irreducible weakly integrable and an irreducible nonintegrable module. In table 2, the level-one modules related to each case are listed. There the weakly integrable module appears left of the nonintegrable irreducible module. The irreducible, nonintegrable level-one module $`V\left(s\mathrm{\Lambda }_1+(s+1)\mathrm{\Lambda }_2(1+2s^{})\mathrm{\Lambda }_4\right)`$ of $`U_q\left(\widehat{sl}(2|2)\right)/`$ accounts for the choice $`\overline{h}^{C,1}`$.
A pair of irreducible modules can be assembled in a reducible module in various ways. In each case found in table 2, a particular reducible module is distinguished by a simple free boson realization in terms of the scheme given in section 2.2. These modules are obtained from the vectors
(30)
$$\kappa _0=e^{\beta ^1}|0\kappa _1=e^{i\phi ^1+\beta ^1}|0\kappa _3=\beta _1^2e^{\phi ^4+\beta ^1+\beta ^2}|0$$
with the boson Fock vacuum $`|0`$ characterised by the properties
(31)
$$\phi _0^l|0=\beta _0^{\overline{l}}|0=0\phi _n^l|0=\beta _n^{\overline{l}}|0=0n>0,l=1,2,3,4,\overline{l}=1,2$$
For the irreducible module $`V\left(s\mathrm{\Lambda }_1+(1+s)\mathrm{\Lambda }_2(1+2s^{})\mathrm{\Lambda }_4\right)`$, the corresponding vector reads
(32)
$$\kappa _4=e^{s^{}i\phi ^1(s^{}s)\phi ^2+(1+s^{})(i\phi ^3+\phi ^4)s\beta ^1}|0$$
The action of the grading operator on $`\kappa _I`$ reads $`d\kappa _I=\frac{1}{2}(\delta _{I,0}1)\kappa _I`$ for $`I=0,1,3`$ and $`d\kappa _4=\frac{1}{2}s(1+2s^{})\kappa _4`$. Arbitrary polynomials of the $`U_q\left(\widehat{gl}(2|2)\right)`$-generators (4),(5) and (8) applied on $`\kappa _I`$ give rise to vectors with the maximal value of their grades given by $`\frac{1}{2}(\delta _{I,0}+2\delta _{I,3}1)`$ for $`I=0,1,3`$ and by $`\frac{1}{2}s(1+2s^{})`$ for $`I=4`$. If $`I=0`$, the vectors with grade $`0`$ form a reducible $`U_q\left(gl(2|2)\right)`$-module which decomposes into the infinite-dimensional, irreducible module $`V(\overline{\mathrm{\Lambda }}_2+\overline{\mathrm{\Lambda }}_4)`$ and the one-dimensional module with weight $`(0,0,0,0)`$. For $`I=1`$, the vectors with grade $`\frac{1}{2}`$ furnish a reducible $`U_q\left(gl(2|2)\right)`$-module which is decomposed into the four-dimensional module $`V(\overline{\mathrm{\Lambda }}_1+\overline{\mathrm{\Lambda }}_4)`$ and the infinite-dimensional module $`V(\overline{\mathrm{\Lambda }}_3+\overline{\mathrm{\Lambda }}_4)`$. The latter are both irreducible. In the case $`I=3`$, the vectors with grade $`\frac{1}{2}`$ form the infinite-dimensional, irreducible $`U_q\left(gl(2|2)\right)`$-module $`V(\overline{\mathrm{\Lambda }}_1+2\overline{\mathrm{\Lambda }}_2+\overline{\mathrm{\Lambda }}_4)`$. Similarly, for $`I=4`$, the vectors with grade $`\frac{1}{2}s(1+2s^{})`$ constitute the infinite-dimensional, irreducible $`U_q\left(gl(2|2)\right)`$-module $`V\left(s\overline{\mathrm{\Lambda }}_1+(1+s)\overline{\mathrm{\Lambda }}_2(1+2s^{})\overline{\mathrm{\Lambda }}_4\right)`$. The $`U_q\left(gl(2|2)\right)`$-modules obtained in this way can be viewed as the maximal-grade subspaces of level-one modules of $`U_q\left(\widehat{sl}(2|2)\right)/`$. This suggests to examine the Fock spaces $`\stackrel{̊}{}_I`$ with $`I=0,1,3,4`$ defined by
(33)
$$\begin{array}{c}\stackrel{̊}{}_I=[\stackrel{~}{\phi }_1^{\overline{l}_1},\beta _1^{\overline{l}_2},\stackrel{~}{\phi }_2^{\overline{l}_3},\beta _2^{\overline{l}_4},\mathrm{}]\hfill \\ \hfill \left(_{S_1,S_2,S_3}e^{S_1(i\phi ^1+\phi ^2\beta ^1)+S_2(\phi ^2i\phi ^3\beta ^1)+S_3(i\phi ^3+\phi ^4\beta ^2)+\beta ^1+\varpi _I}|0\right)\end{array}$$
with $`\overline{l}_i=1,2`$,
(34)
$$\begin{array}{cc}\hfill \stackrel{~}{\phi }_n^1& \phi _n^1i\phi _n^2+\phi _n^3i\phi _n^4\hfill \\ \hfill \stackrel{~}{\phi }_n^2& i\phi _n^2+\phi _n^3\hfill \end{array}$$
and
(35)
$$\begin{array}{cc}\hfill \varpi _0& =0\varpi _1=\phi ^2+\beta ^1\varpi _3=\phi ^2\beta ^1\hfill \\ \hfill \varpi _4& =s(\phi ^2\beta ^1)(1+s^{})\left(i\phi ^1+\phi ^2i\phi ^3\phi ^4\right)\hfill \end{array}$$
The remaining two linear combinations in $`\phi _n^l`$ are not included in (33) since the Fock space $`\stackrel{̊}{}_I`$ should accommodate a module of $`U_q\left(\widehat{sl}(2|2)\right)/`$ rather than a $`U_q\left(\widehat{gl}(2|2)\right)`$-module. Investigation of the next lower grades indicates that the required reducible modules are realized as restricted Fock spaces:
(36)
$$\begin{array}{cc}\hfill \stackrel{̊}{V}(\mathrm{\Lambda }_0)=Ker_{\eta _0^2}\stackrel{̊}{}_0\stackrel{̊}{V}(\mathrm{\Lambda }_1+\mathrm{\Lambda }_4)& =Ker_{\eta _0^2}\stackrel{̊}{}_1\stackrel{̊}{V}(2\mathrm{\Lambda }_0\mathrm{\Lambda }_3+\mathrm{\Lambda }_4)=Ker_{\eta _0^2}\stackrel{̊}{}_3\hfill \\ \hfill V\left(s\mathrm{\Lambda }_1+(1+s)\mathrm{\Lambda }_2(1+2s^{})\mathrm{\Lambda }_4\right)& =Ker_{\eta _0^2}\stackrel{̊}{}_4\hfill \end{array}$$
The expression (11) for the grading operator $`d`$ and the properties (31) yield
(37)
$$\begin{array}{c}de^{S_1(i\phi ^1+\phi ^2\beta ^1)+S_2(\phi ^2i\phi ^3\beta ^1)+S_3(i\phi ^3+\phi ^4\beta ^2)+\beta ^1+\varpi _I}|0=\hfill \\ \hfill \frac{1}{2}(S_1(S_11)+(S_2S_3)(S_2S_31)+\delta _{I,1}\delta _{I,3}+s(1+2s^{})\delta _{I,4})\\ \hfill e^{S_1(i\phi ^1+\phi ^2\beta ^1)+S_2(\phi ^2i\phi ^3\beta ^1)+S_3(i\phi ^3+\phi ^4\beta ^2)+\beta ^1+\varpi _I}|0\end{array}$$
for $`I=0,1,3,4`$. With (14) and (37), the $`U_q\left(gl(2|2)\right)`$-weights of all vectors in $`Ker_{\eta _0^2}\stackrel{̊}{}_I`$ with grade $`n`$ are easily collected for small $`n`$. Then the one-to-one correspondence between the $`U_q\left(gl(2|2)\right)`$-weights $`\overline{h}^{A,1}`$, $`\overline{h}^{A,1}`$, $`\overline{h}^{B,1}`$, $`\overline{h}^{B,1}`$ or $`\overline{h}^{C,1}`$ of the configurations in $`\mathrm{\Omega }_A^{(1)}`$ or $`\mathrm{\Omega }_B^{(1)}`$ and the weights of the vectors in the modules (36) is readily verified for $`n=0,1,2`$.
In the cases $`I=0,1,3`$, the direct sum decomposition $`Ker_{\eta _0^2}\stackrel{̊}{}_I=\xi _0^1\eta _0^1Ker_{\eta _0^2}\stackrel{̊}{}_I\eta _0^1\xi _0^1Ker_{\eta _0^2}\stackrel{̊}{}_I`$ allows to separate the irreducible components of $`Ker_{\eta _0^2}\stackrel{̊}{}_I`$. Boson realizations for the irreducible $`U_q\left(\widehat{sl}(2|2)\right)/`$-modules in table 2 are provided by
(38) $`V(\mathrm{\Lambda }_0)`$ $`=\eta _0^1Ker_{\eta _0^2}\stackrel{̊}{}_0`$ $`V(\mathrm{\Lambda }_2+\mathrm{\Lambda }_4)`$ $`=Ker_{\eta _0^1}Ker_{\eta _0^2}\stackrel{̊}{}_0`$
$`V(\mathrm{\Lambda }_1+\mathrm{\Lambda }_4)`$ $`=\eta _0^1Ker_{\eta _0^2}\stackrel{̊}{}_1`$ $`V(\mathrm{\Lambda }_3+\mathrm{\Lambda }_4)`$ $`=Ker_{\eta _0^1}Ker_{\eta _0^2}\stackrel{̊}{}_1`$
$`V(2\mathrm{\Lambda }_0\mathrm{\Lambda }_3+\mathrm{\Lambda }_4)`$ $`=\eta _0^1Ker_{\eta _0^2}\stackrel{̊}{}_3`$ $`V(\mathrm{\Lambda }_1+2\mathrm{\Lambda }_2+\mathrm{\Lambda }_4)`$ $`=Ker_{\eta _0^1}Ker_{\eta _0^2}\stackrel{̊}{}_3`$
Hence, the zero mode $`\eta _0^1`$ annihilates the nonintegrable components. Some features of the reducible modules can be read from (30) and (38). Since $`\eta _0^1\kappa _I0I`$, the vectors $`\kappa _I`$ belong to the weakly integrable irreducible submodules. The coupling of both irreducible submodules is described by
(39) $`E_0^{1,+}\kappa _I`$ $`0`$ $`\eta _0^1E_0^{1,+}\kappa _I=E_0^{1,+}\eta _0^1\kappa _I=0`$ $`I=0,3`$
$`E_0^{2,}\kappa _J`$ $`0`$ $`\eta _0^1E_0^{2,}\kappa _J=E_0^{2,}\eta _0^1\kappa _J=0`$ $`J=0,1`$
Thus $`E_0^{1,+}\kappa _0`$, $`E_0^{2,}\kappa _1`$ and $`E_0^{1,+}\kappa _3`$ are contained in the nonintegrable submodules. A vector with maximal grade in $`V(2\mathrm{\Lambda }_0\mathrm{\Lambda }_3+\mathrm{\Lambda }_4)`$ is given by $`E_1^{2,}E_0^{3,}\kappa _3`$, for example. In the following, the notations $`\stackrel{̊}{V}(\mathrm{\Lambda }_0)`$, $`\stackrel{̊}{V}(\mathrm{\Lambda }_1+\mathrm{\Lambda }_4)`$ and $`\stackrel{̊}{V}(2\mathrm{\Lambda }_0\mathrm{\Lambda }_3+\mathrm{\Lambda }_4)`$ refer to the reducible modules realized by (36).
Provided that the correspondence between half-infinite configurations and weight states found for $`n=0,1,2`$ remains valid for all $`n0`$, the characters of the reducible modules can be written in terms of the energy functions (20). For the nonintegrable modules described above, well defined characters are introduced by
(40)
$$ch_{V()}(\varrho ,\rho _0,\rho _2,\rho _3)tr_{V()}\varrho ^d\rho _0^{\frac{1}{2}(h_1+h_2+h_3+h_4)}\rho _2^{\frac{1}{2}(h_1+h_2h_3h_4)}\rho _3^{\frac{1}{2}(h_1+h_2+h_3h_4)}$$
where $`|\rho _3|<1`$ and $`V()`$ denotes any of the modules listed in (36).
###### Conjecture 1.
The characters of the $`U_q\left(\widehat{sl}(2|2)\right)/`$-modules $`\stackrel{̊}{V}(\mathrm{\Lambda }_0)`$, $`\stackrel{̊}{V}(\mathrm{\Lambda }_1+\mathrm{\Lambda }_4)`$, $`\stackrel{̊}{V}(2\mathrm{\Lambda }_0\mathrm{\Lambda }_3+\mathrm{\Lambda }_4)`$ and $`V\left(s\mathrm{\Lambda }_1+(1+s)\mathrm{\Lambda }_2(1+2s^{})\mathrm{\Lambda }_4\right)`$ are given by
(41)
$$\begin{array}{cc}\hfill ch_{\stackrel{̊}{V}(\mathrm{\Lambda }_0)}(\varrho ,\rho _0,\rho _2,\rho _3)& =\underset{\{\mathrm{},k_6,k_4,k_2\}}{}\varrho ^{_{r>0}ry_{k_{2r+2},k_{2r}}}\underset{l=0,2,3}{}\rho _l^{_{r>0}\delta _{k_{2r},l}}\hfill \\ & \underset{\{\mathrm{},k_5,k_3,k_1\}}{}\varrho ^{_{\overline{r}>0}\overline{r}x_{k_{2\overline{r}+1},k_{2\overline{r}1}}}\underset{\overline{l}=0,2,3}{}\rho _{\overline{l}}^{_{r>0}\delta _{k_{2\overline{r}1},\overline{l}}}\hfill \end{array}$$
and
(42)
$$\begin{array}{cc}\hfill ch_{\stackrel{̊}{V}(\mathrm{\Lambda }_1+\mathrm{\Lambda }_4)}(\varrho ,\rho _0,\rho _2,\rho _3)& =\varrho ^{\frac{1}{2}}ch_{\stackrel{̊}{V}(\mathrm{\Lambda }_0)}(\varrho ,\rho _0,\rho _2,\rho _3)\hfill \\ \hfill ch_{\stackrel{̊}{V}(2\mathrm{\Lambda }_2\mathrm{\Lambda }_3+\mathrm{\Lambda }_4)}(\varrho ,\rho _0,\rho _2,\rho _3)& =\varrho ^{\frac{1}{2}}ch_{\stackrel{̊}{V}(\mathrm{\Lambda }_0)}(\varrho ,\rho _0,\rho _2,\rho _3)\hfill \\ \hfill ch_{\stackrel{̊}{V}\left(s\mathrm{\Lambda }_1+(1+s)\mathrm{\Lambda }_2(1+2s^{})\mathrm{\Lambda }_4\right)}(\varrho ,\rho _0,\rho _2,\rho _3)& =\varrho ^{\frac{1}{2}s(1+2s^{})}ch_{\stackrel{̊}{V}(\mathrm{\Lambda }_0)}(\varrho ,\rho _0,\rho _2,\rho _3)\hfill \end{array}$$
with $`|\rho _3|<1`$. The sums in (41) are restricted by the requirement that $`k_r=1`$ for almost all $`r>0`$.
In most of the remainder, the assignment $`\overline{h}^{A,1}`$ is considered.
## 4. Infinite border strips
Within the framework proposed in , , , the eigenstates of the row-to-row transfer matrix of the vertex model are created on the ground state by means of type II vertex operators. This general feature can be expected to apply to the $`U_q\left(\widehat{gl}(2|2)\right)`$-model as well. Unfortunately, there is no simple method to single out the appropriate vertex operators among all those existing at the given level. In and , the space of states of the $`XXZ`$-model and its higher spin generalisation has been decomposed at $`q=0`$. The analysis involves a creation algebra whose defining relations can be viewed as formal $`q0`$ limits of the commutation relations satisfied by the appropriate type II vertex operators. A relation between the space of states at $`q=0`$ and the creation algebra is established by means of the domain wall picture and the crystal theory associated with the paths. The resulting expressions in terms of the generators of the creation algebra are interpreted as the $`q0`$-limits of the $`n`$-particle eigenstates of the model. In the next section, two creation algebras relevant to the present model will be considered. Infinite border strips prove a useful tool for setting up the relation between their generators and the configuration space.
All following considerations refer to infinite configurations $`(\mathrm{}w_{j_2}w_{j_1}^{}w_{j_0}w_{j_1}^{}w_{j_2}w_{j_3}^{}\mathrm{})`$ subject to suitable boundary conditions. The set of all infinite configurations satisfying $`k_r=kr>r_+>0`$ and $`k_r=k^{}r<r_{}<0`$ with $`k,k^{}=1,3`$ is denoted by $`𝒦_{k,k^{}}`$. Where convenient, the infinite components $`(\mathrm{}w_{j_2}w_{j_0}w_{j_2}\mathrm{})`$ and $`(\mathrm{}w_{j_1}^{}w_{j_1}^{}w_{j_3}^{}\mathrm{})`$ will be abbreviated by $`(\mathrm{},j_2,j_0,j_2,\mathrm{})`$ and $`(\mathrm{},j_1,j_1,j_3,\mathrm{})`$, respectively. All four choices of $`k,k^{}`$ allow for a well-defined generalisation of the expression (19):
(43)
$$h_{(\mathrm{},k_2,k_0,k_2,k_4,\mathrm{});(\mathrm{},k_2,k_0,k_2,k_4,\mathrm{})}=\underset{r}{}ry_{k_{2r+2},k_{2r}}$$
(44)
$$h_{(\mathrm{},k_3,k_1,k_1,k_3,\mathrm{});(\mathrm{},k_3,k_1,k_1,k_3,\mathrm{})}=\underset{r}{}rx_{k_{2r+1},k_{2r1}}$$
with $`x_{k_{2r+1},k_{2r1}}`$ and $`y_{k_{2r+2},k_{2r}}`$ defined by (20). $`U_q\left(gl(2|2)\right)`$-weights $`(\overline{h}_1,\overline{h}_2,\overline{h}_3,\overline{h}_4)`$ compatible with the assignment $`\overline{h}^{A,1}`$ defined by (22) and (25) are introduced for the infinite components of $`𝒦_{1,1}`$ by
(45)
$$\begin{array}{cc}\hfill \overline{h}_1(\mathrm{},k_2,k_0,k_2,k_4,\mathrm{})=& \overline{h}_3(\mathrm{},k_2,k_0,k_2,k_4,\mathrm{})\hfill \\ & =\underset{r}{}(\delta _{k_{2r},2}+\delta _{k_{2r},3})(\mathrm{},k_2,k_0,k_2,k_4,\mathrm{})\hfill \\ \hfill \overline{h}_2(\mathrm{},k_2,k_0,k_2,k_4,\mathrm{})=& \underset{r}{}(\delta _{k_{2r},0}+\delta _{k_{2r},3})(\mathrm{},k_2,k_0,k_2,k_4,\mathrm{})\hfill \\ \hfill \overline{h}_4(\mathrm{},k_2,k_0,k_2,k_4,\mathrm{})=& \underset{r}{}(\delta _{k_{2r},0}\delta _{k_{2r},3})(\mathrm{},k_2,k_0,k_2,k_4,\mathrm{})\hfill \end{array}$$
and
(46)
$$\begin{array}{cc}\hfill \overline{h}_1(\mathrm{},k_3,k_1,k_1,k_3,\mathrm{})=& \overline{h}_3(\mathrm{},k_3,k_1,k_1,k_3,\mathrm{})\hfill \\ & =\underset{r}{}(\delta _{k_{2r1},2}+\delta _{k_{2r1},3})(\mathrm{},k_3,k_1,k_1,k_3,\mathrm{})\hfill \\ \hfill \overline{h}_2(\mathrm{},k_3,k_1,k_1,k_3,\mathrm{})=& \underset{r}{}(\delta _{k_{2r1},0}+\delta _{k_{2r1},3})(\mathrm{},k_3,k_1,k_1,k_3,\mathrm{})\hfill \\ \hfill \overline{h}_4(\mathrm{},k_3,k_1,k_1,k_3,\mathrm{})=& \underset{r}{}(\delta _{k_{2r1},0}\delta _{k_{2r1},3})(\mathrm{},k_3,k_1,k_1,k_3,\mathrm{})\hfill \end{array}$$
For homogeneous vertex models related to quantum affine algebras, a one-to-one correspondence between the spin configurations and semi-standard super tableaux of skew Young diagrams has been demonstrated in . A similar one-to-one correspondence between the half-infinite spin configurations of the present model and two types of semi-standard super tableaux of finite and half-infinite border strips is pointed out in . Making use of this correspondence, pairs of infinite border strips can be related to the components $`(\mathrm{}w_{k_4}w_{k_2}w_{k_0}w_{k_2}w_{k_4}\mathrm{})`$ and $`(\mathrm{}w_{k_3}^{}w_{k_1}^{}w_{k_1}^{}w_{k_3}^{}w_{k_5}^{}\mathrm{})`$ of the infinite configurations $`(\mathrm{}w_{k_2}w_{k_1}^{}w_{k_0}w_{k_1}^{}w_{k_2}w_{k_3}^{}\mathrm{})`$. These border strips consist of finitely many rows and columns of finite length assembled between either two half-infinite rows or two half-infinite columns. In the following, they will be referred to as horizontal or vertical border strips, respectively. As an example, figure 2 shows a horizontal border strip. The horizontal border strips are related to the components $`(\mathrm{}w_{k_4}w_{k_2}w_{k_0}w_{k_2}w_{k_4}\mathrm{})`$ and the vertical border strips to the components $`(\mathrm{}w_{k_3}^{}w_{k_1}^{}w_{k_1}^{}w_{k_3}^{}w_{k_5}^{}\mathrm{})`$. A semi-standard super tableau of a horizontal or vertical border strip is obtained by assigning one of the numbers $`0,1,2,3`$ to each box such that the numbers of each two neighbouring boxes satisfy two rules:
1. If the side common to both boxes is vertical, then the number $`k_1`$ in the left box and the number $`k_2`$ in the right box fulfil
(47)
$$k_1>k_2\text{or}k_1=k_2=\{\begin{array}{cc}0,2\hfill & \text{vertical strip}\hfill \\ 1,3\hfill & \text{horizontal strip}\hfill \end{array}$$
2. If the side common to both boxes is horizontal, then the number $`k_1`$ in the upper box and the number $`k_2`$ in the lower box fulfil
(48)
$$k_1>k_2\text{or}k_1=k_2=\{\begin{array}{cc}1,3\hfill & \text{vertical strip}\hfill \\ 0,2\hfill & \text{horizontal strip}\hfill \end{array}$$
Almost all numbers attributed to the infinite strip are fixed by a boundary condition. The set of all semi standard super tableaux of horizontal (vertical) border strips with the number $`k`$ given to almost all boxes of the lower half-infinite row (left half-infinite column) and the number $`k^{}`$ given to almost all boxes in the upper half-infinite row (right half-infinite column) will be called $`_{k,k^{}}^h`$ ($`_{k,k^{}}^v`$). Each of the four choices $`k,k^{}=1,3`$ is consistent with the rules (47), (48). The set $`_{k,k}^h`$ ($`_{k,k}^v`$) includes exactly one semi standard super tableau for the infinite border strip consisting of one single row (column). This tableau attributes the number $`k`$ to each box. Excluding this tableau from $`_{k,k}^h`$ $`(_{k,k}^v`$) yields a set called $`_{k,k}^{h\backslash 0}`$ $`\left(_{k,k}^{v\backslash 0}\right)`$.
A horizontal border strip with at least one finite row is characterised by the set $`R,(p_1,p_2,\mathrm{},p_R;\overline{p}_1,\overline{p}_2,\mathrm{},\overline{p}_{R1})`$ with $`R>1`$ and $`p_i,\overline{p}_i`$. This border strip contains $`R`$ finite columns with more than one box. A vertical border strip with at least one finite column is described by the set $`R,(p_1,p_2,\mathrm{},p_{R1};\overline{p}_1,\overline{p}_2,\mathrm{},\overline{p}_R)`$ with $`R>1`$ and $`p_i,\overline{p}_i`$. In both cases, the parameter $`p_i+1`$ ($`\overline{p}_i+1`$) specifies the number of the boxes contained in the $`i`$-th finite column (finite row) composed of at least two boxes. Here the counting proceeds from the right to the left end of a horizontal border strip or from the upper to the lower end of a vertical border strip. In the remainder, the counting of rows or columns refers only to rows or columns consisting of more than one box. The set $`R=1,(p_1;\mathrm{})`$ with $`p_1`$ specifies a horizontal border strip composed of the lower and upper half-infinite row and a column containing $`p_1+12`$ boxes. Similarly, the set $`R=1,(\mathrm{};\overline{p}_1)`$ characterises a vertical border strip built from the left and right half-infinite columns and a row with $`\overline{p}_1+12`$ boxes. Finally, the value $`R=0`$ refers to the border strips consisting of one single infinite row or column.
A horizontal (vertical) border strip with $`R1`$ has exactly one semi-standard super tableau in $`_{1,1}^h`$ $`(_{1,1}^v`$) involving only the numbers $`0`$ and $`1`$. In a horizontal border strip, the $`p_i`$ lower boxes of each column with $`p_i+12`$ boxes receive the number $`0`$. All other boxes obtain the number $`1`$. In case of a vertical border strip, the rightmost $`\overline{p}_i`$ boxes in a row with $`\overline{p}_i+12`$ boxes receive the value $`0`$. The number $`1`$ is given to all other boxes. These tableaux will be called the reference labellings in the following. An example is illustrated by figure 2.
In order to define a mapping from a semi standard super tableau onto a component of an infinite configuration, a counting of the boxes needs to be specified. The leftmost box of the upper half-infinite row of a horizontal border strip or the lowest box in the right half-infinite column of a vertical border strip is counted as the $`r_0`$-th box for some $`r_0`$. The $`(r+1)`$-th box is left of or below the $`r`$-th box for any $`r`$. Then for a horizontal border strip with parameters $`R=0`$, $`R=1,(p_1;\mathrm{})`$ or $`R>1,(p_1,\mathrm{},p_R;\overline{p}_1,\mathrm{},\overline{p}_{R1})`$ an arbitrary semi-standard super tableau in $`_{k,k^{}}^h`$ may be considered. Given a fixed value of $`r_0`$, it is convenient to introduce numbers $`s_i`$, $`1iR`$, by
(49)
$$\begin{array}{cc}\hfill s_1& =r_0\hfill \\ \hfill s_i& =r_0+\underset{i_<=1}{\overset{i1}{}}(p_{i_<}+\overline{p}_{i_<})2iR\hfill \end{array}$$
and
(50)
$$d^h(p_1,\mathrm{},p_R;\overline{p}_1,\mathrm{},\overline{p}_{R1};r_0)=\underset{i=1}{\overset{R}{}}p_i\left(s_i+\frac{1}{2}(p_i1)\right)$$
A component $`(\mathrm{}w_{k_4}w_{k_2}w_{k_0}w_{k_2}w_{k_4}\mathrm{})`$ is associated with the semi-standard super tableau by identifying the number attributed to the $`(r+1)`$-th box with $`k_{2r}`$. With (20) and (43) it is easily verified that the contribution of this component to the diagonal element of the CTM Hamiltonian coincides with (50).
For a vertical border strip with parameters $`R^{}=0`$, $`R^{}=1,(\mathrm{};\overline{p}_1^{})`$ or $`R^{}>1,(p_1^{},\mathrm{},p_{R^{}1}^{};\overline{p}_1^{},\mathrm{},\overline{p}_R^{}^{})`$ and the box counting fixed by the number $`r_0^{}`$, the number
(51)
$$d^v(p_1^{},\mathrm{},p_{R^{}1}^{};\overline{p}_1^{},\mathrm{},\overline{p}_R^{}^{};r_0)=\underset{i=1}{\overset{R^{}}{}}\overline{p}_i^{}\left(s_i^{}+\frac{1}{2}(\overline{p}_i^{}1)\right)$$
is introduced. Here $`s_i^{}`$ is defined by (49) with $`r_0`$, $`p_{i_<}`$, $`\overline{p}_{i_<}`$ replaced by $`r_0^{}`$, $`p_{i_<}^{}`$ and $`\overline{p}_{i_<}^{}`$, respectively. A semi standard super tableau in $`_{k,k^{}}^v`$ is related to a component $`(\mathrm{}w_{k_3}^{}w_{k_1}^{}w_{k_1}^{}w_{k_3}^{}w_{k_5}^{}\mathrm{})`$ by identifying the number attributed to the $`(r+1)`$-th box with $`k_{(2r1)}`$. According to (20) and (44), the contribution of this component to the diagonal element of the CTM Hamiltonian equals (51).
Analogous statements apply to the components $`(\mathrm{}w_1w_1w_1\mathrm{})`$ and $`(\mathrm{}w_1^{}w_1^{}w_1^{}\mathrm{})`$ provided that
(52)
$$d^h(0,\mathrm{})=d^v(\mathrm{},0)=0$$
This is readily verified comparing the definitions (20) for $`y_{k_{2r+2},k_{2r}}`$ and $`x_{k_{2r+1},k_{2r1}}`$ with the rules (47), (48). A $`U_q\left(gl(2|2)\right)`$-weight is assigned to each semi standard super tableau in $`_{1,1}^h`$ or $`_{1,1}^v`$ via (46) or (45) and the above identifications.
The above prescription gives a one-to-one correspondence between the components $`(\mathrm{}w_{k_4}w_{k_2}w_{k_0}w_{k_2}w_{k_4}\mathrm{})`$ and $`(\mathrm{}w_{k_3}^{}w_{k_1}^{}w_{k_1}^{}w_{k_3}^{}w_{k_5}^{}\mathrm{})`$ of the configurations in $`𝒦_{k,k^{}}`$ and the sets $`B_{k,k^{}}^h`$ and $`B_{k,k^{}}^v`$ of semi-standard super tableaux associated with infinite border strips. This correspondence will be used in the next section to describe the configuration space in terms of two creation algebras.
## 5. The creation algebras
###### Definition 1.
The creation algebra $`𝒜^{}`$ is generated by $`\{\varphi _{j,t;m}^{}|0j3,t_0,m\}`$ over $`𝒵`$ subject to the defining relations
(53)
$$\varphi _{j_2,t_2;m_2}^{}\varphi _{j_1,t_1;m_1}^{}=\varphi _{j_2,t_2;m_1+t_2+\theta _{j_2,j_1}}^{}\varphi _{j_1,t_1;m_2t_2\theta _{j_2,j_1}}^{}$$
with
(54)
$$\theta _{j_2,j_1}=\{\begin{array}{cc}0,\hfill & \text{if }j_2=2;\hfill \\ 0,\hfill & \text{if }j_2=0,3\text{ and }j_1=0,1\text{;}\hfill \\ 1,\hfill & \text{otherwise.}\hfill \end{array}$$
A special case of the defining relations (53) is
(55)
$$\varphi _{j_2,t_2;m+t_2+\theta _{j_2,j_1}}^{}\varphi _{j_1,t_1;m}^{}=0t_1_0,m$$
For the subsequent analysis, it is useful to introduce the notion of normal forms . The product
(56)
$$\varphi _{j_n,t_n;m_n}^{}\mathrm{}\varphi _{j_2,t_2;m_2}^{}\varphi _{j_1,t_1;m_1}^{}$$
is called a normal form (of $`𝒜^{}`$) iff
(57)
$$m_{i+1}>m_i+t_{i+1}+\theta _{j_{i+1},j_i}\text{for}\mathrm{\hspace{0.33em}\hspace{0.33em}1}i<n$$
The set
(58)
$$B^{}=_n\left\{\varphi _{j_n,t_n;m_n}^{}\mathrm{}\varphi _{j_1,t_1;m_1}^{}\right|j_i=0,1,2,3,t_i_0,m_i\text{satisfy (}\text{57}\text{)}\}$$
provides a $``$-linear base of the algebra $`𝒜^{}`$. This statement is a special case of corollary 2 proven in . A $`U_q\left(gl(2|2)\right)`$-weight $`(h_1^\varphi ^{},h_2^\varphi ^{},h_3^\varphi ^{},h_4^\varphi ^{})`$ is introduced for the generators $`\varphi _{j,t;m}^{}`$ via
(59)
$$\begin{array}{cc}\hfill [h_l,\varphi _{j,t;m}^{}]& =h_l^\varphi ^{}(j,t)\varphi _{j,t;m}^{},l=1,2,3,4\hfill \\ \hfill h_1^\varphi ^{}(j,t)& =h_3^\varphi ^{}(j,t)=(t+1\delta _{j,2})\hfill \\ \hfill h_2^\varphi ^{}(j,t)& =t+1\delta _{j,0}\hfill \\ \hfill h_4^\varphi ^{}(j,t)& =(t1+\delta _{j,0}+2\delta _{j,1})\hfill \end{array}$$
The commutator with the grading operator $`d`$ is defined by
(60)
$$[d,\varphi _{j,t;m}^{}]=m\varphi _{j,t;m}^{}$$
Taking into account (59), (60) and (46), the set $`B^{}`$ of normal forms may be compared to the set $`_{1,1}^{h\backslash 0}`$ of semi standard super tableaux.
For fixed values of $`r_0`$ and $`p_1`$, the set of all normal forms
(61)
$$\varphi _{j_{p_1},t_{p_1};r_0+p_11}^{}\mathrm{}\varphi _{j_2,t_2;r_0+1}^{}\varphi _{j_1,t_1;r_0}^{}$$
is called $`B^{}(p_1;\mathrm{};r_0)`$.
For $`r_0`$ and $`(p_1,p_2,\mathrm{},p_R;\overline{p}_1,\overline{p}_2,\mathrm{},\overline{p}_{R1})`$ with $`p_i,\overline{p}_i`$ and $`R>1`$, the set of all normal forms $`\varphi _{j_N,t_N;m_N}^{}\mathrm{}\varphi _{j_2,t_2;m_2}^{}\varphi _{j_1,t_1;m_1}^{}`$ with
(62)
$$\begin{array}{cc}\hfill N& =\underset{i=1}{\overset{R}{}}p_i,m_1=r_0,\hfill \\ \hfill m_{i+1}m_i& =\{\begin{array}{cc}\overline{p}_{\widehat{R}}+1\hfill & \text{if }i=p_1+p_2+\mathrm{}+p_{\widehat{R}}\text{ for }1\widehat{R}<R,\hfill \\ 1\hfill & \text{otherwise.}\hfill \end{array}\hfill \end{array}$$
is denoted by $`B^{}(p_1,p_2,\mathrm{},p_R;\overline{p}_1,\overline{p}_2,\mathrm{},\overline{p}_{R1};r_0)`$.
The set of all semi-standard super tableaux associated with the horizontal border strip $`R>1,(p_1,\mathrm{},p_R;\overline{p}_1,\mathrm{},\overline{p}_{R1})`$ or $`R=1,(p_1,\mathrm{})`$ with the box counting fixed by $`r_0`$ can be mapped onto $`B^{}(p_1,\mathrm{},p_R;\overline{p}_1,\mathrm{},\overline{p}_{R1};r_0)`$ or $`B^{}(p_1;\mathrm{};r_0)`$, respectively. In particular, the reference labelling introduced in the previous section is mapped onto the normal form with $`j_r=2,t_r=0`$ $`r`$. This amounts to attributing the generator $`\varphi _{2,0;r}^{}`$ to the $`r+1`$-th box if the reference labelling assigns the number $`0`$ to this box. According to (46) and the correspondence between configurations and semi-standard super tableaux, the reference labelling is the only tableau of the border strip with the $`U_q\left(gl(2|2)\right)`$-weight given by $`(0,N,0,N)`$. Equation (59) specifies the only normal form in $`B^{}(p_1,\mathrm{},p_R;\overline{p}_1,\mathrm{},\overline{p}_{R1};r_0)`$ or $`B^{}(p_1;\mathrm{};r_0)`$ with the same values.
An arbitrary semi-standard super tableau in $`_{1,1}^{h\backslash 0}`$ with $`R1`$ is mapped onto a normal form (62) in three steps. It is convenient to rewrite the normal form (62) as
(63)
$$\begin{array}{cc}& \varphi _{j_{p_1},t_{p_1};s_1+p_11}^{}\mathrm{}\varphi _{j_2,t_2;s_1+1}^{}\varphi _{j_1,t_1;s_1}^{}\hfill \\ & \varphi _{j_{p_1+p_2},t_{p_1+p_2};s_2+p_21}^{}\mathrm{}\varphi _{j_{p_1+2},t_{p_1+2};s_2+1}^{}\varphi _{j_{p_1+1},t_{p_1+1};s_2}^{}\hfill \\ & \varphi _{j_{p_1+p_2+p_3},t_{p_1+p_2+p_3};s_3+p_31}^{}\mathrm{}\varphi _{j_{p_1+p_2+2},t_{p_1+p_2+2};s_3+1}^{}\varphi _{j_{p_1+p_2+1},t_{p_1+p_2+1};,s_3}^{}\hfill \\ & \mathrm{}\hfill \\ & \varphi _{j_N,t_N;s_R+p_R1}^{}\mathrm{}\varphi _{j_{p_1+\mathrm{}+p_{R1}+2},t_{p_1+\mathrm{}+p_{R1}+2};s_R+1}^{}\varphi _{j_{p_1+\mathrm{}+p_{R1}+1},t_{p_1+\mathrm{}+p_{R1}+1};s_R}^{}\hfill \end{array}$$
with the numbers $`s_i`$ defined in (49). The factors in the $`i`$-th line in (63) are related to the numbers given to the boxes in the $`i`$-th column and the right neighbouring row as follows.
1. For $`1iR`$ and $`p_i2`$, the factors
(64)
$$\begin{array}{c}\varphi _{j_{p_1+\mathrm{}+p_i},t_{p_1+\mathrm{}+p_i};s_i+p_i1}^{}\mathrm{}\varphi _{j_{p_1+\mathrm{}+p_{i1}+3},t_{p_1+\mathrm{}+p_{i1}+3};s_i+2}^{}\hfill \\ \hfill \varphi _{j_{p_1+\mathrm{}+p_{i1}+2},t_{p_1+\mathrm{}+p_{i1}+2};s_i+1}^{}\end{array}$$
are determined according to the numbers attributed to the $`p_i`$ lower boxes of the $`i`$-th column. Figure 3 specifies these factors for all cases allowed by the rules (47) and (48) for the tableaux of horizontal strips.
2. To obtain the factor $`\varphi _{j_{p_1+\mathrm{}+p_{i1}+1},t_{p_1+\mathrm{}+p_{i1}+1};s_i}^{}`$ with $`i>1`$ and $`p_i1`$, the $`i`$-th column and the $`(i1)`$-th finite row with length $`\overline{p}_{i1}+1`$ is taken into account. Figure 4 collects all possible cases together with the associated factors.
3. Depending on the numbers given to the boxes of the upper infinite column, the rightmost factor $`\varphi _{j_1,t_1;s_1}^{}`$ is determined. All cases are listed in figure 5.
For a semi-standard super tableau of the border strip $`(p_1,\mathrm{})`$ with $`p_1>1`$, the normal form (61) is determined by steps (2) and (3). In the case $`(1;\mathrm{})`$, only step (3) is required. A box on the lhs of figures 3-5 is drawn boldly if it is the $`(r+1)`$th box in the complete border strip and the corresponding expression on the rhs contains a generator $`\varphi _{j,t;r}^{}`$.
Each correspondence listed in figures 3-5 matches the $`U_q\left(gl(2|2)\right)`$-weights introduced in (46) and (59). Moreover, the sum $`m_1+m_2+\mathrm{}+m_N`$ coincides with the number $`d^h(p_1,\mathrm{},p_R;\overline{p}_1,\mathrm{},\overline{p}_{R1};r_0)`$ introduced in (50). As a consequence of the rules (47)-(48), an arbitrary tableau in $`_{1,1}^{h\backslash 0}`$ assigns the number $`1`$ to all boxes of the lower half-infinite row except the rightmost of them. Thus the prescriptions given in figures 3-5 specify one normal form of $`B^{}(p_1,\mathrm{},p_R;\overline{p}_1,\mathrm{},\overline{p}_{R1};r_0)`$ $`\left(\text{or}B^{}(p_1;\mathrm{};r_0)\right)`$ for each tableau of the border strip with the parameters $`(p_1,\mathrm{},p_R;\overline{p}_1,\mathrm{},\overline{p}_{R1})`$ $`\left(\text{or}(p_1;\mathrm{})\right)`$ and $`r_0`$. Due to the definition (57), the same prescriptions attribute exactly one semi-standard super tableau related for the parameters $`(p_1,\mathrm{},p_R;\overline{p}_1,\mathrm{},\overline{p}_{R1})`$ $`\left(\text{or}(p_1;\mathrm{})\right)`$ and $`r_0`$ to each normal form in $`B^{}(p_1,\mathrm{},p_R;\overline{p}_1,\mathrm{},\overline{p}_{R1};r_0)`$ $`\left(\text{or}B^{}(p_1;\mathrm{};r_0)\right)`$.
###### Result 1.
The semi-standard super tableaux in $`_{1,1}^{h\backslash 0}`$ related a horizontal border strip with the parameters $`r_0`$, $`(p_1,\mathrm{},p_R;\overline{p}_1,\mathrm{},\overline{p}_{R1})\left(\text{or}(p_1;\mathrm{})\right)`$, and the normal forms in $`B^{}(p_1,\mathrm{},p_R;\overline{p}_1,\mathrm{},\overline{p}_{R1};r_0)\left(\text{or}B^{}(p_1;\mathrm{};r_0)\right)`$ are in one-to-one correspondence. Hence, the tableaux in $`_{1,1}^{h\backslash 0}`$ and the normal forms in $`B^{}`$ are in one-to-one correspondence.
The set $`_{1,1}^{v\backslash 0}`$ of semi-standard super tableaux associated with the vertical border strips is related to the set of normal forms of the creation algebra $`𝒜`$.
###### Definition 2.
The creation algebra $`𝒜`$ is generated by $`\{\varphi _{j,t;m}|0j3,t_0,m\}`$ over $`𝒵`$ subject to the defining relations
(65)
$$\varphi _{j_2,t_2;m_2}\varphi _{j_1,t_1;m_1}=\varphi _{j_2,t_2;m_1+t_1+\theta _{j_1,j_2}}\varphi _{j_1,t_1;m_2t_1\theta _{j_1,j_2}}$$
with $`\theta _{j,j^{}}`$ defined by (54).
The procedure applying to the case of vertical border strips is quite analogous to the one outlined above for the horizontal border strips. A product
(66)
$$\varphi _{j_n,t_n;m_n}\mathrm{}\varphi _{j_2,t_2;m_2}\varphi _{j_1,t_1;m_1}$$
is referred to as a normal form (of $`𝒜`$) iff
(67)
$$m_{i+1}>m_i+t_i+\theta _{j_i,j_{i+1}}\text{for}\mathrm{\hspace{0.33em}\hspace{0.33em}1}i<n$$
According to corollary 2 in , the set
(68)
$$B=_n\left\{\varphi _{j_n,t_n;m_n}\mathrm{}\varphi _{j_1,t_1;m_1}\right|j_i=0,1,2,3,t_i_0,m_i\text{satisfy (}\text{67}\text{)}\}$$
provides a $``$-linear basis of $`𝒜`$. A $`U_q\left(gl(2|2)\right)`$-weight is defined for the generators $`\varphi _{j,t;m}`$ by
(69)
$$[h_l,\varphi _{j,t;m}]=h_l^\varphi ^{}(j,t)\varphi _{j,t;m},l=1,2,3,4$$
with $`h_l^\varphi ^{}`$ given by (59). A mapping of the set $`_{1,1}^{v\backslash 0}`$ of semi-standard super tableaux onto the set $`B`$ similar to the mapping of $`_{1,1}^{h\backslash 0}`$ on $`B^{}`$ is specified in Appendix A. The set of all normal forms $`\varphi _{j_{\overline{p}_1},t_{\overline{p}_1};r_0+\overline{p}_11}\mathrm{}\varphi _{j_1,t_2;r_0+1}\varphi _{j_1,t_1;r_0}`$ with fixed $`r_0`$ and $`\overline{p}_1`$ is denoted by $`B(\mathrm{};\overline{p}_1)`$. For $`r_0`$ and $`(p_1,p_2,\mathrm{},p_R;\overline{p}_1,\overline{p}_2,\mathrm{},\overline{p}_{R1})`$ with $`p_i,\overline{p}_i`$ and $`R>1`$, the set of all normal forms $`\varphi _{j_{\overline{N}},t_{\overline{N}};m_{\overline{N}}}\mathrm{}\varphi _{j_2,t_2;m_2}\varphi _{j_1,t_1;m_1}`$ with
(70)
$$\begin{array}{cc}\hfill \overline{N}& =\underset{i=1}{\overset{R}{}}\overline{p}_i,m_1=r_0,\hfill \\ \hfill m_{i+1}m_i& =\{\begin{array}{cc}p_{\widehat{R}}+1\hfill & \text{if }i=\overline{p}_1+\overline{p}_2+\mathrm{}+\overline{p}_{\widehat{R}}\text{ for }1\widehat{R}<R,\hfill \\ 1\hfill & \text{otherwise.}\hfill \end{array}\hfill \end{array}$$
is denoted by $`B(p_1,p_2,\mathrm{},p_{R1};\overline{p}_1,\overline{p}_2,\mathrm{},\overline{p}_R;r_0)`$.
###### Result 2.
The semi-standard super tableaux in $`_{1,1}^{v\backslash 0}`$ related to a vertical border strip with parameters $`r_0`$, $`(p_1,\mathrm{},p_{R1};\overline{p}_1,\mathrm{},\overline{p}_R)`$ $`\left(\text{or}(\mathrm{},\overline{p}_1)\right)`$ and the normal forms in $`B(p_1,\mathrm{},p_{R1};\overline{p}_1,\mathrm{},\overline{p}_R;r_0)`$ $`\left(\text{or}B(\mathrm{},\overline{p}_1)\right)`$ are in one-to-one correspondence .
For a normal form $`\varphi _{j_{\overline{N}},t_{\overline{N}};m_{\overline{N}}}\mathrm{}\varphi _{j_1,t_1;m_1}`$ in $`B(p_1,\mathrm{},p_{R1};\overline{p}_1,\mathrm{},\overline{p}_R;r_0)`$ or $`B(\mathrm{},\overline{p}_1)`$, the sum $`m_1+\mathrm{}+m_{\overline{N}}`$ equals the number $`d^v(p_1,\mathrm{},p_{R1};\overline{p}_1,\mathrm{},\overline{p}_R;r_0)`$ introduced in (51).
In the following section, the definition of the algebra $`𝒜^{}`$ (or $`𝒜`$) is motivated by a naive limit of the commutation relations of type II vertex operators. Besides the components related to the generators $`\varphi _{j,t;m}^{}`$ (or $`\varphi _{j,t;m}`$), the vertex operators have a further component equal to the unit. Each of the sets $`B^{}`$, $`B`$, may be supplemented by the unit. The resulting sets are denoted by $`B_1^{}`$ and $`B_1`$, respectively. Formally, the unit may be assigned to the tableau in $`_{1,1}^h`$ or $`_{1,1}^v`$ attributing the number one to each box of the border strip given by the single infinite or half-infinite row or column. Application of the map between the components and the semi-standard supertableaux specified in section 4 yields a relation between the enlarged sets and the infinite configurations.
###### Result 3.
The infinite components $`(\mathrm{}w_{k_4}w_{k_2}w_{k_0}w_{k_2}w_{k_4}\mathrm{})`$ with $`k_r=1`$ for almost all $`r`$ are in one-to-one correspondence with the set $`B_1^{}`$. The infinite components $`(\mathrm{}w_{k_3}^{}w_{k_1}^{}w_{k_1}^{}w_{k_3}^{}w_{k_5}^{}\mathrm{})`$ with $`k_r=1`$ for almost all $`r`$ are in one-to-one correspondence with the set $`B_1`$.
## 6. Type II vertex operators
The type II vertex operators considered below are intertwiners of $`U`$-modules of the form
(71)
$$\mathrm{\Phi }_{V_I}^{VV_J}(z):V_IV_zV_J$$
where $`V_I`$ and $`V_J`$ are level-one $`U_q\left(\widehat{sl}(2|2)\right)/`$-modules. $`V`$ with basis $`\{v_j\}`$ denotes a $`U_q^{}\left(\widehat{sl}(2|2)\right)`$-module obtained from an infinite-dimensional $`U_q\left(gl(2|2)\right)`$-module by means of the evaluation homomorphism , , . The evaluation module $`V_z=V[z,z^1]`$ is endowed with a $`U_q\left(\widehat{sl}(2|2)\right)`$-structure via
(72) $`E_n^{k,\pm }\left(v_jz^m\right)`$ $`=E_n^{k,\pm }v_jz^{m+n}`$ $`n`$
$`\mathrm{\Psi }_{\pm n}^{l,\pm }\left(v_jz^m\right)`$ $`=\mathrm{\Psi }_{\pm n}^{l,\pm }v_jz^{m\pm n}`$ $`n0`$
$`d\left(v_jz^m\right)`$ $`=mv_jz^m`$
$`c\left(v_jz^m\right)`$ $`=0`$
The vertex operator (71) is introduced as the formal series
(73)
$$\begin{array}{cc}\hfill \mathrm{\Phi }_{V_I}^{VV_J}(z)& =\underset{j}{}v_j\left(\mathrm{\Phi }_{V_I}^{VV_J}\right)_j(z)\hfill \\ \hfill \left(\mathrm{\Phi }_{V_I}^{VV_J}\right)_j(z)& =\underset{m}{}\left(\mathrm{\Phi }_{V_I}^{VV_J}\right)_{j;m}z^m\hfill \end{array}$$
In terms of the maps $`\left(\mathrm{\Phi }_{V_I}^{VV_J}\right)_{j;m}:V_IV_J`$ the intertwining property reads
(74)
$$\begin{array}{c}\mathrm{\Delta }(a)\left\{\underset{j}{}\underset{m}{}\left(v_jz^m\right)\left(\mathrm{\Phi }_{V_I}^{VV_J}\right)_{j;m}\right\}\hfill \\ \hfill =\underset{j}{}\underset{m}{}(1)^{|a|(|v_j|+|\mathrm{\Phi }_j|)}\left(v_jz^m\right)\left(\mathrm{\Phi }_{V_I}^{VV_J}\right)_{j;m}a\end{array}$$
$`aU_q\left(\widehat{gl}(2|2)\right)`$. Here $`|v_j|`$ and $`|\mathrm{\Phi }_j|`$ denote the $`_2`$-gradings of $`v_j`$ and the component $`\left(\mathrm{\Phi }_{V_I}^{VV_J}\right)_{j;m}`$. In the remainder, their relation is fixed by $`|\mathrm{\Phi }_j|=|v_j|`$. For the generators of $`U_q\left(gl(2|2)\right)`$, $`q^{\pm \frac{1}{2}c}`$ and the grading operator $`d`$, the coproduct $`\mathrm{\Delta }`$ is defined by
(75) $`\mathrm{\Delta }\left(E_0^{k,+}\right)`$ $`=E_0^{k,+}1+q^{h_k}E_0^{k,+}`$ $`\mathrm{\Delta }\left(E_0^{k,}\right)`$ $`=E_0^{k,}q^{h_k}+1E_0^{k,}`$
$`\mathrm{\Delta }\left(q^{\pm h_l}\right)`$ $`=q^{\pm h_l}q^{\pm h_l}`$ $`k=1,2,3,`$ $`l=1,2,3,4`$
and
(76)
$$\mathrm{\Delta }\left(q^{\pm \frac{1}{2}c}\right)=q^{\pm \frac{1}{2}c}q^{\pm \frac{1}{2}c}\mathrm{\Delta }(d)=d1+1d$$
A partial information on the coproduct of the remaining Drinfeld generators of $`U_q\left(\widehat{gl}(2|2)\right)`$ is provided by the formulae
(77) $`\mathrm{\Delta }\left(E_n^{k,+}\right)`$ $`=E_n^{k,+}q^{nc}+q^{2nc+h_k}E_n^{k,+}`$
$`+{\displaystyle \underset{n^{}=0}{\overset{n1}{}}}q^{\frac{1}{2}(n+3n^{})c}\mathrm{\Psi }_{nn^{}}^{k,+}q^{(nn^{})c}E_n^{}^{k,+}modN_{}N_+^2`$
$`\mathrm{\Delta }\left(E_n^{k,+}\right)`$ $`=E_n^{k,+}q^{nc}+q^{h_k}E_n^{k,+}`$
$`+{\displaystyle \underset{n^{}=1}{\overset{n1}{}}}q^{\frac{1}{2}(nn^{})c}\mathrm{\Psi }_{n+n^{}}^{k,}q^{(nn^{})c}E_n^{}^{k,+}modN_{}N_+^2`$
$`\mathrm{\Delta }\left(E_n^{k,}\right)`$ $`=E_n^{k,}q^{h_k}+q^{nc}E_n^{k,}`$
$`+{\displaystyle \underset{n^{}=1}{\overset{n1}{}}}q^{(nn^{})c}E_n^{}^{k,}q^{\frac{1}{2}(nn^{})c}\mathrm{\Psi }_{nn^{}}^{k,+}modN_{}^2N_+`$
$`\mathrm{\Delta }\left(E_n^{k,}\right)`$ $`=E_n^{k,}q^{2mch_k}+q^{nc}E_m^{k,}`$
$`+{\displaystyle \underset{n^{}=0}{\overset{n1}{}}}q^{(nn^{})c}E_n^{}^{k,}q^{\frac{1}{2}(n+3n^{})}\mathrm{\Psi }_{n+n^{}}^{k,}modN_{}^2N_+`$
$`\mathrm{\Delta }\left(H_n^l\right)`$ $`=H_n^lq^{\frac{1}{2}nc}+q^{\frac{3}{2}nc}H_n^lmodN_{}N_+`$
$`\mathrm{\Delta }\left(H_n^l\right)`$ $`=H_n^lq^{\frac{3}{2}nc}+q^{\frac{1}{2}nc}H_n^lmodN_{}N_+`$
for $`k=1,2,3`$, $`l=1,2,3,4`$, and $`n>0`$. $`N_\pm `$ and $`N_\pm ^2`$ are left $`(q)[q^{\pm c},q^{\pm h_l},\mathrm{\Psi }_{\pm n}^{l,\pm }]`$-modules generated by $`\left\{E_m^{k,\pm }\right\}_{k=1,2,3,m}`$ and $`\left\{E_m^{k,\pm }E_m^{}^{k^{},\pm }\right\}_{k,k^{}=1,2,3,m,m^{}}`$, respectively. Equations (77) are readily verified using the Hopf algebra structure of $`U_q\left(\widehat{gl}(2|2)\right)`$ given in . In context with vertex models based on the quantum affine algebra $`U_q\left(\widehat{gl}(N)\right)`$, type II vertex operators associated with finite-dimensional evaluation modules of the algebra are considered. Expressions analogous to (77) allow to derive free boson expressions for a particular component of a vertex operator (see , , and , for example). Similar results have been obtained in for $`U_q\left(\widehat{gl}(2|2)\right)`$-vertex operators related to the four-dimensional modules $`W_z`$ and $`W_z^{}`$ defined in section 2.2.
For two $`U_q\left(gl(2|2)\right)`$-modules $`V^{(1)}`$ and $`V^{(2)}`$ with bases $`\{v^{(1)}\}_j`$ and $`\{v^{(2)}\}_j^{}`$, the $`R`$-matrix $`R_{V^{(1)}V^{(2)}}\left(\frac{z}{w}\right)End(V_z^{(1)}V_w^{(2)})`$ intertwines the action of $`U_q\left(\widehat{sl}(2|2)\right)`$ on the tensor product of the two evaluation modules $`V_z^{(1)}`$ and $`V_w^{(2)}`$:
(78)
$$R_{V^{(1)}V^{(2)}}\left(\frac{z}{w}\right)\mathrm{\Delta }(a)=\mathrm{\Delta }^{}(a)R_{V^{(1)}V^{(2)}}\left(\frac{z}{w}\right)aU_q\left(\widehat{sl}(2|2)\right)$$
where $`\mathrm{\Delta }^{}(a)=P^{gr}\mathrm{\Delta }(a)`$ with $`P^{gr}(x_1x_2)=(1)^{|x_1||x_2|}x_2x_1`$. It is convenient to introduce a second matrix $`\overline{R}_{V^{(1)}V^{(2)}}\left(\frac{z}{w}\right)`$ by
(79)
$$\begin{array}{cc}\hfill R_{V^{(1)}V^{(2)}}\left(\frac{z}{w}\right)\left(v_{j_1}^{(1)}v_{j_2}^{(2)}\right)& =\underset{j_3,j_4}{}(1)^{|v_{j_3}^{(1)}||v_{j_4}^{(2)}|}\left(\overline{R}_{V^{(1)}V^{(2)}}\right)_{j_1,j_2}^{j_3,j_4}\left(\frac{z}{w}\right)v_{j_3}^{(1)}v_{j_4}^{(2)}\hfill \\ \hfill \overline{R}_{V^{(1)}V^{(2)}}\left(\frac{z}{w}\right)\left(v_{j_1}^{(1)}v_{j_2}^{(2)}\right)& =\underset{j_3,j_4}{}\left(\overline{R}_{V^{(1)}V^{(2)}}\right)_{j_1,j_2}^{j_3,j_4}\left(\frac{z}{w}\right)v_{j_3}v_{j_4}\hfill \end{array}$$
Uniqueness of the normalised vertex operators (71) provided, their commutation relation reads
(80)
$$\begin{array}{c}\left(\mathrm{\Phi }_{V_J}^{V^{(1)}V_K}\right)_{j_1}(z)\left(\mathrm{\Phi }_{V_I}^{V^{(2)}V_J}\right)_{j_2}(w)=\hfill \\ \hfill c\left(\frac{z}{w}\right)\underset{j_3,j_4}{}\left(\overline{R}_{V^{(1)}V^{(2)}}\right)_{j_1,j_2}^{j_3,j_4}\left(\frac{z}{w}\right)\left(\mathrm{\Phi }_{V_J}^{V^{(2)}V_K}\right)_{j_4}(w)\left(\mathrm{\Phi }_{V_I}^{V^{(1)}V_J}\right)_{j_3}(z)\end{array}$$
with a scalar function $`c\left(\frac{z}{w}\right)`$. The relation (80) follows from the argument given in chapter 6 of for a pair of evaluation modules $`V_z^{(1)}`$ and $`V_w^{(2)}`$ with an intertwiner $`R_{V^{(1)}V^{(2)}}\left(\frac{z}{w}\right)`$. For quantum affine algebras, the normalisation of the $`R`$-matrix governing the analogous commutations relations has been determined via free boson realizations or from the related two-point functions obtained as solutions of certain $`q`$-difference equations . A free boson realization has also been used in to fix the commutation relations of two $`U_q\left(\widehat{gl}(2|2)\right)`$-vertex operators related to $`W_z`$ or $`W_z^{}`$.
The $`U_q^{}\left(\widehat{sl}(2|2)\right)`$-module relevant to the configuration spaces $`\mathrm{\Omega }_A^{(1)}`$ and $`\mathrm{\Omega }_B^{(1)}`$ of the present vertex model is constructed from the infinite-dimensional, irreducible $`U_q\left(gl(2|2)\right)`$-module $`V_{\overline{\mathrm{\Lambda }}_2+\overline{\mathrm{\Lambda }}_4}`$ with basis $`\{v_{j,t}\}_{0j3,t_0}`$. Its $`U_q\left(gl(2|2)\right)`$-weights are
(81) $`h_1v_{j,t}`$ $`=(t+1\delta _{j,2})v_{j,t}`$ $`h_2v_{j,t}`$ $`=(t+1\delta _{j,0})v_{j,t}`$
$`h_3v_{j,t}`$ $`=(t+1\delta _{j,2})v_{j,t}`$ $`h_4v_{j,t}`$ $`=(t1+\delta _{j,0}+2\delta _{j,1})v_{j,t}`$
A $`_2`$-grading is defined by $`|v_{0,t}|=|v_{2,t}|=1`$ and $`|v_{1,t}|=|v_{3,t}|=0`$ $`t`$. The evaluation homomorphism found in yields a $`U_q^{}\left(\widehat{sl}(2|2)\right)`$-structure given by the action of $`h_1,h_2,h_3`$ specified in (81) and
(82)
$$\begin{array}{cc}\hfill H_m^1v_{j,t}& =\frac{1}{m}q^{m(t+1\delta _{j,2})}\left[m(t+1\delta _{j,2})\right]v_{j,t}\hfill \\ \hfill H_m^2v_{j,t}& =\frac{1}{m}q^{m(t+\delta _{j,0}+2\delta _{j,1})}\left[m(t+1\delta _{j,0})\right]v_{j,t}\hfill \\ \hfill H_m^3v_{j,t}& =H_m^1v_{j,t}\hfill \end{array}$$
for $`m0`$ and
(83) $`E_m^{1,}v_{2,t}`$ $`=v_{1,t1}`$ $`E_m^{1,+}v_{1,t}`$ $`=[t+1]v_{2,t+1}`$
$`E_m^{1,}v_{3,t}`$ $`=v_{0,t}`$ $`E_m^{1,+}v_{0,t}`$ $`=[t+1]v_{3,t}`$
$`E_m^{2,}v_{1,t}`$ $`=q^{m(2t+3)}\frac{[t+1]}{[t+2]}v_{0,t+1}`$ $`E_m^{2,+}v_{0,t+1}`$ $`=q^{m(2t+3)}[t+2]v_{1,t}`$
$`E_m^{2,}v_{2,t}`$ $`=q^{m(2t+1)}v_{3,t}`$ $`E_m^{2,+}v_{3,t}`$ $`=q^{m(2t+1)}[t+1]v_{2,t}`$
$`E_m^{3,}v_{0,t}`$ $`=q^{2m(t+1)}v_{1,t}`$ $`E_m^{3,+}v_{1,t}`$ $`=q^{2m(t+1)}[t+1]v_{0,t}`$
$`E_m^{3,}v_{3,t}`$ $`=q^{2m(t+1)}v_{2,t+1}`$ $`E_m^{3,+}v_{2,t+1}`$ $`=q^{2m(t+1)}[t+1]v_{3,t}`$
$`m`$. The basis $`\{v_{j,t}\}_{0j3}`$ coincides with the basis of the module $`V`$ introduced in . A reducible, infinite-dimensional $`U_q^{}\left(\widehat{sl}(2|2)\right)`$-module $`\stackrel{̊}{V}`$ with basis $`v_{1,1}\{v_{j,t}\}_{0j3,t_0}`$ is defined by (81)-(83) and
(84)
$$\begin{array}{cc}\hfill h_lv_{1,1}& =0l=1,2,3,4\hfill \\ \hfill H_m^kv_{1,1}& =0k=1,2,3,m0\hfill \end{array}$$
and
(85)
$$\begin{array}{cc}\hfill E_m^{1,}v_{1,1}& =E_m^{2,+}v_{1,1}=E_m^{3,\pm }v_{1,1}=0\hfill \\ \hfill E_m^{1,+}v_{1,1}& =\frac{[m]}{m}v_{2,0}\hfill \\ \hfill E_m^{2,}v_{1,1}& =q^m\frac{[m]}{m}v_{0,0}\hfill \end{array}$$
$`m`$. Any action of Drinfeld generators of $`U_q^{}\left(\widehat{sl}(2|2)\right)`$ on $`\stackrel{̊}{V}`$ not listed in (81)-(85) vanishes. Specifying the action of all $`H_m^4`$ on $`v_{j,t}`$ for one pair $`j,t`$ determines a $`U_q^{}\left(\widehat{gl}(2|2)\right)`$-structure on $`\stackrel{̊}{V}`$. The appropriate choice will be given below.
For shorter notation, the components $`\left(\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}V_J}\right)_{j,t}(z)`$ may be written $`\mathrm{\Phi }_{j,t}(z)`$. The intertwining property (74) with $`a=e_k`$, $`k=1,2,3`$, relates all components $`\left(\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}V_J}\right)_{j,t}(z)`$ with $`j=0,1,2,3`$ and $`t0`$:
(86)
$$\begin{array}{cc}\hfill [t+1]\mathrm{\Phi }_{3,t}(z)& =q^{t+1}e_2\mathrm{\Phi }_{2,t}(z)+\mathrm{\Phi }_{2,t}(z)e_2\hfill \\ \hfill [t+1]\mathrm{\Phi }_{2,t+1}(z)& =q^{t+1}e_3\mathrm{\Phi }_{3,t}(z)+\mathrm{\Phi }_{3,t}(z)e_3\hfill \\ \hfill [t+1]\mathrm{\Phi }_{0,t}(z)& =q^{t1}e_1\mathrm{\Phi }_{3,t}(z)\mathrm{\Phi }_{3,t}(z)e_1\hfill \\ \hfill [t+1]\mathrm{\Phi }_{1,t}(z)& =q^{t1}e_1\mathrm{\Phi }_{2,t+1}(z)\mathrm{\Phi }_{2,t+1}(z)e_1\hfill \end{array}$$
According to (75)-(77), the intertwining property (74) imposes the following conditions on the component $`\mathrm{\Phi }_{2,0}(z)`$:
(87)
$$\begin{array}{cc}\hfill E_m^{k,}\mathrm{\Phi }_{2,0}(z)(1)^{|\mathrm{\Phi }_{2,0}|}\mathrm{\Phi }_{2,0}(z)E_m^{k,}& =0k=1,2,3\hfill \\ \hfill E_m^{3,+}\mathrm{\Phi }_{2,0}(z)(1)^{|\mathrm{\Phi }_{2,0}|}\mathrm{\Phi }_{2,0}(z)E_m^{3,+}& =0\hfill \end{array}$$
(88)
$$E_m^{1,+}\mathrm{\Phi }_{2,0}(z)(1)^{|\mathrm{\Phi }_{2,0}|}\mathrm{\Phi }_{2,0}(z)E_m^{1,+}=(1)^{|v_{2,0}|}q^mz^m\mathrm{\Phi }_{1,1}(z)$$
$`m`$ and
(89)
$$\begin{array}{cc}\hfill [h_l,\mathrm{\Phi }_{2,0}(z)]& =(\delta _{l,2}+\delta _{l,4})\mathrm{\Phi }_{2,0}(z)l=1,2,3,4\hfill \\ \hfill w^d\mathrm{\Phi }_{2,0}(z)w^d& =\mathrm{\Phi }_{2,0}(wz)\hfill \\ \hfill [H_{\pm n}^1,\mathrm{\Phi }_{2,0}(z)]& =[H_{\pm n}^3,\mathrm{\Phi }_{2,0}(z)]=0\hfill \\ \hfill [H_n^2,\mathrm{\Phi }_{2,0}(z)]& =q^{\frac{1}{2}n}\frac{[n]}{n}z^n\mathrm{\Phi }_{2,0}(z)\hfill \\ \hfill [H_n^2,\mathrm{\Phi }_{2,0}(z)]& =q^{\frac{3}{2}n}\frac{[n]}{n}z^n\mathrm{\Phi }_{2,0}(z)\hfill \end{array}$$
for $`n>0`$. The conditions (87) and (89) are fulfilled by
(90)
$$\begin{array}{cc}& \mathrm{\Phi }_{2,0}(z)=\hfill \\ & :\mathrm{exp}\{\beta ^1(qz)i\phi ^1\phi ^2\mathrm{ln}(qz)(\phi _0^1i\phi _0^2)i\pi \phi _0^3+C(\phi _0^1i\phi _0^2\phi _0^3+i\phi _0^4)\hfill \\ & +\underset{m0}{}C_m(\phi _m^1i\phi _m^2\phi _m^3+i\phi _m^4)z^m+\underset{m0}{}\frac{1}{[m]}q^{\frac{1}{2}|m|m}(\phi _m^3i\phi _m^4)z^m\}:\hfill \end{array}$$
and
(91)
$$\begin{array}{c}\mathrm{\Phi }_{1,1}(z)=:\mathrm{exp}\{C(\phi _0^1i\phi _0^2\phi _0^3+i\phi _0^4)\hfill \\ \hfill +\underset{m0}{}(C_m\frac{1}{[m]}q^{\frac{1}{2}|m|m})(\phi _m^1i\phi _m^2\phi _m^3+i\phi _m^4)z^m\}:\end{array}$$
for $`C,C_m`$. Here the free boson realizations (5), (8) and (11) for the Drinfeld generators have been employed. The expressions (90) and (91) satisfy
(92)
$$\mathrm{\Phi }_{2,0}(z)\mathrm{\Phi }_{2,0}(w)=\mathrm{\Phi }_{2,0}(w)\mathrm{\Phi }_{2,0}(z)$$
and
(93)
$$\mathrm{\Phi }_{1,1}(z)\mathrm{\Phi }_{1,1}(w)=\mathrm{\Phi }_{1,1}(w)\mathrm{\Phi }_{1,1}(z)\mathrm{\Phi }_{1,1}(z)\mathrm{\Phi }_{2,0}(w)=\mathrm{\Phi }_{2,0}(w)\mathrm{\Phi }_{1,1}(z)$$
As required by (74), $`\mathrm{\Phi }_{1,1}(z)`$ commutes with $`h_l`$, $`H_{\pm n}^k`$ and $`E_m^{k,\pm }`$ with $`k=1,2,3,l=1,2,3,4,n>0,m0`$. Due to (77), the property (74) is satisfied for $`a=H_m^4`$ provided that
(94)
$$H_m^4v_{2,0}=\frac{[m]}{m}\left(2q^{m+\frac{1}{2}|m|}[m]C_m+1\right)z^mv_{2,0}m0$$
The components $`\mathrm{\Phi }_{1,1;m}`$ and $`\mathrm{\Phi }_{2,0;m}`$ defined by $`\mathrm{\Phi }_{1,1}(z)=_m\mathrm{\Phi }_{1,1;m}z^m`$ and $`\mathrm{\Phi }_{2,0}(z)=_m\mathrm{\Phi }_{2,0;m}z^m`$ may be applied to the level-one modules $`\stackrel{̊}{V}(\mathrm{\Lambda }_0)`$, $`\stackrel{̊}{V}(\mathrm{\Lambda }_1+\mathrm{\Lambda }_4)`$, $`\stackrel{̊}{V}(2\mathrm{\Lambda }_2\mathrm{\Lambda }_3+\mathrm{\Lambda }_4)`$ or $`V\left(s\mathrm{\Lambda }_1+(1+s)\mathrm{\Lambda }_2(1+2s^{})\mathrm{\Lambda }_4\right)`$ introduced in section 3. Use of the realizations (36) and taking into account the property (31) reveals that the components $`\mathrm{\Phi }_{1,1;m}`$ and $`\mathrm{\Phi }_{2,0;m}`$ map each of these level-one modules onto itself if
(95)
$$[m]C_m=q^{\frac{1}{2}|m|m}\text{for}m<0$$
This suggests the existence of intertwiners of the form (71) with $`V=\stackrel{̊}{V}`$ and $`V_I=V_J`$ given by any of the four level-one modules. The validity of the entire set of equations (74) remains to be demonstrated.
###### Conjecture 2.
There exist type II vertex operators
(96)
$$\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}V_I}(z):V_I\stackrel{̊}{V}_zV_II=0,1,3,4$$
for $`V_0=\stackrel{̊}{V}(\mathrm{\Lambda }_0)`$, $`V_1=\stackrel{̊}{V}(\mathrm{\Lambda }_1+\mathrm{\Lambda }_4)`$, $`V_3=\stackrel{̊}{V}(2\mathrm{\Lambda }_0\mathrm{\Lambda }_3+\mathrm{\Lambda }_4)`$ and $`V_4=V\left(s\mathrm{\Lambda }_1+(1+s)\mathrm{\Lambda }_2(1+2s^{})\mathrm{\Lambda }_4\right)`$. They can be normalised such that
(97)
$$\begin{array}{c}\left(\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}V_I}\right)_{1,1}(z)=:\mathrm{exp}\{C(\phi _0^1i\phi _0^2\phi _0^3+i\phi _0^4)\hfill \\ \hfill +\underset{m>0}{}(C_m\frac{1}{[m]}q^{\frac{1}{2}|m|m})(\phi _m^1i\phi _m^2\phi _m^3+i\phi _m^4)z^m\}:\end{array}$$
with arbitrary complex numbers $`C`$ and $`C_m`$ with $`m>0`$. A further component is given by
(98)
$$\begin{array}{c}\left(\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}V_I}\right)_{2,0}(z)=\hfill \\ \hfill :\mathrm{exp}\{\beta ^1(qz)i\phi ^{1,+}(qz)\phi ^{2,+}(qz)i\pi \phi _0^3+C(\phi _0^1i\phi _0^2\phi _0^3+i\phi _0^4)\\ \hfill +\underset{m>0}{}(C_m\frac{1}{[m]}q^{\frac{1}{2}|m|m})(\phi _m^1i\phi _m^2\phi _m^3+i\phi _m^4)z^m\}:\end{array}$$
All the remaining components $`\left(\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}V_I}\right)_{j,t}(z)`$, $`0j3`$, $`t\delta _{j,2}0`$, follow from (86) and (98). The vertex operators satisfy the commutation relation
(99)
$$\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}V_I}(z)\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}V_I}(w)=\overline{R}_{\stackrel{̊}{V}\stackrel{̊}{V}}\left(\frac{z}{w}\right)\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}V_I}(w)\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}V_I}(z)$$
with the $`R`$-matrix $`R_{\stackrel{̊}{V}\stackrel{̊}{V}}\left(\frac{z}{w}\right):\stackrel{̊}{V}_z\stackrel{̊}{V}_w\stackrel{̊}{V}_z\stackrel{̊}{V}_w`$ normalised by
(100)
$$\left(\overline{R}_{\stackrel{̊}{V}\stackrel{̊}{V}}\right)_{2,0;\mathrm{\hspace{0.17em}2},0}^{2,0;\mathrm{\hspace{0.17em}2},0}\left(\frac{z}{w}\right)=1$$
The normalisation (100) stems from the relation (92). All $`R`$-matrix elements involving the index pair $`1,1`$ are listed by $`\left(\overline{R}_{\stackrel{̊}{V}\stackrel{̊}{V}}\right)_{1,1;\mathrm{\hspace{0.17em}1},1}^{1,1;\mathrm{\hspace{0.17em}1},1}(z)=1`$ and $`\left(\overline{R}_{\stackrel{̊}{V}\stackrel{̊}{V}}\right)_{1,1;j,t}^{1,1;j,t}(z)=\left(\overline{R}_{\stackrel{̊}{V}\stackrel{̊}{V}}\right)_{j,t;\mathrm{\hspace{0.17em}1},1}^{j,t;\mathrm{\hspace{0.17em}1},1}(z)=1j,t`$.
The action of $`H_m^4`$ on $`v_{j,t}`$ follows from the values of $`C_m`$. With
(101)
$$C=0,[m]C_m=q^{\frac{|m|}{2}m}m0$$
equation (94) yields
(102) $`H_m^4v_{1,1}`$ $`=0`$
$`H_m^4v_{0,t}`$ $`=\frac{1}{m}q^{m(t+1)}[mt]v_{0,t}`$ $`H_m^4v_{1,t}`$ $`=\frac{1}{m}q^{m(t+2)}[m(t+1)]v_{1,t}`$
$`H_m^4v_{2,t}`$ $`=\frac{1}{m}\left([m]q^{m(t+1)}[tm]\right)v_{2,t}`$ $`H_m^4v_{3,t}`$ $`=\frac{1}{m}\left([m]q^{m(t+1)}[tm]\right)v_{3,t}`$
$`m0`$ and $`t_0`$. (101) is a convenient choice kept for the remainder of this section.
A similar analysis indicates the existence of type II vertex operators
(103)
$$\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}^{}V_I}(z):V_I\stackrel{̊}{V}_z^{}V_II=0,1,3,4$$
for the dual module $`\stackrel{̊}{V}^{}`$. In terms of the dual basis $`v_{1,1}^{}\{v_{j,t}^{}\}_{0j3,t_0}`$ with $`v_{j_1,t_1}^{}|v_{j_2,t_2}=\delta _{j_1,j_2}\delta _{t_1,t_2}`$, the $`U_q^{}\left(\widehat{gl}(2|2)\right)`$-structure on $`\stackrel{̊}{V}^{}`$ is introduced by
(104)
$$av^{}|v=(1)^{|a||v^{}|}v^{}|S(a)vaU_q^{}\left(\widehat{gl}(2|2)\right)$$
where $`S`$ denotes the antipode. On $`U_q\left(gl(2|2)\right)`$, the antipode is defined by
(105)
$$S(E_0^{k,+})=q^{h_k}E_0^{k,+}S(E_0^{k,})=E_0^{k,}q^{h_k}S(h_l)=h_l$$
for $`k=1,2,3`$ and $`l=1,2,3,4`$. Its definition on $`U_q\left(\widehat{gl}(2|2)\right)`$ can be found in . For two vertex operators (103), the commutation relation reads
(106)
$$\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}^{}V_I}(z)\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}^{}V_I}(w)=\overline{R}_{\stackrel{̊}{V}^{}\stackrel{̊}{V}^{}}\left(\frac{z}{w}\right)\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}^{}V_I}(w)\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}^{}V_I}(z)$$
In the following, the index pairs of the R-matrix elements will be indicated by $`j_i,t_i`$ if associated with $`\stackrel{̊}{V}`$ and by $`j_i,t_i^{}`$ if associated to $`\stackrel{̊}{V}^{}`$. Then the subscript specifying $`V^{(1)}V^{(2)}`$ in (79) and (80) can be skipped. The R-matrix elements of $`\overline{R}_{\stackrel{̊}{V}^{}\stackrel{̊}{V}^{}}(z)`$ can be obtained from the matrix elements of $`\overline{R}_{\stackrel{̊}{V}\stackrel{̊}{V}}(z)`$ by means of the relation
(107)
$$\overline{R}_{j_1,t_1^{};j_2,t_2^{}}^{j_3,t_3^{};j_4,t_4^{}}(z)=\overline{R}_{j_3,t_3;j_4,t_4}^{j_1,t_1;j_2,t_2}(z)$$
Explicite expressions for the matrix elements of $`\overline{R}_{\stackrel{̊}{V}\stackrel{̊}{V}}(z)`$ are found in appendix B. In order to obtain well-defined limits for the R-matrix elements in the limit $`q0`$, the basis of the modules $`V`$ and $`V^{}`$ needs to be changed slightly. A new basis $`\{\stackrel{~}{v}_{j,t}\}_{0j3}`$ or $`\{\stackrel{~}{v}_{j,t}^{}\}_{0j3}`$ is introduced by
(108)
$$\stackrel{~}{v}_{j,t}=q^{t(t+\alpha _j)}v_{j,t}\stackrel{~}{v}_{j,t}^{}=q^{t(t+\alpha _j)}v_{j,t}^{}$$
with $`\alpha _j=\frac{1}{2}+\delta _{j,1}\delta _{j,2}`$. Then the expressions collected in formulae (163)-(187) yield
(109)
$$\underset{q0}{lim}\stackrel{~}{R}_{j_1,t_1;j_2,t_2}^{j_3,t_3;j_4,t_4}(z)=\delta _{j_1,j_4}\delta _{j_2,j_3}\delta _{t_1,t_4}\delta _{t_2,t_3}z^{t_2+\theta _{j_2,j_1}}$$
Here $`\stackrel{~}{R}_{j_1,t_1;j_2,t_2}^{j_3,t_3;j_4,t_4}(z)`$ denotes the R-matrix elements in the basis $`\{\stackrel{~}{v}_{j,t}\}_{0j3}`$ as given by (189) and $`\theta _{j_2,j_1}`$ is defined by (54). Changing the basis according to (108) does not affect the expression (98). With respect to the basis $`\{\stackrel{~}{v}_{j,t}\}_{0j3,t_0}`$ or $`\{\stackrel{~}{v}_{j,t}^{}\}_{0j3,t_0}`$, the commutation relations (80) for $`V=\stackrel{̊}{V}`$ or $`V=\stackrel{̊}{V}^{}`$ with $`0j3`$, $`t_0`$ may be viewed as equations for objects with small $`q`$ expansions $`\left(\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}V_I}\right)_{j,t}(z^1)=\varphi _{j,t}(z)+O(q)`$ or $`\left(\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}^{}V_I}\right)_{j,t}(z^1)=\varphi _{j,t}^{}(z)+O(q)`$. Then in the limit of vanishing $`q`$ the commutation relations reads
(110)
$$\begin{array}{cc}\hfill \varphi _{j_1,t_1}(z)\varphi _{j_2,t_2}(w)& =\left(\frac{z}{w}\right)^{t_2\theta _{j_2,j_1}}\varphi _{j_1,t_1}(w)\varphi _{j_2,t_2}(z)\hfill \\ \hfill \varphi _{j_1,t_1}^{}(z)\varphi _{j_2,t_2}^{}(w)& =\left(\frac{z}{w}\right)^{t_1\theta _{j_1,j_2}}\varphi _{j_1,t_1}^{}(w)\varphi _{j_2,t_2}^{}(z)\hfill \end{array}$$
With the mode expansions $`\varphi _{j,t}(z)=_m\varphi _{j,t;m}z^m`$, $`\varphi _{j,t}^{}(z)=_m\varphi _{j,t;m}^{}z^m`$, the relations (110) yield the defining equations (53), (65) for the creation algebras given in section 5:
(111)
$$\begin{array}{cc}\hfill \varphi _{j_1,t_1;m_1}\varphi _{j_2,t_2;m_2}& =\varphi _{j_1,t_1;m_2+t_2+\theta _{j_2,j_1}}\varphi _{j_2,t_2;m_1t_2\theta _{j_2,j_1}}\hfill \\ \hfill \varphi _{j_1,t_1;m_1}^{}\varphi _{j_2,t_2;m_2}^{}& =\varphi _{j_1,t_1;m_2+t_1+\theta _{j_1,j_2}}^{}\varphi _{j_2,t_2;m_1t_1\theta _{j_1,j_2}}^{}\hfill \end{array}$$
The properties (59) and (69) follow from (89). $`\mathrm{\Phi }_{1,1;0}^{}`$ $`\left(\mathrm{\Phi }_{1,1;0}\right)`$ can be viewed as the component corresponding to the unique semi-standard super tableau in $`_{1,1}^h`$ $`(_{1,1}^v)`$ associated with the border strip consisting of a single row (column).
A free boson realization for $`\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}^{}V_I}(z)`$ is obtained employing the coproduct of the antipode of the $`U_q\left(\widehat{gl}(2|2)\right)`$-generators. Replacing $`q`$ by $`q^1`$ on the rhs of (14) gives a free boson realization of $`q^mS\left(H_m^l\right)`$. On the lhs of (9), $`X^{1,\pm }(z)`$ is replaced by $`X^{2,\pm }(z)`$ and vice versa. Then substitution of $`q`$ by $`q^1`$ on the rhs of (8) yields a suitable free boson realization of $`q^{\pm h_k}S\left(E^{k,\pm }(q^1z)\right)`$. Coproduct and antipode of the algebra satisfy $`\mathrm{\Delta }S(a)=(SS)\mathrm{\Delta }^{}(a)`$. Due to this property, the coproduct formulae (77) give sufficient information to provide a free boson expression for $`\left(\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}^{}V_I}\right)_{2,0}(z)`$ analogous to (98). The result yields the commutation relation
(112)
$$\left(\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}^{}V_I}\right)_{2,0}(z)\left(\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}^{}V_I}\right)_{2,0}(w)=\left(\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}^{}V_I}\right)_{2,0}(w)\left(\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}^{}V_I}\right)_{2,0}(z)$$
Together with (107), this leads to (106).
The relation with the creation algebras suggests that the particular vertex operators discussed above create eigenstates of the row-to-row transfer matrix.
###### Conjecture 3.
The type II vertex operators $`\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}V_I}(z)`$ and $`\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}^{}V_I}(z)`$ provide a part of the quasi-particle structure of the vertex model.
Further details will be given together with a more complete account in a forthcoming publication.
The limiting behaviour of the R-matrix elements $`\stackrel{~}{R}_{j_1,t_1;j_2,t_2}^{j_3,t_3;j_4,t_4}(z)`$ specified in (109) reflects the eigenvalue structure of the intertwiner $`\stackrel{ˇ}{R}_{VV}(z)=P^{gr}R_{VV}(z)`$. The decomposition of the tensor product $`VV`$ into irreducible $`U_q\left(gl(2|2)\right)`$-modules is given by
(113)
$$VV=V_{2\overline{\mathrm{\Lambda }}_2+2\overline{\mathrm{\Lambda }}_4}\left(\underset{n=0}{\overset{\mathrm{}}{}}V_{(n+1)\overline{\mathrm{\Lambda }}_1+(n+2)\overline{\mathrm{\Lambda }}_2+(n+1)\overline{\mathrm{\Lambda }}_3+(2n)\overline{\mathrm{\Lambda }}_4}\right)$$
Using the explicit expressions for the matrix elements listed in appendix B, the eigenvalues of $`\stackrel{ˇ}{R}_{VV}(z)`$ are readily evaluated. On any vector $`_{j_1,j_2,t_1,t_2}\varkappa _{j_1,t_1;j_2,t_2}v_{j_1,t_1}v_{j_2,t_2}`$ contained in the component $`V_{(n+1)\overline{\mathrm{\Lambda }}_1+(n+2)\overline{\mathrm{\Lambda }}_2+(n+1)\overline{\mathrm{\Lambda }}_3+(2n)\overline{\mathrm{\Lambda }}_4}`$, the linear operator $`\stackrel{ˇ}{R}_{VV}(z)`$ acts as the scalar
(114)
$$r_n(z)=\frac{zq^{2(n+1)}}{1q^{2(n+1)}z}\mathrm{}\frac{zq^4}{1q^4z}\frac{zq^2}{1q^2z}$$
Here the normalisation is chosen such that $`\stackrel{ˇ}{R}_{VV}(z)`$ acts as the unit on the component $`V_{2\overline{\mathrm{\Lambda }}_2+2\overline{\mathrm{\Lambda }}_4}`$. The $`U_q\left(gl(2|2)\right)`$-weight space of the component $`V_{(n+1)\overline{\mathrm{\Lambda }}_1+(n+2)\overline{\mathrm{\Lambda }}_2+(n+1)\overline{\mathrm{\Lambda }}_3+(2n)\overline{\mathrm{\Lambda }}_4}`$ is in one-to-one correspondence with the pairs $`(j_1,t_1;j_2,t_2)`$ attributed to the power $`z^n`$ by (109). Hence the $`U_q\left(gl(2|2)\right)`$-weights attributed to all normal forms $`\varphi _{j_2,t_2;m+n+1}\varphi _{j_1,t_1;m}`$ with $`n>0`$ and $`m`$ fixed exactly corresponds to the $`U_q\left(gl(2|2)\right)`$-weights of the components $`V_{2\overline{\mathrm{\Lambda }}_2+2\overline{\mathrm{\Lambda }}_4}`$ and $`V_{(n^{}+1)\overline{\mathrm{\Lambda }}_1+(n^{}+2)\overline{\mathrm{\Lambda }}_2+(n^{}+1)\overline{\mathrm{\Lambda }}_3+(2n^{})\overline{\mathrm{\Lambda }}_4}`$ with $`n^{}<n`$. The $`U_q\left(gl(2|2)\right)`$-weights of $`\varphi _{j_r,t_r;m+r1}\mathrm{}\varphi _{j_2,t_2;t_2}\varphi _{j_1,t_1;m_1}`$ correspond to the weights of $`V_{r\overline{\mathrm{\Lambda }}_2+r\overline{\mathrm{\Lambda }}_4}`$. Generally, a finite set of infinite-dimensional $`U_q\left(gl(2|2)\right)`$-modules can be associated with each infinite vertical or horizontal border strip with parameters $`r_0`$, $`(p_1,\mathrm{},p_{R1};\overline{p}_1,\mathrm{},\overline{p}_R)`$ such that the weights found in the modules coincide with the weights of the normal forms in the set $`B(p_1,\mathrm{},p_{R1};\overline{p}_1,\mathrm{},\overline{p}_R;r_0)`$. Analogous statements apply to horizontal border strips.
In contrast, the infinite configuration space of the vertex model does not seem to indicate the structure of $`VV^{}`$ or $`V^{}V`$ in an analogous way. Within a certain region of the spectral parameter $`z`$, the action of the R-matrix on the tensor product $`VV^{}`$ is analysed in the following section.
## 7. The mixed case
The description of the space of infinite configurations outlined in sections 4 and 5 is achieved by means of two sets of normal forms. Each set is provided by a creation algebra and applies to the description of one space of components. The two creation algebras can be taken independent. As fare as described by the previous section, the limiting procedure connecting the type II vertex operators to the creation algebras treats both pairs $`(\mathrm{\Phi }_{V_{I_1}}^{VV_{J_1}}(z),𝒜)`$ and $`(\mathrm{\Phi }_{V_{I_2}}^{V^{}V_{J_2}}(w),𝒜^{})`$ separately. This raises the questions whether combining both vertex operators for the description of the entire configuration space is compatible with a description by two independent creation algebras. Assuming that $`V_zV_w^{}`$ and $`V_w^{}V_z`$ are intertwined by $`R_{VV^{}}\left(\frac{z}{w}\right)`$ according to (78), the commutation relations involving both types of vertex operators can be written in the form (80) with $`V^{(1)}=V`$ and $`V^{(2)}=V^{}`$. However, this choice of evaluation modules does not allow for the introduction of an intertwiner in the general sense. Though the intertwining conditions (78) can be solved formally for the single R-matrix elements (79), the corresponding action of the R-matrix on the $`U_q\left(gl(2|2)\right)`$-irreducible components of $`VV^{}`$ is well defined only for $`|\frac{q^2w}{z}|<1`$. In the following, it will be shown that within this region, this action does not depend on the irreducible component in the limit $`q0`$.
In certain cases, consideration of single matrix elements of $`\stackrel{ˇ}{R}_{V^{(1)}V^{(2)}}(z)=P^{gr}R_{V^{(1)}V^{(2)}}(z)`$ is less instructive in the case $`V^{(1)}V^{(2)}`$. The normalisation of a vector in an irreducible component of $`V^{(1)}V^{(2)}`$ can be chosen independently of the normalisation of a vector in the corresponding component of $`V^{(2)}V^{(1)}`$. Choices of the normalisations depend on the purpose, e.g. taking a certain limit, and is not apparent from a single matrix element. In fact, it seems hard to extract some formal limit of vanishing $`q`$ from the equations (80) with $`V^{(1)}=V`$ and $`V^{(2)}=V^{}`$.
The decomposition of $`VV^{}`$ can be found studying the action of the quadratic Casimir $`C^{(2)}`$ on the tensor products.
The operator $`C^{(2)}`$ given by
(115)
$$\begin{array}{cc}& \frac{C^{(2)}}{(qq^1)^2}=\hfill \\ & E_0^{1,}E_0^{1,+}q^{h_2+h_3+h_4}q^2E_0^{2,}E_0^{2,+}q^{h_1+h_3+h_4}E_0^{3,}E_0^{3,+}q^{h_1h_2+h_4}\hfill \\ & +q^1[E_0^{1,},E_0^{2,}]_{q^1}[E_0^{1,+},E_0^{2,+}]_qq^{h_3+h_4}q^1[E_0^{2,},E_0^{3,}]_{q^1}[E_0^{2,+},E_0^{3,+}]_qq^{h_1+h_4}\hfill \\ & (E_0^{3,}[E_0^{1,},E_0^{2,}]_{q^1}q^1[E_0^{1,},E_0^{2,}]_{q^1}E_0^{3,})\hfill \\ & \left(E_0^{3,+}[E_0^{1,+},E_0^{2,+}]_qq[E_0^{1,+},E_0^{2,+}]_qE_0^{3,+}\right)q^{h_4}\hfill \\ & +\frac{[2]}{(qq^1)^2}\frac{q^{h_41}}{(qq^1)^2}\left(q^{h_1+h_2+h_3}q^{h_1+h_2+h_3}+q^{h_1h_2+h_3}q^{h_1h_2h_3}\right)\hfill \end{array}$$
is a central element of $`U_q\left(gl(2|2)\right)`$. Here the $`q`$-deformed commutators are defined by
(116)
$$\begin{array}{cc}\hfill [E_0^{i,+},E_0^{i+1,+}]_q& =E_0^{i,+}E_0^{i+1,+}+q^{(1)^{i+1}}E_0^{i+1,+}E_0^{i,+}\hfill \\ \hfill [E_0^{i,},E_0^{i+1,}]_{q^1}& =E_0^{i,}E_0^{i+1,}+q^{(1)^i}E_0^{i+1,}E_0^{i,}\hfill \end{array}$$
Generalised eigenvectors of $`\mathrm{\Delta }\left(C^{(2)}\right)`$ with eigenvalue $`2y`$ are provided by
(117) $`{\displaystyle \underset{t=0}{\overset{\mathrm{}}{}}}q^{t(2r1)}{\displaystyle \frac{(q^2;q^2)_{t+1}(q^2;q^2)_r}{(q^2;q^2)_{t+r}(q^2;q^2)_1}}\stackrel{~}{P}_t(y;q^{2r+1},q,q|q^2)v_{2,t+r}v_{0,t}^{}`$ $`r0`$
(118) $`{\displaystyle \underset{t=0}{\overset{\mathrm{}}{}}}q^{t(2r+1)}{\displaystyle \frac{(q^2;q^2)_{t+r+1}}{(q^2;q^2)_t(q^2;q^2)_{r+1}}}\stackrel{~}{P}_t(y;q^{2r+1},q,q|q^2)v_{2,t}v_{0,t+r}^{}`$ $`r>0`$
and
(119) $`{\displaystyle \underset{t=0}{\overset{\mathrm{}}{}}}q^{t(2r1)}{\displaystyle \frac{(q^2;q^2)_{t+1}(q^2;q^2)_r}{(q^2;q^2)_{t+r}(q^2;q^2)_1}}\stackrel{~}{P}_t(y;q^{2r1},q^1,q^1|q^2)v_{0,t}^{}v_{2,t+r}`$ $`r0`$
(120) $`{\displaystyle \underset{t=0}{\overset{\mathrm{}}{}}}q^{t(2r+1)}{\displaystyle \frac{(q^2;q^2)_{t+r+1}}{(q^2;q^2)_t(q^2;q^2)_{r+1}}}\stackrel{~}{P}_t(y;q^{2r1},q^1,q^1|q^2)v_{0,t+r}^{}v_{2,t}`$ $`r>0`$
where $`r0`$, $`(q^2;q^2)_n=(1q^2)(1q^4)\mathrm{}(1q^{2n})`$ and $`\stackrel{~}{P}_t(y;a,c,d|q^2)\stackrel{~}{P}_t(y)`$ denote orthonormal continuous dual Hahn polynomials in base $`q^2`$ . They satisfy the recurrence relation
(121)
$$2y\stackrel{~}{P}_t(y)=a_t\stackrel{~}{P}_{t+1}(y)+b_t\stackrel{~}{P}_t(y)+a_{t1}\stackrel{~}{P}_{t1}(y)$$
where
(122)
$$\begin{array}{cc}\hfill a_t& =\sqrt{\left(1q^{2(t+1)}\right)\left(1acq^{2t}\right)\left(1adq^{2t}\right)\left(1cdq^{2t}\right)}\hfill \\ \hfill b_t& =a+a^1a^1\left(1acq^{2t}\right)\left(1adq^{2t}\right)a\left(1q^{2t}\right)\left(1cdq^{2(t1)}\right)\hfill \end{array}$$
and $`t1`$, $`\stackrel{~}{P}_0(y)=1`$ and $`\stackrel{~}{P}_1(y)=0`$. Replacing $`q`$ by $`q^1`$ in the coefficients (122) yields the recursion relations for the polynomials $`\stackrel{~}{P}_t(y;a,c,d|q^2)`$. They are supplemented by the initial condition $`\stackrel{~}{P}_1(y;a,c,d|q^2)=0`$ and a convenient choice of the value of $`\stackrel{~}{P}_0(y;a,c,d|q^2)`$ to be specified below. The polynomials $`\stackrel{~}{P}_t(y;a,c,d|q^2)`$ are orthonormal with respect to the measure $`dm(;a,0,c,d|q^2)`$. The measure $`dm(;a,b,c,d|q^2)`$ is defined by
(123)
$$_{}f(y)𝑑m(y;a,b,c,d|q^2)=_0^\pi f(\mathrm{cos}\theta )w(\mathrm{cos}\theta )𝑑\theta +\underset{k}{}f(y_k)w_k$$
with
$$\begin{array}{c}w(\mathrm{cos}\theta )=w(\mathrm{cos}\theta ;a,b,c,d|q^2)\hfill \\ \hfill =\frac{1}{2\pi }\frac{(q^2,ab,ac,ad,bc,bd,cd;q^2)_{\mathrm{}}}{(abcd;q^2)_{\mathrm{}}}\frac{(e^{2i\theta },e^{2i\theta };q^2)_{\mathrm{}}}{(ae^{i\theta },ae^{i\theta },be^{i\theta },be^{i\theta },ce^{i\theta },ce^{i\theta },de^{i\theta },de^{i\theta };q^2)_{\mathrm{}}}\end{array}$$
and $`y_k=\frac{1}{2}(eq^{2k}+e^1q^{2k})`$ for $`e`$ any of the parameters $`a,b,c,d`$. The sum includes all $`k_0`$ with the property $`|eq^{2k}|>1`$. For $`e=b`$, the coefficients are given by
(124)
$$\begin{array}{c}w_k=w_k(b;a,c,d|q^2)=\hfill \\ \hfill \frac{1b^2q^{4k}}{1b^2}\frac{(b^2,ac,ad,cd;q^2)_{\mathrm{}}}{(a/b,c/b,d/b,abcd;q^2)_{\mathrm{}}}\frac{(b^2,ab,bc,bd;q^2)_k}{(q^2,q^2b/a,q^2b/c,q^2b/d;q^2)_k}\left(\frac{q^2}{abcd}\right)^k\end{array}$$
Throughout the following, the parameters $`a,c,d`$ satisfy the condition $`max(|a|,|c|,|d|)<1`$. Thus the measure $`dm(;a,0,c,d|q^2)`$ contains only the continuous part. In (123) and (124), the notation for the $`q`$-shifted factorials is adopted from :
(125) $`(a_1,\mathrm{},a_s;q^2)_{\mathrm{}}`$ $`=(a_1;q^2)_{\mathrm{}}\mathrm{}(a_s;q^2)_{\mathrm{}},`$ $`(a;q^2)_{\mathrm{}}`$ $`={\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}\left(1aq^{2j}\right)`$
(126) $`(a_1,\mathrm{},a_s;q^2)_n`$ $`=(a_1;q^2)_n\mathrm{}(a_s;q^2)_n,`$ $`(a;q^2)_n`$ $`={\displaystyle \underset{j=0}{\overset{n1}{}}}\left(1aq^{2j}\right)`$
(127) $`(a_1,\mathrm{},a_s;q^2)_0`$ $`=1`$
The measure defined in (123) provides an orthogonality measure for the Askey-Wilson polynomials $`p_t(y;a,b,c,d|q^2)`$ introduced by
(128)
$$\begin{array}{c}p_t(y;a,b,c,d|q^2)=\hfill \\ \hfill a^t(ab,ac,ad;q^2)_t{}_{4}{}^{}\varphi _{3}^{}(q^{2t},abcdq^{2(t1)},a\gamma ,a\gamma ^1;ab,ac,ad;q^2,q^2)\end{array}$$
with $`y=\frac{1}{2}(\gamma +\gamma ^1)`$. Here the basic hypergeometric series is defined by
(129)
$${}_{n+1}{}^{}\varphi _{n}^{}(a_1,a_1,\mathrm{},a_{n+1};b_1,\mathrm{},b_n;q^2,\alpha )=\underset{t=0}{\overset{\mathrm{}}{}}\frac{(a_1,a_2,\mathrm{},a_{n+1};q^2)_t}{(q^2,b_1,\mathrm{},b_n;q^2)_t}\alpha ^t$$
Due to Sears’ transformation (, eq. 2.10.4), the polynomials $`p_t(y;a,b,c,d|q^2)`$ are symmetric in $`a,b,c,d`$. The orthonormal continuous dual $`q^2`$-Hahn polynomials and the Askey-Wilson polynomials are related by
(130)
$$\stackrel{~}{P}_t(y;a,b,c,d|q^2)=\frac{1}{(q^2,ac;q^2)_t}p_t(y;a,0,c,d|q^2)$$
The formal intertwining conditions on the R-matrix elements introduced on $`V_zV_w^{}`$ or $`V_z^{}V_w`$ are given by (78) and (79) with $`V^{(1)}=V`$, $`V^{(2)}=V^{}`$ or $`V^{(1)}=V^{}`$, $`V^{(2)}=V`$. With the $`U_q\left(gl(2|2)\right)`$-weight structure of $`V`$ and $`V^{}`$ specified by (81) and (104), these conditions imply
(131)
$$\begin{array}{cc}\hfill \stackrel{ˇ}{R}_{VV^{}}\left(\frac{z}{w}\right)(v_{2,t_1}v_{0,t_2}^{})& =\underset{t_3,t_4}{}\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}0},t_2^{}}^{0,t_4^{};\mathrm{\hspace{0.17em}2},t_3}\left(\frac{z}{w}\right)v_{0,t_4}^{}v_{2,t_3}\hfill \\ \hfill \stackrel{ˇ}{R}_{V^{}V}\left(\frac{w}{z}\right)(v_{0,t_2}^{}v_{2,t_1})& =\underset{t_3,t_4}{}\stackrel{ˇ}{R}_{0,t_2^{};\mathrm{\hspace{0.17em}2},t_1}^{2,t_3;\mathrm{\hspace{0.17em}0},t_4^{}}\left(\frac{w}{z}\right)v_{2,t_3}v_{0,t_4}^{}\hfill \end{array}$$
where $`t_1t_2=t_3t_4`$. Equations (159) and (192) imply
(132)
$$\begin{array}{cc}\hfill \stackrel{ˇ}{R}_{0,t_2^{};\mathrm{\hspace{0.17em}2},t_2+r}^{2,t_1+r;\mathrm{\hspace{0.17em}0},t_1^{}}(z)& =q^{2(t_1t_2)(2r1)}\left(\frac{(q^2;q^2)_{t_2+r}(q^2;q^2)_{t_1+1}}{(q^2;q^2)_{t_1+r}(q^2;q^2)_{t_2+1}}\right)^2\stackrel{ˇ}{R}_{2,t_1+r;\mathrm{\hspace{0.17em}0},t_1^{}}^{0,t_2^{};\mathrm{\hspace{0.17em}2},t_2+r}(z)\hfill \\ \hfill \stackrel{ˇ}{R}_{0,t_2+r^{};\mathrm{\hspace{0.17em}2},t_2}^{2,t_1;\mathrm{\hspace{0.17em}0},t_1+r^{}}(z)& =q^{2(t_2t_1)(2r+1)}\left(\frac{(q^2;q^2)_{t_1+r+1}(q^2;q^2)_{t_2}}{(q^2;q^2)_{t_2+r+1}(q^2;q^2)_{t_1}}\right)^2\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}0},t_1+r^{}}^{0,t_2+r^{};\mathrm{\hspace{0.17em}2},t_2}(z)\hfill \end{array}$$
The matrix elements on the rhs of (131) may be viewed as the $`\stackrel{ˇ}{R}`$-matrix elements associated with two evaluation modules of $`U_q\left(\widehat{sl}(2)\right)`$. A convenient choice of $`U_q\left(\widehat{sl}(2)\right)`$-generators is
(133)
$$\begin{array}{cc}\hfill \widehat{h}_1& =h_2+h_3\hfill \\ \hfill \widehat{h}_0& =h_0+h_1\hfill \\ \hfill \widehat{e}_1& =\left(E_0^{3,}E_0^{2,}+q^1E_0^{2,}E_0^{3,}\right)q^{h_2+h_3}\hfill \\ \hfill \widehat{f}_1& =q^{h_2h_3}\left(E_0^{2,+}E_0^{3,+}+qE_0^{3,+}E_0^{2,+}\right)\hfill \\ \hfill S(\widehat{e}_0)& =q^{h_2+h_3}\left(E_0^{2,+}E_1^{3,+}+q^1E_1^{3,+}E_0^{2,+}\right)\hfill \\ \hfill S(\widehat{f}_0)& =\left(E_1^{3,}E_0^{2,}+qE_0^{2,}E_1^{3,}\right)q^{h_2h_3}\hfill \end{array}$$
The defining relations obeyed by the generators $`\widehat{e}_0,\widehat{e}_1,\widehat{f}_0,\widehat{f}_1,\widehat{h}_0,\widehat{h}_1`$ (see , for example) are inherited from the defining relations of $`U_q\left(\widehat{sl}(2|2)\right)`$ . Two infinite-dimensional $`U_q\left(sl(2)\right)`$-modules $`\widehat{V}`$ and $`\widehat{V}^{}`$ with basis $`\{\widehat{v}_j\}_{j_0}`$ and $`\{\widehat{v}_j^{}\}_{j_0}`$ are introduced by
(134)
$$\widehat{v}_t=q^{\frac{1}{2}t(2t+1)}\frac{(q^2;q^2)_{t+1}}{(1q^2)^{t+1}}v_{0,t}^{}\widehat{v}_t^{}=q^{\frac{1}{2}t(2t1)}\frac{(1q^2)^t}{(q^2;q^2)_t}v_{2,t}$$
Then (133) and (81), (83) yield $`U_q\left(\widehat{sl}(2)\right)`$-structures on $`\widehat{V}`$ and $`\widehat{V}^{}`$. With these structures, $`\{v_{2,t}z^m\}_{t_0,m}`$ and $`\{v_{0,t}^{}w^m\}_{t_0,m}`$ can be regarded as $`U_q\left(\widehat{sl}(2)\right)`$-evaluation modules $`\widehat{V}_z^{}`$ and $`\widehat{V}_w^{}^{}`$. With the definitions of the $`U_q\left(\widehat{sl}(2)\right)`$-structures on $`\widehat{V}_z^{}`$ and $`\widehat{V}_w^{}`$ chosen analogous to (72), the spectral parameters are related by $`z^{}=q^3z^1`$, $`w^{}=q^3w^1`$. The $`U_q\left(\widehat{sl}(2)\right)`$-structures of $`\widehat{V}`$ and $`\widehat{V}^{}`$ read
(135) $`\widehat{e}_0\widehat{v}_t`$ $`=q^{t\frac{3}{2}}[t+1]\widehat{v}_{t+1}`$ $`\widehat{e}_0\widehat{v}_t^{}`$ $`=q^{t\frac{5}{2}}[t]\widehat{v}_{t1}^{}`$
$`\widehat{f}_0\widehat{v}_t`$ $`=q^{t+\frac{5}{2}}[t]\widehat{v}_{t1}`$ $`\widehat{f}_0\widehat{v}_t^{}`$ $`=q^{t+\frac{7}{2}}[t+1]\widehat{v}_{t+1}^{}`$
$`\widehat{h}_0\widehat{v}_t`$ $`=(2t+1)\widehat{v}_t`$ $`\widehat{h}_0\widehat{v}_t^{}`$ $`=(2t+1)\widehat{v}_t^{}`$
$`\widehat{e}_1\widehat{v}_t`$ $`=q^{t\frac{1}{2}}[t]\widehat{v}_{t1}`$ $`\widehat{e}_1\widehat{v}_t^{}`$ $`=q^{t+\frac{1}{2}}[t+1]\widehat{v}_{t+1}^{}`$
$`\widehat{f}_1\widehat{v}_t`$ $`=q^{t+\frac{3}{2}}[t+1]\widehat{v}_{t+1}`$ $`\widehat{f}_1\widehat{v}_t^{}`$ $`=q^{t+\frac{1}{2}}[t]\widehat{v}_{t1}^{}`$
$`\widehat{h}_1\widehat{v}_t`$ $`=(2t+1)\widehat{v}_t`$ $`\widehat{h}_1\widehat{v}_t^{}`$ $`=(2t+1)\widehat{v}_t^{}`$
Apart from simple factors depending on $`r`$, the vectors (117), (119) can be rewritten by
(136)
$$\underset{t=0}{\overset{\mathrm{}}{}}\stackrel{~}{P}_t(y;q^{2r+1},q,q|q^2)\widehat{v}_t^{}\widehat{v}_{t+r},\underset{t=0}{\overset{\mathrm{}}{}}\stackrel{~}{P}_t(y;q^{2r+1},q,q|q^2)\widehat{v}_{t+r}^{}\widehat{v}_t$$
and
(137)
$$\underset{t=0}{\overset{\mathrm{}}{}}\stackrel{~}{P}_t(y;q^{2r1},q^1,q^1|q^2)\widehat{v}_{t+r}\widehat{v}_t^{},\underset{t=0}{\overset{\mathrm{}}{}}\stackrel{~}{P}_t(y;q^{2r1},q^1,q^1|q^2)\widehat{v}_t\widehat{v}_{t+r}^{}$$
They are generalised eigenvectors of the Casimir element of $`U_q\left(sl(2)\right)`$ given by
(138)
$$C=(qq^1)^2\widehat{f}_1\widehat{e}_1+q^{\widehat{h}_1+1}+q^{\widehat{h}_11}$$
with eigenvalue $`2y`$. For fixed $`y=\mathrm{cos}\theta `$, the vectors (136) with $`r=0,1,2,\mathrm{}`$ constitute an infinite-dimensional $`U_q\left(sl(2)\right)`$-module $`\widehat{V}^{(y)}`$ corresponding to the principal unitary series representation $`\pi _{\rho ,0}^P`$ of $`U_q\left(su(1,1)\right)`$ with $`\mathrm{cos}\theta =\frac{1}{2}\left(q^{2i\rho }+q^{2i\rho }\right)`$ , . Each representation $`\pi _{\rho ,0}^P`$ is irreducible . The tensor product $`\widehat{V}^{}\widehat{V}`$ decomposes into an integral over of all modules $`\widehat{V}^{(\mathrm{cos}\theta )}`$ with $`0\theta \pi `$ , .
The vectors in (117) have the $`U_q\left(gl(2|2)\right)`$-weights $`(r+1,r+1,r1,r+1)`$ with $`r`$. Acting with the coproducts of the $`U_q\left(gl(2|2)\right)`$-generators on them produces further generalised eigenvectors of the quadratic Casimir $`C^{(2)}`$. Taking into account the defining relations of $`U_q\left(gl(2|2)\right)`$ (see ) and the explicite expression (115) for $`C^{(2)}`$, the set of all resulting linear independent vectors for fixed $`y`$ can be specified. They provide an irreducible $`U_q\left(gl(2|2)\right)`$-module $`V^{(y)}`$. The weights of the vectors composing $`V^{(y)}`$ are listed by
(139)
$$\begin{array}{cc}& (r+1,r+1,r1,r+1),(r+1,r1,r1,r+3),\hfill \\ & (r+1,r3,r1,r+1),(r+1,r1,r1,r1),\hfill \\ & (r+1,r,r1,r+2)_2,(r+1,r2,r1,r+2)_2,\hfill \\ & (r+1,r2,r1,r)_2,(r+1,r,r1,r)_2,\hfill \\ & (r+1,r1,r1,r+1)_4.\hfill \end{array}$$
Here the subscripts denotes the multiplicity of generalised eigenvectors at the given weight.
A decomposition of the tensor product $`VV^{}`$ contains all $`U_q\left(gl(2|2)\right)`$-modules $`V^{(\mathrm{cos}\theta )}`$ with $`0\theta \pi `$ since the eigenvectors related to the first set of weights in (139) can be regarded as the set of vectors spanning the $`U_q\left(sl(2)\right)`$-module $`\widehat{V}^{(y)}`$. Inspection of $`U_q\left(gl(2|2)\right)`$-weight structure of $`VV^{}`$ and the list (139) shows that no further modules occur in the decomposition of $`VV^{}`$.
The intertwining condition (78) should apply with the particular choice $`a=C^{(2)}`$. Thus, if the action of $`\stackrel{ˇ}{R}_{VV^{}}(z)`$ is well defined on the irreducible components of $`VV^{}`$, it is expected to take the form
(140)
$$\begin{array}{c}\stackrel{ˇ}{R}_{VV^{}}(z)\left(\underset{t=0}{\overset{\mathrm{}}{}}q^{t(2r+1)}\frac{(q^2;q^2)_{t+r+1}}{(q^2;q^2)_t(q^2;q^2)_{r+1}}\stackrel{~}{P}_t(y;q^{2r+1},q,q|q^2)v_{2,t}v_{0,t+r}^{}\right)=\hfill \\ \hfill r(y,z)\underset{t=0}{\overset{\mathrm{}}{}}q^{t(2r+1)}\frac{(q^2;q^2)_{t+r+1}}{(q^2;q^2)_t(q^2;q^2)_{r+1}}\stackrel{~}{P}_t(y;q^{2r1},q^1,q^1|q^2)v_{0,t+r}^{}v_{2,t}\end{array}$$
or
(141)
$$\begin{array}{c}\stackrel{ˇ}{R}_{VV^{}}(z)\left(\underset{t=0}{\overset{\mathrm{}}{}}q^{t(2r1)}\frac{(q^2;q^2)_{t+1}(q^2;q^2)_r}{(q^2;q^2)_{t+r}(q^2;q^2)_1}\stackrel{~}{P}_t(y;q^{2r+1},q,q|q^2)v_{2,t+r}v_{0,t}^{}\right)=\hfill \\ \hfill r(y,z)\underset{t=0}{\overset{\mathrm{}}{}}q^{t(2r1)}\frac{(q^2;q^2)_{t+1}(q^2;q^2)_r}{(q^2;q^2)_{t+r}(q^2;q^2)_1}\stackrel{~}{P}_t(y;q^{2r1},q^1,q^1|q^2)v_{0,t}^{}v_{2,t+r}\end{array}$$
with $`|y|=1`$ and a suitable normalisation of the polynomials $`\stackrel{~}{P}_t(y;q^{2r1},q^1,q^1|q^2)`$. Within a certain range of spectral parameters $`z`$, the function $`r(y,z)`$ can be expressed as an infinite sum over elements of $`R_{VV}(q^2z^1)`$ with the continuous dual $`q^2`$-Hahn polynomials as coefficients. If $`|qz^1|<1`$, the sum is absolutely convergent and can be evaluated by means of a sum formula established in (see Appendix C):
(142)
$$r(\mathrm{cos}\theta ,z)=\left(1q^2z^1\right)\frac{(q^2z^1,q^2z^1;q^2)_{\mathrm{}}}{(qz^1e^{i\theta },qz^1e^{i\theta };q^2)_{\mathrm{}}},|qz^1|<1$$
where the normalisation $`\stackrel{~}{P}_0(y;q^{2r1},q^1,q^1|q^2)=q^r`$ has been chosen. Within the range of validity, the rhs of (142) does not possess zeros or poles. For $`|qz^1|<1`$, the inverse transformations read:
(143)
$$\begin{array}{c}\stackrel{ˇ}{R}_{2,t_1+r;\mathrm{\hspace{0.17em}0},t_1^{}}^{0,t_2^{};\mathrm{\hspace{0.17em}2},t_2+r}(z)=q^{(t_2t_1)(2r1)}\frac{(q^2;q^2)_{t_1+r}(q^2;q^2)_{t_2+1}}{(q^2;q^2)_{t_2+r}(q^2;q^2)_{t_1+1}}\hfill \\ \hfill _0^\pi \stackrel{~}{P}_{t_1}(\mathrm{cos}\theta ;q^{2r+1},q,q|q^2)\stackrel{~}{P}_{t_2}(\mathrm{cos}\theta ;q^{2r1},q^1,q^1|q^2)r(\mathrm{cos}\theta ,z)w(\mathrm{cos}\theta )𝑑\theta \end{array}$$
and
(144)
$$\begin{array}{c}\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}0},t_1+r^{}}^{0,t_2+r^{};\mathrm{\hspace{0.17em}2},t_2}(z)=q^{(t_1t_2)(2r+1)}\frac{(q^2;q^2)_{t_2+r+1}(q^2;q^2)_{t_1}}{(q^2;q^2)_{t_1+r+1}(q^2;q^2)_{t_2}}\hfill \\ \hfill _0^\pi \stackrel{~}{P}_{t_1}(\mathrm{cos}\theta ;q^{2r+1},q,q|q^2)\stackrel{~}{P}_{t_2}(\mathrm{cos}\theta ;q^{2r1},q^1,q^1|q^2)r(\mathrm{cos}\theta ,z)w(\mathrm{cos}\theta )𝑑\theta \end{array}$$
with $`r0`$, $`r(\mathrm{cos}\theta ,z)`$ given by (142) and $`w(\mathrm{cos}\theta )=w(\mathrm{cos}\theta ;q^{2r+1},0,q,q|q^2)`$ defined by (123). The last two factors in (143) and (144) may be rewritten as
(145)
$$\begin{array}{c}r(\mathrm{cos}\theta ,z)w(\mathrm{cos}\theta ;q^{2r+1},0,q,q|q^2)=\hfill \\ \hfill \left(1q^2z^1\right)(q^2z^1,q^2z^1;q^2)_{\mathrm{}}w(\mathrm{cos}\theta ;q^{2r+1},qz^1,q,q|q^2)\end{array}$$
According to (123), the function $`w(\mathrm{cos}\theta ;q^{2r+1},qz^1,q,q|q^2)`$ enters the orthogonality measure for the Askey-Wilson polynomials $`p_t(\mathrm{cos}\theta ;q^{2r+1},qz^1,q,q|q^2)`$. For $`|qz^1|<1`$, the orthogonality relation reads (, eqs. 7.5.15-7.5.17)
(146)
$$\begin{array}{c}_0^\pi p_t(\mathrm{cos}\theta ;q^{2r+1},qz^1,q,q|q^2)p_s(\mathrm{cos}\theta ;q^{2r+1},qz^1,q,q|q^2)w_z(\mathrm{cos}\theta )𝑑\theta =\hfill \\ \hfill \frac{\delta _{t,s}}{(q^2z^1,q^2z^1;q^2)_{\mathrm{}}h_t(q^{2r+1},qz^1,q,q|q^2)}\end{array}$$
where $`w_z(\mathrm{cos}\theta )=w(\mathrm{cos}\theta ;q^{2r+1},qz^1,q,q|q^2)`$ and
(147)
$$h_t(q^{2r+1},qz^1,q,q|q^2)=\frac{1q^{2(2t+r+1)}z^1}{(q^2,q^2,q^{2(r+1)},q^{2(r+1)},q^2z^1,q^2z^1;q^2)_t}$$
The continuous dual Hahn polynomials $`\stackrel{~}{P}_t(y;a,0,c,d|q^2)`$ and $`\stackrel{~}{P}_t(y;a^1,0,c^1,d^1|q^2)`$ may be expressed in terms of the polynomials $`p_t(y;q^{2r+1},qz^1,q,q|q^2)`$ making use of (130) and the connection coefficients between Askey-Wilson polynomials established in (see also section 7.6 in ). The required relation is a special case of eqs. 7.6.2 and 7.6.3 in :
(148)
$$\stackrel{~}{P}_t(y;a,c,d|q^2)=\underset{s=0}{\overset{t}{}}c_{s,t}p_s(y;a,b,c,d)$$
with
(149)
$$c_{s,t}=b^{ts}\frac{(cd;q^2)_t(acq^{2s};q^2)_{ts}}{(q^2,ad,cd,abcdq^{2(s1)};q^2)_s(q^2,abcdq^{4s};q^2)_{ts}}$$
Here the parameters $`a,b,c,d`$ are only restricted by the requirement that the denominators in the polynomials and coefficients never vanish. With $`a=q^{2r+1}`$, $`c=d=q`$ and $`b=qz^1`$, the relation reads
(150)
$$\stackrel{~}{P}_t(y;q^{2r+1},qz^1,q,q|q^2)=\underset{s=0}{\overset{t}{}}c_{s,t}(q,z)p_s(y;q^{2r+1},qz^1,q,q|q^2)$$
with
(151)
$$c_{s,t}(q,z)=\left(\frac{q}{z}\right)^{ts}\frac{(q^2;q^2)_t(q^{2(r+s+1)};q^2)_{ts}}{(q^{2(r+1)},q^2,q^2,q^{2(r+s+1)}z^1;q^2)_s(q^2,q^{2(r+2s+2)}z^1;q^2)_{ts}}$$
Replacing $`q`$ by $`q^1`$ and $`z`$ by $`z^1`$ in (150), (151) and changing from base $`q^2`$ to base $`q^2`$ in $`p_s(y;q^{2r1},q^1z,q^1,q^1|q^2)`$ yields
(152)
$$\frac{\stackrel{~}{P}_t(y;q^{2r1},q^1,q^1|q^2)}{\stackrel{~}{P}_0(y;q^{2r1},q^1,q^1|q^2)}=\underset{s=0}{\overset{t}{}}c_{s,t}^{}(q,z)p_s(y;q^{2r+1},qz^1,q,q|q^2)$$
with
(153)
$$c_{s,t}^{}(q,z)=z^sq^{s(3s+2r+1)}(1)^sc_{s,t}(q^1,z^1)$$
Here the normalisation of the polynomials $`\stackrel{~}{P}_t(y;q^{2r1},q^1,q^1|q^2)`$ is left arbitrary. Insertion of (146), (150) and (152) into equation (144) allows to express the R-matrix elements in terms of connection coefficients and the function specifying the corresponding orthogonality relation:
(154)
$$\begin{array}{c}\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}0},t_1+r^{}}^{0,t_2+r^{};\mathrm{\hspace{0.17em}2},t_2}(z)=q^{(t_1t_2)(2r+1)}\frac{(q^2;q^2)_{t_2+r+1}(q^2;q^2)_{t_1}}{(q^2;q^2)_{t_1+r+1}(q^2;q^2)_{t_2}}\stackrel{~}{P}_0(y;q^{2r1},q^1,q^1|q^2)\hfill \\ \hfill (1q^2z^1)\underset{t=0}{\overset{min(t_1,t_2)}{}}c_{t,t_1}(q,z)c_{t,t_2}^{}(q,z)\frac{1}{h_t(q^{2r+1},qz^1,q,q|q^2)}\end{array}$$
Equation (154) is valid for all values of $`z`$ provided that the R-matrix element on the lhs is well defined. The orthogonality relations (146) take a different form for $`|qz^1|>1`$ (see chapter 6 of ). Therefore the expression (142) does not give rise to inverse transformations of the form (143) or (144) if the condition $`|qz^1|<1`$ is not satisfied.
Within the region of validity, the expression (142) for the action of $`R_{VV^{}}(z)`$ has a well-defined limit as $`q`$ tends to zero. The latter does not show a structure which could distinguish the different irreducible components in the limit $`q0`$. Hence in this region the action does not indicate an inconsistency in the quasi-particle description of the space of infinite configurations in $`\mathrm{\Omega }_A^{(1)}`$ in terms of the two type II vertex operators $`\mathrm{\Phi }_{V_I}^{\stackrel{̊}{V}V_I}(z)`$ and $`\mathrm{\Phi }_{V_J}^{\stackrel{̊}{V}^{}V_J}(w)`$.
In contrast, the action of $`\stackrel{ˇ}{R}_{V^{}V}(z)`$ on the irreducible components (119) is well-defined for all $`zq^{2s}`$, $`s`$. The eigenvalue is finite at $`z=q^{2s1}e^{\pm i\theta }`$ for all $`s`$. Details will be given in a forthcoming publication.
## Appendix A The mapping for vertical border stripes
A generator $`\varphi _{j,t;r}`$ with suitable $`j,t`$ is assigned to the $`r`$-th box of a vertical border strip if this box carries the number $`0`$ in the reference labelling defined in section 4. Hence, the normal forms corresponding to a border strip with parameters $`r_0`$ and $`(p_1,\mathrm{},p_{R1};\overline{p}_1,\mathrm{},p_R)`$ or $`(\mathrm{};\overline{p}_1)`$ can be written
(155)
$$\begin{array}{cc}& \varphi _{j_{\overline{p}_1},t_{\overline{p}_1};s_1+\overline{p}_11}\mathrm{}\varphi _{j_2,t_2;s_1+1}\varphi _{j_1,t_1;s_1}\hfill \\ & \varphi _{j_{\overline{p}_1+\overline{p}_2},t_{\overline{p}_1+\overline{p}_2};s_2+\overline{p}_21}\mathrm{}\varphi _{j_{\overline{p}_1+2},t_{\overline{p}_1+2};s_2+1}\varphi _{j_{\overline{p}_1+1},t_{\overline{p}_1+1};s_2}\hfill \\ & \varphi _{j_{\overline{p}_1+\overline{p}_2+\overline{p}_3},t_{\overline{p}_1+\overline{p}_2+\overline{p}_3};s_3+\overline{p}_31}\mathrm{}\varphi _{j_{\overline{p}_1+\overline{p}_2+2},t_{\overline{p}_1+\overline{p}_2+2};s_3+1}\varphi _{j_{\overline{p}_1+\overline{p}_2+1},t_{\overline{p}_1+\overline{p}_2+1};,s_3}\hfill \\ & \mathrm{}\hfill \\ & \varphi _{j_{\overline{N}},t_{\overline{N}};s_R+\overline{p}_R1}\mathrm{}\varphi _{j_{\overline{p}_1+\mathrm{}+\overline{p}_{R1}+2},t_{\overline{p}_1+\mathrm{}+\overline{p}_{R1}+2};s_R+1}\varphi _{j_{\overline{p}_1+\mathrm{}+\overline{p}_{R1}+1},t_{\overline{p}_1+\mathrm{}+\overline{p}_{R1}+1};s_R}\hfill \end{array}$$
where $`\overline{N}=_{i=1}^R\overline{p}_i`$ and the numbers $`s_i`$ are defined in (49). Prescriptions for the indices $`j,t`$ are obtained from modifications of figures 3-5. Each section of a border strip drawn in figure 3 is transposed such that the uppermost box of the column becomes the leftmost box of a row. In figures 4 and 5, each section is transposed such that the leftmost box of the row becomes the uppermost box of a column. Throughout the figures, the $`p_i`$ is replaced by $`\overline{p}_i`$ for all $`i=1,2,\mathrm{}R`$ and $`\overline{p}_i`$ by $`p_i`$ for $`i=1,2,\mathrm{},R1`$. In the right part of the figures, each generator $`\varphi _{j,t;m}^{}`$ is replaced by $`\varphi _{j,t;m}`$. Then in figure 3 the order of the pairs $`j_i,t_i`$ in the product is reversed. The modes $`m_i`$ are kept unchanged. For example, the lowest part of figure 3 is changed to
$$\{\begin{array}{cc}\overline{p}_i=t+1:\hfill & \varphi _{0,0;s_i+\overline{p}_i1}\mathrm{}\varphi _{0,0;s_i+2}\varphi _{0,0;s_i+1}\hfill \\ \overline{p}_i>t+1:\hfill & \varphi _{0,0;s_i+\overline{p}_i1}\mathrm{}\varphi _{0,0;s_i+\overline{p}_it+1}\varphi _{0,0;s_i+\overline{p}_it}\hfill \\ & \varphi _{2,0;s_i+\overline{p}_it1}\mathrm{}\varphi _{2,0;s_i+2}\varphi _{2,0;s_i+1}\hfill \end{array}$$
The above changes transform figures 3, 4 and 5 to figures 3’, 4’ and 5’, respectively. Then the factors in the $`i`$-th line in (63) are related to the numbers given to the boxes in the $`i`$-th row and the left neighbouring column by the following three steps.
1. For $`1iR`$ and $`\overline{p}_i2`$, the factors
(156)
$$\begin{array}{c}\varphi _{j_{\overline{p}_1+\mathrm{}+\overline{p}_i1},t_{\overline{p}_1+\mathrm{}+\overline{p}_i1};s_i+\overline{p}_i2}\mathrm{}\varphi _{j_{\overline{p}_1+\mathrm{}+\overline{p}_{i1}+2},t_{\overline{p}_1+\mathrm{}+\overline{p}_{i1}+2};s_i+1}\hfill \\ \hfill \varphi _{j_{\overline{p}_1+\mathrm{}+\overline{p}_{i1}+1},t_{\overline{p}_1+\mathrm{}+\overline{p}_{i1}+1};s_i}\end{array}$$
are determined according to the numbers attributed to the $`\overline{p}_i`$ rightmost boxes of the $`i`$-th row. Figure 3’ specifies these factors for all cases admitted by the rules (47) and (48) for the tableaux of vertical strips.
2. To obtain the factor $`\varphi _{j_{\overline{p}_1+\mathrm{}+\overline{p}_i},t_{\overline{p}_1+\mathrm{}+\overline{p}_i};s_i+\overline{p}_i}`$ with $`i<R`$ and $`\overline{p}_i1`$, the $`i`$-th row and the $`(i+1)`$-th finite column with length $`p_{i+1}+1`$ is taken into account. Figure 4’ gives all possible cases together with the associated factors.
3. Depending on the numbers given to the boxes of the left infinite column, the leftmost factor $`\varphi _{j_{\overline{N}},t_{\overline{N}};s_R+\overline{p}_R1}`$ is determined. All cases are found in figure 5’.
For a semi-standard super tableau of the border strip $`(\mathrm{},\overline{p}_1)`$ with $`\overline{p}_1>1`$, the normal form (61) is determined by the second and third step. In the case $`(\mathrm{};1)`$, only the last step is needed. The set of all normal forms (155) may be called $`B(p_1,\mathrm{},p_{R1};\overline{p}_1,\mathrm{},\overline{p}_R;r_0)`$ if $`R>1`$ and $`B(\mathrm{};\overline{p}_1)`$ for $`R=1`$. Steps 156-3 describe the one-to-one correspondence claimed in Result 2.
## Appendix B The R-matrix $`\stackrel{ˇ}{R}_{VV}(z)`$
The intertwiner $`\stackrel{ˇ}{R}:V_zV_wV_wV_z`$ satisfies the conditions
(157)
$$\stackrel{ˇ}{R}\mathrm{\Delta }(a)=\mathrm{\Delta }(a)\stackrel{ˇ}{R}aU_q\left(\widehat{gl}(2|2)\right).$$
This set of linear equations has a solution
(158)
$$\stackrel{ˇ}{R}_{VV}\left(\frac{z}{w}\right)=P^{gr}R_{VV}\left(\frac{z}{w}\right)$$
where $`P^{gr}`$ is the graded transposition defined in section 6 and $`R_{VV}\left(\frac{z}{w}\right)`$ denotes the R-matrix in the first line of (80). In the following, the nonvanishing matrix elements $`\stackrel{ˇ}{R}_{j_1,t_1;j_2,t_2}^{j_3,t_3;j_4,t_4}(z)`$ of $`\stackrel{ˇ}{R}_{\stackrel{̊}{V}\stackrel{̊}{V}}(z)`$ with $`t0`$ and the basis of $`\stackrel{̊}{V}`$ specified by (82), (83) are collected. They are related by
(159)
$$\begin{array}{c}\stackrel{ˇ}{R}_{j_1,t_1;j_2,t_2}^{j_3,t_3;j_4,t_4}(z)=q^{\frac{1}{8}(1(1)^{j_3})(1(1)^{j_4})\frac{1}{8}(1(1)^{j_1})(1(1)^{j_2})}\hfill \\ \hfill q^{(t_1+\frac{1}{2}\delta _{j_1,1}+\frac{1}{2}\delta _{j_1,3})^2(t_2+\frac{1}{2}\delta _{j_2,1}+\frac{1}{2}\delta _{j_2,3})^2+(t_3+\frac{1}{2}\delta _{j_3,1}+\frac{1}{2}\delta _{j_3,3})^2+(t_4+\frac{1}{2}\delta _{j_4,1}+\frac{1}{2}\delta _{j_4,3})^2}\\ \hfill \left(\frac{[t_1]![t_2]!}{[t_3]![t_4]!}\right)^2\frac{[t_1+1]^{22\delta _{j_1,2}+\delta _{j_1,1}\delta _{j_1,3}}[t_2+1]^{22\delta _{j_2,2}+\delta _{j_2,1}\delta _{j_2,3}}}{[t_3+1]^{22\delta _{j_3,2}+\delta _{j_3,1}\delta _{j_3,3}}[t_4+1]^{22\delta _{j_4,2}+\delta _{j_4,1}\delta _{j_4,3}}}\stackrel{ˇ}{R}_{j_3,t_3;j_4,t_4}^{j_1,t_1;j_2,t_2}(z)\end{array}$$
and
(160)
$$\begin{array}{c}\stackrel{ˇ}{R}_{j_1,t_1;j_2,t_2}^{j_3,t_3;j_4,t_4}(z)=\hfill \\ \hfill z^{t_4t_1+\delta _{j_4,1}\delta _{j_1,1}}\frac{[t_1+1]^{\delta _{j_1,0}\delta _{j_1,2}}[t_2+1]^{\delta _{j_2,0}\delta _{j_2,2}}}{[t_3+1]^{\delta _{j_3,0}\delta _{j_3,2}}[t_4+1]^{\delta _{j_4,0}\delta _{j_4,2}}}\stackrel{ˇ}{R}_{\sigma (j_2),t_2;\sigma (j_1),t_1}^{\sigma (j_4),t_4;\sigma (j_3),t_3}(z)\end{array}$$
with $`\sigma (0)=2`$, $`\sigma (2)=0`$, $`\sigma (1)=1`$, $`\sigma (3)=3`$ and
(161)
$$\begin{array}{c}\stackrel{ˇ}{R}_{j_1,t_1;j_2,t_2}^{j_3,t_3;j_4,t_4}(z)=\hfill \\ \hfill z^{t_4t_1\delta _{j_4,2}+\delta _{j_1,2}}\frac{[t_1+1]^{\delta _{j_1,1}\delta _{j_1,3}}[t_2+1]^{\delta _{j_2,1}\delta _{j_2,3}}}{[t_3+1]^{\delta _{j_3,1}\delta _{j_3,3}}[t_4+1]^{\delta _{j_4,1}\delta _{j_4,3}}}\stackrel{ˇ}{R}_{\tau (j_2),t_2;\tau (j_1),t_1}^{\tau (j_4),t_4;\tau (j_3),t_3}(z)\end{array}$$
with $`\tau (0)=0`$, $`\tau (2)=2`$, $`\tau (1)=3`$ and $`\tau (3)=1`$. Here the notation $`[t]!=[t][t1]\mathrm{}[1]`$ is used.
The matrix elements $`\stackrel{ˇ}{R}_{j,t_1;j,t_2}^{j,t_3;j,t_4}(z)`$ are nonvanishing for $`t_1+t_2=t_3+t_4`$. The solution of the intertwining condition (157) is unique up to a scalar. Normalising $`\stackrel{ˇ}{R}_{\stackrel{̊}{V}\stackrel{̊}{V}}(z)`$ by
(162)
$$\stackrel{ˇ}{R}_{2,0;\mathrm{\hspace{0.17em}2},0}^{2,0;\mathrm{\hspace{0.17em}2},0}(z)=1,$$
explicit solution of the intertwining condition (157) gives
(163)
$$\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}2},t_2}^{2,0;\mathrm{\hspace{0.17em}2},t_1+t_2}(z)=z^{t_2}q^{t_1(2t_2+1)}(q^2;q^2)_{t_2}\frac{z1}{q^21}\underset{r=0}{\overset{t_2}{}}\frac{q^21}{q^{2(t_1+r)}z1}$$
and
(164)
$$\begin{array}{cc}& \stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}2},t_2}^{2,s;\mathrm{\hspace{0.17em}2},t_1+t_2s}(z)=\frac{1}{(q^2;q^2)_s}\frac{q^{2t_1}z1}{zq^{2(t_1s)}}q^{(t_1s)(2t_22s+1)t_1(2t_2+1)}\overline{R}_{2,t_1;\mathrm{\hspace{0.17em}2},t_2}^{2,0;\mathrm{\hspace{0.17em}2},t_1+t_2}(z)\hfill \\ & \underset{r=0}{\overset{s}{}}z^{rs}q^{(sr)(sr+1)}\frac{1}{(q^2;q^2)_r(q^2;q^2)_{sr}}\frac{(q^2;q^2)_{t_2+r}}{(q^2;q^2)_{t_2s+r}}\left(\frac{(q^2;q^2)_{t_1}}{(q^2;q^2)_{t_1r}}\right)^2\hfill \\ & \underset{r_1=0}{\overset{sr}{}}(zq^{2(t_1rr_1)})\underset{r_2=0}{\overset{r}{}}\frac{1}{q^{2(t_1r_2)}z1}\hfill \end{array}$$
for $`t_1s>0`$ and $`t_2s0`$. In (163), (164) and the remainder, the notation $`(a;q^2)_n=\left(1q^{2(n1)}a\right)\mathrm{}(1q^2a)(1a)`$ is used. The remaining elements $`\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}2},t_2}^{2,s;\mathrm{\hspace{0.17em}2},t_1+t_2s}(z)`$ are obtained from the relations
(165)
$$\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}2},t_2}^{2,t_3;\mathrm{\hspace{0.17em}2},t_4}(z)=z^{t_4t_1}\stackrel{ˇ}{R}_{2,t_2;\mathrm{\hspace{0.17em}2},t_1}^{2,t_4;\mathrm{\hspace{0.17em}2},t_3}(z)$$
and
(166)
$$\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}2},t_2}^{2,t_3;\mathrm{\hspace{0.17em}2},t_4}(z)=q^{2(t_1^2+t_2^2t_3^2t_4^2)}\left(\frac{(q^2;q^2)_{t_1}(q^2;q^2)_{t_2}}{(q^2;q^2)_{t_3}(q^2;q^2)_{t_4}}\right)^2\stackrel{ˇ}{R}_{2,t_3;\mathrm{\hspace{0.17em}2},t_4}^{2,t_1;\mathrm{\hspace{0.17em}2},t_2}(z)$$
which are special cases of (159)-(161). All matrix elements $`\stackrel{ˇ}{R}_{j,t_1;j,t_2}^{j,t_3;j,t_4}(z)`$ can be expressed in terms of $`\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}2},t_2}^{2,t_3;\mathrm{\hspace{0.17em}2},t_4}(z)`$:
(167)
$$\begin{array}{c}\stackrel{ˇ}{R}_{3,t_1;\mathrm{\hspace{0.17em}3},t_2}^{3,s;\mathrm{\hspace{0.17em}3},t_1+t_2s}(z)=\frac{1}{(1z)\left(q^{2(t_2+1)}1\right)}\hfill \\ \hfill \{q^{st_1}(q^{2(st_1)}1)(q^{2(t_1+t_2s+1)}1)\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}2},t_2+1}^{2,s;\mathrm{\hspace{0.17em}2},t_1+t_2s+1}(z)\\ \hfill +q^{t_1+2t_23s1}(zq^{2(st_1+1)})(q^{2(s+1)}1)\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}2},t_2+1}^{2,s+1;\mathrm{\hspace{0.17em}2},t_1+t_2s}(z)\}\end{array}$$
(168)
$$\begin{array}{cc}\hfill \stackrel{ˇ}{R}_{0,t_1;\mathrm{\hspace{0.17em}0},t_2}^{0,t_3;\mathrm{\hspace{0.17em}0},t_4}(z)& =q^{t_3+t_4t_1t_2}\frac{\left(q^{2(t_1+1)}1\right)\left(q^{2(t_2+1)}1\right)}{\left(q^{2(t_3+1)}1\right)\left(q^{2(t_4+1)}1\right)}\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}2},t_2}^{2,t_3;\mathrm{\hspace{0.17em}2},t_4}(z)\hfill \\ \hfill \stackrel{ˇ}{R}_{1,t_1;\mathrm{\hspace{0.17em}1},t_2}^{1,t_3;\mathrm{\hspace{0.17em}1},t_4}(z)& =q^{t_3+t_4t_1t_2}\frac{\left(q^{2(t_1+1)}1\right)\left(q^{2(t_2+1)}1\right)}{\left(q^{2(t_3+1)}1\right)\left(q^{2(t_4+1)}1\right)}\stackrel{ˇ}{R}_{3,t_1;\mathrm{\hspace{0.17em}3},t_2}^{3,t_3;\mathrm{\hspace{0.17em}3},t_4}(z)\hfill \end{array}$$
The R-matrix elements $`\stackrel{ˇ}{R}_{j_1,t_1;j_2,t_2}^{j_3,t_3;j_4,t_4}(z)`$ with $`j_1=j_3j_2=j_4`$ or $`j_1=j_4j_2=j_3`$ are given by
(169)
$$\begin{array}{c}\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}3},t_2}^{2,s;\mathrm{\hspace{0.17em}3},t_1+t_2s}(z)=\frac{q^{t_1+1}}{q^{2(t_1+t_2+2)}1}\hfill \\ \hfill \left\{q^{t_1+1}\left(q^{2(t_2+1)}1\right)\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}2},t_2}^{2,s;\mathrm{\hspace{0.17em}2},t_1+t_2s}(z)+q^{s1}\left(q^{2(s+1)}1\right)\stackrel{ˇ}{R}_{3,t_1;\mathrm{\hspace{0.17em}3},t_2}^{3,s;\mathrm{\hspace{0.17em}3},t_1+t_2s}(z)\right\}\end{array}$$
(170)
$$\begin{array}{c}\stackrel{ˇ}{R}_{3,t_1;\mathrm{\hspace{0.17em}2},t_2}^{3,s;\mathrm{\hspace{0.17em}2},t_1+t_2s}(z)=\frac{q^{s+1}}{q^{2(t_1+t_2+2)}1}\hfill \\ \hfill \left\{q^{t_11}\left(q^{2(t_1+1)}1\right)\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}2},t_2}^{2,s_1;\mathrm{\hspace{0.17em}2},s_2}(z)+q^{s+1}\left(q^{2(t_1+t_2s+1)}1\right)\stackrel{ˇ}{R}_{3,t_1;\mathrm{\hspace{0.17em}3},t_2}^{3,s;\mathrm{\hspace{0.17em}3},t_1+t_2s}(z)\right\}\end{array}$$
(171)
$$\begin{array}{c}\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}3},t_2}^{3,s;\mathrm{\hspace{0.17em}2},t_1+t_2s}(z)=\frac{q^{st_2+t_1+1}}{q^{2(t_1+t_2+2)}1}\hfill \\ \hfill \left\{\left(q^{2(t_2+1)}1\right)\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}2},t_2}^{2,s;\mathrm{\hspace{0.17em}2},t_1+t_2s}(z)q^{st_1}\left(q^{2(t_1+t_2s+1)}1\right)\stackrel{ˇ}{R}_{3,t_1;\mathrm{\hspace{0.17em}3},t_2}^{3,s;\mathrm{\hspace{0.17em}3},t_1+t_2s}(z)\right\}\end{array}$$
(172)
$$\stackrel{ˇ}{R}_{3,t_1;\mathrm{\hspace{0.17em}2},t_2}^{2,s;\mathrm{\hspace{0.17em}3},t_1+t_2s}(z)=\left(\frac{z}{q}\right)^{t_2s}\stackrel{ˇ}{R}_{2,t_2;\mathrm{\hspace{0.17em}3},t_1}^{3,t_1+t_2s;\mathrm{\hspace{0.17em}2},s}(z)$$
(173)
$$\begin{array}{c}\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}3},t_2}^{3,s;\mathrm{\hspace{0.17em}2},t_1+t_2s}(z)=\hfill \\ \hfill \frac{q^{2(s+1)}1}{q^{2(t_2+1)}1}\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}1},t_2}^{1,s;\mathrm{\hspace{0.17em}2},t_1+t_2s}(z)=q^{2(st_2)}\frac{q^{2(t_1+t_2s+1)}1}{q^{2(t_1+1)}1}\stackrel{ˇ}{R}_{0,t_1;\mathrm{\hspace{0.17em}3},t_2}^{3,s;\mathrm{\hspace{0.17em}0},t_1+t_2s}(z)=\\ \hfill \frac{\left(q^{2(s+1)}1\right)\left(q^{2(t_1+t_2s+1)}1\right)}{\left(q^{2(t_1+1)}1\right)\left(q^{2(t_2+1)}1\right)}\stackrel{ˇ}{R}_{0,t_1;\mathrm{\hspace{0.17em}1},t_2}^{1,s;\mathrm{\hspace{0.17em}0},t_1+t_2s}(z)\end{array}$$
(174)
$$\begin{array}{c}\stackrel{ˇ}{R}_{3,t_1;\mathrm{\hspace{0.17em}2},t_2}^{2,s;\mathrm{\hspace{0.17em}3},t_1+t_2s}(z)=\hfill \\ \hfill \frac{q^{2(s+1)}1}{q^{2(t_2+1)}1}\stackrel{ˇ}{R}_{3,t_1;\mathrm{\hspace{0.17em}0},t_2}^{0,s;\mathrm{\hspace{0.17em}3},t_1+t_2s}(z)=q^{2(st_2)}\frac{q^{2(t_1+t_2s+1)}1}{q^{2(t_1+1)}1}\stackrel{ˇ}{R}_{1,t_1;\mathrm{\hspace{0.17em}2},t_2}^{2,s;\mathrm{\hspace{0.17em}1},t_1+t_2s}(z)=\\ \hfill \frac{\left(q^{2(s+1)}1\right)\left(q^{2(t_1+t_2s+1)}1\right)}{\left(q^{2(t_1+1)}1\right)\left(q^{2(t_2+1)}1\right)}\stackrel{ˇ}{R}_{1,t_1;\mathrm{\hspace{0.17em}0},t_2}^{0,s;\mathrm{\hspace{0.17em}1},t_1+t_2s}(z)\end{array}$$
(175)
$$\begin{array}{c}\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}3},t_2}^{2,s;\mathrm{\hspace{0.17em}3},t_1+t_2s}(z)=\frac{\left(q^{2(s+1)}1\right)\left(q^{2(t_1+t_2s+1)}1\right)}{q^{2(t_1+1)}1\left)\right(q^{2(t_2+1)}1)}\stackrel{ˇ}{R}_{0,t_1;\mathrm{\hspace{0.17em}1},t_2}^{0,s;\mathrm{\hspace{0.17em}1},t_1+t_2s}(z)=\hfill \\ \hfill q^{st_1}z^{t_2s+1}\frac{q^{2(t_1+t_2s+1)}1}{q^{2(t_2+1)}1}\stackrel{ˇ}{R}_{1,t_2;\mathrm{\hspace{0.17em}2},t_1}^{1,t_1+t_2s;\mathrm{\hspace{0.17em}2},s}(z)=\\ \hfill q^{t_1s}z^{t_2s}\frac{q^{2(s+1)}1}{q^{2(t_1+1)}1}\stackrel{ˇ}{R}_{3,t_2;\mathrm{\hspace{0.17em}0},t_1}^{3,t_1+t_2s;\mathrm{\hspace{0.17em}0},s}(z)\end{array}$$
(176)
$$\begin{array}{c}\stackrel{ˇ}{R}_{3,t_1;\mathrm{\hspace{0.17em}2},t_2}^{3,s;\mathrm{\hspace{0.17em}2},t_1+t_2s}(z)=\frac{\left(q^{2(s+1)}1\right)\left(q^{2(t_1+t_2s+1)}1\right)}{q^{2(t_1+1)}1\left)\right(q^{2(t_2+1)}1)}\stackrel{ˇ}{R}_{1,t_1;\mathrm{\hspace{0.17em}0},t_2}^{1,s;\mathrm{\hspace{0.17em}0},t_1+t_2s}(z)=\hfill \\ \hfill q^{t_1s}z^{t_2s1}\frac{q^{2(s+1)}1}{q^{2(t_1+1)}1}\stackrel{ˇ}{R}_{2,t_2;\mathrm{\hspace{0.17em}1},t_1}^{2,t_1+t_2s;\mathrm{\hspace{0.17em}1},s}(z)=\\ \hfill q^{st_1}z^{t_2s}\frac{q^{2(t_1+t_2s+1)}1}{q^{2(t_2+1)}1}\stackrel{ˇ}{R}_{0,t_2;\mathrm{\hspace{0.17em}3},t_1}^{0,t_1+t_2s;\mathrm{\hspace{0.17em}3},s}(z)\end{array}$$
(177)
$$\begin{array}{c}\stackrel{ˇ}{R}_{2,t_2;\mathrm{\hspace{0.17em}0},t_1}^{0,t_1+t_2s;\mathrm{\hspace{0.17em}2},s}(z)=z^{st_2}\stackrel{ˇ}{R}_{0,t_1;\mathrm{\hspace{0.17em}2},t_2}^{2,s;\mathrm{\hspace{0.17em}0},t_1+t_2s}(z)=\hfill \\ \hfill \frac{q^{t_2}z^{st_2}}{q^{2(t_1+t_2+1)}1}\{q^{t_1s+1}\frac{\left(q^{2(t_1+1)}1\right)\left(q^{2t_2}1\right)}{q^{2(t_1+t_2s+1)}1}\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}3},t_21}^{3,s1;\mathrm{\hspace{0.17em}2},t_1+t_2s}(z)\\ \hfill +(q^{2(t_1+1)}1)\stackrel{ˇ}{R}_{3,t_1;\mathrm{\hspace{0.17em}2},t_2}^{2,s;\mathrm{\hspace{0.17em}3},t_1+t_2s}(z)\}\end{array}$$
(178)
$$\begin{array}{c}\stackrel{ˇ}{R}_{3,t_2;\mathrm{\hspace{0.17em}1},t_1}^{1,t_1+t_2s;\mathrm{\hspace{0.17em}3},s}(z)=z^{st_2}\stackrel{ˇ}{R}_{1,t_1;\mathrm{\hspace{0.17em}3},t_2}^{3,s;\mathrm{\hspace{0.17em}1},t_1+t_2s}(z)=\hfill \\ \hfill \frac{q^{t_2+1}z^{st_2}}{q^{2(t_1+t_2+2)}1}\{q^{t_1s}\frac{\left(q^{2(t_1+1)}1\right)\left(q^{2(t_2+1)}1\right)}{q^{2(t_1+t_2s+1)}1}\stackrel{ˇ}{R}_{3,t_1;\mathrm{\hspace{0.17em}2},t_2}^{2,s;\mathrm{\hspace{0.17em}3},t_1+t_2s}(z)\\ \hfill +(q^{2(t_1+1)}1)\stackrel{ˇ}{R}_{2,t_1+1;\mathrm{\hspace{0.17em}3},t_2}^{3,s;\mathrm{\hspace{0.17em}2},t_1+t_2s+1}(z)\}\end{array}$$
(179)
$$\begin{array}{c}\stackrel{ˇ}{R}_{2,t_2;\mathrm{\hspace{0.17em}0},t_1}^{2,t_1+t_2s;\mathrm{\hspace{0.17em}0},s}(z)=z^{st_2+1}\stackrel{ˇ}{R}_{0,t_1;\mathrm{\hspace{0.17em}2},t_2}^{0,s;\mathrm{\hspace{0.17em}2},t_1+t_2s}(z)\hfill \\ \hfill =z^{st_2+1}\frac{q^{2(t_1+1)}1}{q^{2(t_1+t_2+1)}1}\{q^{t_1+2t_2s}\stackrel{ˇ}{R}_{3,t_1;\mathrm{\hspace{0.17em}2},t_2}^{3,s;\mathrm{\hspace{0.17em}2},t_1+t_2s}\\ \hfill +\frac{q^{2t_2}1}{q^{2(s+1)}1}\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}3},t_21}^{2,s;\mathrm{\hspace{0.17em}3},t_1+t_1s1}(z)\}\end{array}$$
(180)
$$\stackrel{ˇ}{R}_{3,t_2;\mathrm{\hspace{0.17em}1},t_1}^{3,t_1+t_2s;\mathrm{\hspace{0.17em}1},s}(z)=z^{st_2}q^{2(st_1)}\frac{\left(q^{2(t_1+1)}1\right)\left(q^{2(t_1+t_2s+1)}1\right)}{\left(q^{2(s+1)}1\right)\left(q^{2(t_2+1)}1\right)}\stackrel{ˇ}{R}_{3,t_1;\mathrm{\hspace{0.17em}1},t_2}^{3,s;\mathrm{\hspace{0.17em}1},t_1+t_2s}(z)$$
(181)
$$\begin{array}{c}\stackrel{ˇ}{R}_{3,t_1;\mathrm{\hspace{0.17em}1},t_2}^{3,s;\mathrm{\hspace{0.17em}1},t_1+t_2s}(z)=\frac{q^{2(t_2+1)}1}{q^{2(t_1+t_2+2)}1}\hfill \\ \hfill \left\{\stackrel{ˇ}{R}_{3,t_1;\mathrm{\hspace{0.17em}2},t_2+1}^{3,s;\mathrm{\hspace{0.17em}2},t_1+t_2s+1}(z)+q^{t_1+2t_2s+2}\frac{q^{2(t_1+1)}1}{q^{2(t_1+t_2s+1)}1}\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}3},t_2}^{2,s;\mathrm{\hspace{0.17em}3},t_1+t_2s}(z)\right\}\end{array}$$
(182)
$$\stackrel{ˇ}{R}_{1,t_2;\mathrm{\hspace{0.17em}3},t_1}^{1,t_1+t_2s;\mathrm{\hspace{0.17em}3},s}(z)=z^{st_2}q^{2(t_1s)}\frac{\left(q^{2(s+1)}1\right)\left(q^{2(t_2+1)}1\right)}{\left(q^{2(t_1+1)}1\right)\left(q^{2(t_1+t_2s+1)}1\right)}\stackrel{ˇ}{R}_{1,t_1;\mathrm{\hspace{0.17em}3},t_2}^{1,s;\mathrm{\hspace{0.17em}3},t_1+t_2s}(z)$$
(183)
$$\begin{array}{c}\stackrel{ˇ}{R}_{1,t_1;\mathrm{\hspace{0.17em}3},t_2}^{1,s;\mathrm{\hspace{0.17em}3},t_1+t_2s}(z)=\hfill \\ \hfill \frac{1}{z}\frac{q^{t_1+2t_2+s+1}}{q^{2(t_1+t_2+2)}1}\frac{q^{2(t_1+1)}1}{q^{2(s+1)}1}\{(q^{2(t_1+1)}1)\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}3},t_2}^{2,s;\mathrm{\hspace{0.17em}3},t_1+t_2s}(z)\\ \hfill +q^{t_12t_2+s1}(q^{2(t_1+t_2s+1)}1)\stackrel{ˇ}{R}_{3,t_1;\mathrm{\hspace{0.17em}2},t_2+1}^{3,s;\mathrm{\hspace{0.17em}2},t_1+t_2s+1}(z)\}\end{array}$$
The above expressions are valid for $`0st_1+t_2`$. In (179), the second contribution on the rhs is dropped for $`s=t_1+t_2`$. The remaining nonvanishing R-matrix elements are
(184)
$$\begin{array}{c}\stackrel{ˇ}{R}_{3,t_2;\mathrm{\hspace{0.17em}1},t_1}^{0,t_1+t_2s+1;\mathrm{\hspace{0.17em}2},s}(z)=\hfill \\ \hfill q^{2(t_2s)+1}\frac{\left(q^{2(t_1+1)}1\right)\left(q^{2(s+1)}1\right)}{\left(q^{2(t_2+1)}1\right)\left(q^{2(t_1+t_2s+2)}1\right)}\stackrel{ˇ}{R}_{1,t_2;\mathrm{\hspace{0.17em}3},t_1}^{2,t_1+t_2s+1;\mathrm{\hspace{0.17em}0},s}(z)=\\ \hfill z^{st_21}q^{t_2t_1}\frac{q^{2(t_1+1)}1}{q^{2(t_2+1)}1}\stackrel{ˇ}{R}_{3,t_1;\mathrm{\hspace{0.17em}1},t_2}^{2,s;\mathrm{\hspace{0.17em}0},t_1+t_2s+1}(z)=\\ \hfill z^{st_2}q^{t_1+t_22s+1}\frac{q^{2(s+1)}1}{q^{2(t_1+t_2s+2)}1}\stackrel{ˇ}{R}_{1,t_1;\mathrm{\hspace{0.17em}3},t_2}^{0,s;\mathrm{\hspace{0.17em}2},t_1+t_2s+1}(z)\end{array}$$
(185)
$$\begin{array}{c}\stackrel{ˇ}{R}_{1,t_1;\mathrm{\hspace{0.17em}3},t_2}^{0,s;\mathrm{\hspace{0.17em}2},t_1+t_2s+1}(z)=\frac{q^{2(t_1+1)}1}{q^{2(t_1+t_2+2)}1}\hfill \\ \hfill \left\{q^{t_1+2t_2s+2}\stackrel{ˇ}{R}_{2,t_1+1;\mathrm{\hspace{0.17em}3},t_2}^{3,s;\mathrm{\hspace{0.17em}2},t_1+t_2s+1}(z)\frac{q^{2(t_2+1)}1}{q^{2(s+1)}1}\stackrel{ˇ}{R}_{3,t_1;\mathrm{\hspace{0.17em}2},t_2}^{2,s;\mathrm{\hspace{0.17em}3},t_1+t_2s}(z)\right\}\end{array}$$
with $`0st_1+t_2+1`$ and
(186)
$$\begin{array}{c}\stackrel{ˇ}{R}_{2,t_2;\mathrm{\hspace{0.17em}0},t_1}^{1,t_1+t_2s;\mathrm{\hspace{0.17em}3},s1}(z)=\hfill \\ \hfill q^{2(t_2s)+1}\frac{\left(q^{2(t_1+1)}1\right)\left(q^{2s}1\right)}{(q^{2(t_2+1)}1)(q^{2(t_1+t_2s+1)}1}\stackrel{ˇ}{R}_{0,t_2;\mathrm{\hspace{0.17em}2},t_1}^{3,t_1+t_2s;\mathrm{\hspace{0.17em}1},s1}(z)=\\ \hfill z^{st_2}q^{t_1+t_22s+1}\frac{q^{2s}1}{q^{2(t_1+t_2s+1)}1}\stackrel{ˇ}{R}_{0,t_1;\mathrm{\hspace{0.17em}2},t_2}^{1,s1;\mathrm{\hspace{0.17em}3},t_1+t_2s}(z)=\\ \hfill z^{st_21}q^{t_2t_1}\frac{q^{2(t_1+1)}1}{q^{2(t_2+1)}1}\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}0},t_2}^{3,s1;\mathrm{\hspace{0.17em}1},t_1+t_2s}(z)\end{array}$$
(187)
$$\begin{array}{c}\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}0},t_2}^{3,s1;\mathrm{\hspace{0.17em}1},t_1+t_2s}(z)=\frac{q^{2(t_2+1)}1}{q^{2(t_1+t_2+1)}1}\hfill \\ \hfill \left\{q^{t_1+2t_2s+2}\frac{q^{2t_1}1}{q^{2(t_1+t_2s+1)}1}\stackrel{ˇ}{R}_{3,t_11;\mathrm{\hspace{0.17em}2},t_2}^{2,s1;\mathrm{\hspace{0.17em}3},t_1+t_2s}(z)\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}3},t_2}^{3,s1;\mathrm{\hspace{0.17em}2},t_1+t_2s+1}(z)\right\}\end{array}$$
with $`t_1+t_2>0`$ and $`1st_1+t_2`$. In the limit $`q0`$, well-defined expressions are found for the R-matrix elements with respect to a different basis $`\{\stackrel{~}{v}_{j,t}\}_{0j3,t_0}`$ given by
(188)
$$\stackrel{~}{v}_{j,t}=q^{t(t+\alpha _j)}v_{j,t}$$
with $`\alpha _j=\frac{1}{2}+\delta _{j,1}\delta _{j,2}`$. The matrix elements $`\stackrel{~}{R}_{j_1,t_1;j_2,t_2}^{j_3,t_3;j_4,t_4}(z)`$ of the corresponding R-matrix are related to the above matrix elements by
(189)
$$\stackrel{~}{R}_{j_1,t_1;j_2,t_2}^{j_3,t_3;j_4,t_4}(z)=q^{t_1(t_1+\alpha _{j_1})+t_2(t_2+\alpha _{j_2})t_3(t_3+\alpha _{j_3})t_4(t_4+\alpha _{j_4})}\overline{R}_{j_1,t_1;j_2,t_2}^{j_3,t_3;j_4,t_4}(z)$$
Straightforward analysis of (163)-(187) with properties (159)-(161) yields
(190)
$$\underset{q0}{lim}\stackrel{~}{R}_{j_1,t_1;j_2,t_2}^{j_3,t_3;j_4,t_4}(z)=\delta _{j_1,j_4}\delta _{j_2,j_3}\delta _{t_1,t_4}\delta _{t_2,t_3}z^{t_2+\theta _{j_2,j_1}}$$
## Appendix C Evaluation of $`r(y,z)`$
The function $`r(y,z)`$ in (140) and (141) can be obtained by collecting all contributions to the term $`v_{0,0}^{}v_{2,r}`$ on the rhs of (141):
(191)
$$\begin{array}{c}r(y,z)=\frac{1}{\stackrel{~}{P}_0(y;q^{2r1},q^1,q^1|q^2)}\hfill \\ \hfill \underset{t=0}{\overset{\mathrm{}}{}}q^{t(2r+1)}\frac{(q^2;q^2)_{t+1}(q^2;q^2)_r}{(q^2;q^2)_{t+r}(q^2;q^2)_1}\stackrel{~}{P}_t(y;q^{2r+1},q,q|q^2)\stackrel{ˇ}{R}_{2,t+r;\mathrm{\hspace{0.17em}0},t^{}}^{0,0^{};\mathrm{\hspace{0.17em}2},r}(z)\end{array}$$
A choice of $`\stackrel{~}{P}_0(y;q^{2r1},q^1,q^1|q^2)`$ sets the normalisation of the polynomials $`\stackrel{~}{P}_t(y;q^{2r1},q^1,q^1|q^2)`$. Its value may depend on $`r`$. Comparison of the intertwining conditions yields a relation between the matrix elements $`\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}0},t_2^{}}^{0,t_3^{};\mathrm{\hspace{0.17em}2},t_1t_2+t_3}(z)`$ or $`\stackrel{ˇ}{R}_{0,t_2^{};\mathrm{\hspace{0.17em}2},t_1}^{2,t_1t_2+t_3;\mathrm{\hspace{0.17em}0},t_3^{}}(z)`$ and matrix elements of $`\stackrel{ˇ}{R}_{VV}(q^2z^1)`$:
(192)
$$\begin{array}{cc}\hfill \stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}0},t_2^{}}^{0,t_3^{},\mathrm{\hspace{0.17em}2},t_1t_2+t_3}(z)& =q^{t_3t_2}\frac{[t_3+1]}{[t_2+1]}\stackrel{ˇ}{R}_{2,t_3;\mathrm{\hspace{0.17em}2},t_1}^{2,t_1t_2+t_3;\mathrm{\hspace{0.17em}2},t_2}(q^2z^1)\hfill \\ \hfill \stackrel{ˇ}{R}_{0,t_2^{};\mathrm{\hspace{0.17em}2},t_1}^{2,t_1t_2+t_3;\mathrm{\hspace{0.17em}0},t_3^{}}(z)& =q^{t_2t_3}\frac{[t_3+1]}{[t_2+1]}\stackrel{ˇ}{R}_{2,t_1;\mathrm{\hspace{0.17em}2},t_3}^{2,t_2;\mathrm{\hspace{0.17em}2},t_1t_2+t_3}(q^2z^1)\hfill \end{array}$$
Here an overall normalisation is fixed by equating the R-matrix elements on the lhs with those on rhs for $`t_1=t_2=t_3=0`$. Using (192) with (159), (163), the sum in (191) can be written
(193)
$$q^r\frac{1q^2z^1}{1q^{2(r+1)}z^1}\underset{t=0}{\overset{\mathrm{}}{}}\stackrel{~}{P}_t(y;q^{2r+1},q,q|q^2)\frac{(q^{2(r+1)};q^2)_t}{(q^{2(r+2)}z^1;q^2)_t}\left(qz^1\right)^t$$
The infinite sum in (193) is related to a generating function for Askey-Wilson polynomials derived in . In base $`q^2`$, eq. 2.10 in reads
(194)
$$\begin{array}{c}\underset{t=0}{\overset{\mathrm{}}{}}\frac{(abcd;q^2)_{2t}p_t(\mathrm{cos}\theta ;a,b,c,d|q^2)}{(q^2,ab,ac,bc,abcdq^{2(t1)};q^2)_t}\frac{(f/g,abc/f;q^2)_t}{(abcdg/f,fd;q^2)_t}g^t=\hfill \\ \hfill \frac{(abcd,dg,abcdge^{i\theta }/f,fe^{i\theta };q^2)_{\mathrm{}}}{(abcdg/f,df,abce^{i\theta },ge^{i\theta };q^2)_{\mathrm{}}}{}_{8}{}^{}W_{7}^{}(abce^{i\theta }/q^2;ae^{i\theta },be^{i\theta },ce^{i\theta },f/g,abc/f;q^2,ge^{i\theta })\end{array}$$
for $`|g|<1`$. For the definitions of the very-well-poised basic hypergeometric series $`{}_{8}{}^{}W_{7}^{}(a_1;a_4,a_5,a_6,a_7,a_8;q^2,\alpha )`$, the reader is referred to . The orthonormal continuous dual $`q^2`$-Hahn polynomials $`\stackrel{~}{P}_t(\mathrm{cos}\theta ;a,c,d|q^2)`$ are related to the Askey-Wilson polynomials by
(195)
$$\stackrel{~}{P}_t(\mathrm{cos}\theta ;a,c,d|q^2)=\frac{p_t(\mathrm{cos}\theta ;a,0,c,d|q^2)}{(q^2,ac;q^2)_t}$$
For $`b=0`$, equation (194) simplifies:
(196)
$$\begin{array}{c}\underset{t=0}{\overset{\mathrm{}}{}}\stackrel{~}{P}_t(\mathrm{cos}\theta ;a,c,d|q^2)\frac{(f/g;q^2)_t}{(df;q^2)_t}g^t=\hfill \\ \hfill \frac{(dg,fe^{i\theta };q^2)_{\mathrm{}}}{(df,ge^{i\theta };q^2)_{\mathrm{}}}{}_{3}{}^{}\varphi _{2}^{}(ae^{i\theta },ce^{i\theta },f/g;ac,fe^{i\theta };q^2,ge^{i\theta }),|g|<1\end{array}$$
With $`a=q^{2r+1}`$, $`c=d=q`$, $`f=q^{2r+3}z^1`$ and $`g=qz^1`$, the rhs of (196) coincides with the sum (193) while the rhs becomes
(197)
$$\frac{(qg,fe^{i\theta };q^2)_{\mathrm{}}}{(qf,ge^{i\theta };q^2)_{\mathrm{}}}{}_{2}{}^{}\varphi _{1}^{}(q^{2r+1}e^{i\theta },qe^{i\theta };q^{2r+3}z^1e^{i\theta };q^2,qz^1e^{i\theta })$$
The basic hypergeometric series in (197) has the form $`{}_{2}{}^{}\varphi _{1}^{}(A,B;C;q^2,C/AB)`$. If $`|C/AB|<1`$, the latter can be summed according to eq. 1.5.1 in :
(198)
$$_2\varphi _1(A,B;C;q^2,C/AB)=\frac{(C/A,C/B;q^2)_{\mathrm{}}}{(C,C/AB;q^2)_{\mathrm{}}}$$
Use of (198) in (197) and collecting terms yields
(199)
$$\begin{array}{c}r(y,z)=\frac{1}{\stackrel{~}{P}_0(y;q^{2r1},q^1,q^1|q^2)}q^r\left(1q^2z^1\right)\frac{(q^2z^1,q^2z^1;q^2)_{\mathrm{}}}{(qz^1e^{i\theta },qz^1e^{i\theta };q^2)_{\mathrm{}}},\hfill \\ \hfill |qz^1|<1\end{array}$$
Setting $`\stackrel{~}{P}_0(y;q^{2r1},q^1,q^1|q^2)=q^r`$ leads to (142). Equations (196) and (198) allow to handle the contributions to the terms $`v_{0,s}^{}v_{2,r+s}`$ on the rhs of (141) or $`v_{0,r+s}^{}v_{2,s}`$ on the rhs of (140) in a similar way. The condition $`qz^1<1`$ proves sufficient to reproduce the result (142) in all cases. |
warning/0507/gr-qc0507002.html | ar5iv | text | # Deviation equations in spaces with torsion
( Short title: Deviation equations in spaces with torsion
| Basic ideas$`:`$ | 1980–1981 |
| --- | --- |
| Last update$`:`$ | June 28, 2005 |
| Produced$`:`$ | |
http://www.arXiv.org e-Print archive No.$`:`$ gr-qc/0507002
Published in
Proceedings of the 5-th Soviet (USSR) Gravity Conference
“Modern theoretical and experimental problems of relativity theory
and gravitation”, Moscow Univ. Press, Moscow, 1981, p. 122 (In Russian)
BO/$``$$``$ HO$`^{}_{}{}^{\text{®}}\text{}`$
Subject Classes:
General relativity, Differential geometry
Key-Words:
Deviation equations
Spaces with affine (linear) connection, Spaces with torsion
)
Abstract. The most general form of the deviation equations in spaces with linear connection with arbitrary torsion is derived.
1. The deviation equations of two trajectories (paths) in a Riemannian space $`V_n`$ find applications in number of problems in the gravity theory \[Yano/LieDerivatives\]. The method of their derivation can be generalized to spaces $`L_n`$ with affine (linear) connection $``$, which is, of course, applicable to spaces with torsion $`U_n`$.
2. For any vector $`\xi T_x(L_n)`$ in the space $`T_x(L_n)`$ tangent to $`L_n`$ at $`xL_n`$, the following identity holds:
$$\frac{D^2\xi }{\mathrm{d}s}=\widehat{R}(u,\xi )+C_1^1(FD\xi )+\frac{D\widehat{T}(u,\xi ))}{\mathrm{d}s}T(F,\xi )+𝔏_\xi V\frac{D𝔏_\xi u}{\mathrm{d}s}C_1^2(Du𝔏_\xi u).$$
(1)
Here: $`\frac{D}{\mathrm{d}s}:=u^i_{E_i}`$ $`u^i:=\frac{\mathrm{d}x^i(s)}{\mathrm{d}s}`$, $`_{E_i}`$ is the covariant derivative along $`E_i`$ with $`\{E_i`$} being a frame in a neighborhood of $`x`$, $`u=u^iE_i`$, $`Du:=(_{E_i}u)\mathrm{\Theta }^i`$ with $`\{\mathrm{\Theta }^i\}`$ being a frame dual to $`\{E_i\}`$ ($`\mathrm{\Theta }^i(E_j)=\delta _j^i`$), $`s`$ is a parameter of a path $`\gamma (s)L_n`$, $`s`$, $`\widehat{R}`$ is the curvature operator with $`\widehat{R}(u,\xi )v:=\left(_u_\xi _\xi _u_{[u,\xi ]}\right)v=R_{jkl}^iv^ju^k\xi ^lE_i,`$ $`F:=\frac{Du}{\mathrm{d}s}=u^i_{E_i}u`$, $`\widehat{T}`$ is the torsion operator with $`\widehat{T}(u,\xi ):=_u\xi _\xi u[u,\xi ]=T_{jk}^iu^j\xi ^kE_i,`$ $`𝔏_\xi `$ is the Lie derivative along $`\xi `$, and $`[u,\xi ]:=u\xi \xi u`$. The identity (1) in component form reads
$$\frac{\overline{D}^2\xi ^k}{\mathrm{d}s}=R_{ijl}^ku^iu^j\xi ^l+\xi _{|j}^kF^j+u^j\frac{\overline{D}(T_{jl}^k\xi ^l)}{\mathrm{d}s}+𝔏_\xi F^kT(F,\xi )+\frac{\overline{D}(𝔏_\xi u)^k}{\mathrm{d}s}u_{|i}^k𝔏_\xi u^i,$$
(2)
where $`\frac{Dv}{\mathrm{d}s}:=u^i\overline{}_{E_i}`$, $`_{E_i}v=:(\overline{}_{E_i}v^j)E_j`$ and $`v_{|j}^i:=)_{E_j}v^i`$ is the covariant derivative of the components $`v^i`$ along $`E_i`$.
3. If we impose some additional restrictions, which in certain cases may be first integrals of the deviation equation, on the quantities entering into (1) (e.g. $`F=0`$, $`𝔏_\xi u=0`$, $`𝔏_\xi F=F`$), then we get corresponding deviation equations in space with torsion $`U_n`$, in which the vector $`\xi `$ is considered as an infinitesimal one. |
warning/0507/cond-mat0507320.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The bending energy of a two-dimensional surface, quadratic in its extrinsic curvature, is invariant under scaling; size does not matter. What is less obvious is that this energy is also invariant under transformations of the three-dimensional background which preserve angles; it is conformally invariant. Any two surface geometries related to each other by inversion in a point have the same energy. This property was first studied systematically by Willmore in the 60s . More recently, it was discovered to lie at the heart of some fascinating connections between differential geometry and integrable systems . At mesoscopic scales, the physics of a fluid membrane — formed by the spontaneous aggregration of amphiphilic molecules into bilayers in water — is captured completely by its geometrical degrees of freedom . On such scales the membrane can be modeled as a two-dimensional surface; at lowest order, the principal cost in energy is associated with bending this surface . Remarkably, all of the molecular details get telescoped into a single rigidity modulus. A role was also found for a relativistic counterpart in the eighties: an addition quadratic in extrinsic curvature to the action of a relativistic string accounts for the behaviour of colour flux tubes in QCD .
Higher-dimensional analogues of the two-dimensional bending energy are also of potential interest both as statistical field theories, and relativistically, as braneworld actions. If the surface dimension differs from two, however, the energy quadratic in extrinsic curvature is no longer scale invariant much less conformally invariant: a higher-dimensional sphere without a constraint on its area will collapse; tension is necessarily introduced. Conformal invariant energies, polynomial in the extrinsic curvature, are however simple to construct: the building block is the traceless part of the extrinsic curvature tensor which transforms by a multiplicative factor. In four dimensions, for a hypersurface of fixed topology, there are two independent conformal invariant energies quartic in the extrinsic curvature.The snag is that these invariants alone cannot accurately describe a conformally invariant theory of bending. The reason is simple: when expanded as a power series in terms of a height function, they begin with a term quartic in this function. Thus, they vanish in the Gaussian approximation to the energy truncating it at the quadratic in the height function. In particular, there is no harmonic regime to describe fluctuations about a flat geometry. On dimensional grounds, the relevant invariant must involve curvature gradients. This invariant will play a role in the formulation of a consistent statistical field theory of four-dimensional hypersurfaces. Its identification, however, is somewhat less trivial than that of its quartic counterparts. To do this, it will be useful to approach the problem from a global point of view which focuses directly on the transformation properties of the surface energy rather than the individual tensors which appear within it. For the sake of simplicity, we will focus on a closed hypersurface of fixed topology; more simple still, think topological four-sphere.
Consider any energy, constructed using the metric and the extrinsic curvature, which is invariant under reparametrizations of the hypersurface and Euclidean motions of space. If the theory is invariant under translations the response of this energy to an arbitrary infinitesimal deformation of the hypersurface can be expressed in terms of a stress tensor . In it was shown that the tangential stress $`f^{ab}`$ has two contributions: one of these is the metric stress tensor $`T^{ab}`$ which determines the response of the energy to changes in the intrinsic geometry; the second, which determines the response to changes in the extrinsic geometry, involves the functional derivative of the energy with respect to the extrinsic curvature, $`^{ab}`$. One is interested, in particular, in determining the response of the energy to the deformation of the hypersurface induced by an infinitesimal conformal transformation. It is possible to characterize this response in a remarkably succinct way in terms of traces: that of $`f^{ab}`$ and that of $`^{ab}`$. Conformal invariance will place constraints on these traces. In contrast to a conformal invariant of the intrinsic geometry which has a representation with $`^{ab}=0`$, these constraints may be satisfied in a very subtle way by an invariant of the the extrinsic geometry.
The task is to identify energies that are consistent with these constraints. While the focus will be on closed four-dimensional hypersurfaces, the techniques developed will be independent of the dimension. We first briefly describe the construction, within this framework, of the two well-known four-dimensional conformally invariant energies quartic in extrinsic curvature. Modulo the Gauss-Codazzi equations, which identify the intrinsic Riemann tensor with a quadratic in the extrinsic curvature, one of these invariants is the Weyl invariant associated with the intrinsic geometry of the hypersurface, and thus insensitive to the particular way the hypersurface is embedded. The third invariant involves a balance of a part quadratic in curvature gradients with a quartic in curvature; neither term on its own is conformally invariant. We show how this constraint can be satisfied by tuning the quartic so that the two trace terms cancel. Intriguingly, this cancellation involves properties of the four-dimensional Willmore energy quadratic in the extrinsic curvature (which is not itself a conformal invariant in this dimension) in an essential way.
## 2 Linear response, the Euler-Lagrange derivative as a divergence, and the stress
Consider a closed $`D`$-dimensional hypersurface embedded in $`R^{D+1}`$. This hypersurface is described locally by the embedding, $`𝐱=𝐗(\xi ^a)`$. Here $`𝐱=(x^1,\mathrm{},x^{D+1})`$ and $`\xi ^a`$, $`a=1,\mathrm{},D`$ parametrize the hypersurface. The metric tensor and extrinsic curvature induced by $`𝐗`$ are respectively $`g_{ab}=𝐞_a𝐞_b`$ and $`K_{ab}=𝐞_a_b𝐧`$, where $`𝐞_a=_a𝐗`$, $`a=1,\mathrm{},D`$ are tangent and $`𝐧`$ is the unit normal. The Gauss-Weingarten equations are $`_a𝐞_b=K_{ab}𝐧`$ and $`_a𝐧=K_a{}_{}{}^{b}𝐞_{b}^{}`$ . $`_a`$ is the covariant derivative compatible with $`g_{ab}`$; spatial indices get raised with the inverse metric $`g^{ab}`$. We are interested in functionals of $`𝐗`$ which are invariant under reparametrizations of the hypersurface.
The metric and extrinsic curvature are both invariant under the change in $`𝐗`$ induced by a Euclidean motion in $`R^{D+1}`$: a Euclidean invariant energy $`H[𝐗]`$ can therefore be cast as a functional of the metric, the extrinsic curvature and its derivatives,
$$H[𝐗]=𝑑A(g_{ab},K_{ab},_aK_{bc},\mathrm{}).$$
(1)
The area element induced on the hypersurface is $`dA=\sqrt{\mathrm{det}g_{ab}}d^D\xi `$. We wish, in particular, to construct an energy which is invariant under deformations induced by a conformal change of the Euclidean background:
$$\delta 𝐱=𝐚+\mathrm{𝐁𝐱}+\lambda 𝐱+𝐱^2𝐜2(𝐜𝐱)𝐱,$$
(2)
where $`𝐚`$ and $`𝐜`$ are two constant vectors ($`𝐜`$ has dimensions of inverse length), $`𝐁`$ is an antisymmetric $`(D+1)\times (D+1)`$ matrix, and $`\lambda `$ is a positive constant. $`\delta 𝐱`$ is the sum of an infinitesimal Euclidean motion, a change of scale and a special conformal transformation. The latter exponentiates to the composition of an inversion $`𝐱𝐱/𝐱^2`$, a translation through the vector $`𝐜`$ , and another inversion, $`𝐱(𝐱+𝐱^2𝐜)/(1+2𝐜𝐱+𝐜^2𝐱^2)`$. Both $`g_{ab}`$ and $`K_{ab}`$ are invariant under Euclidean motions: as a result, any energy of the form (1) will also be by construction. What one now needs to do is characterize the constraints placed on $`H`$ by invariance under scaling and special conformal transformations.
It is useful to first determine the linear response of $`H`$ to any small deformation of the hypersurface. This task is simplified by exploiting invariance under Euclidean motions of the ambient space. While Noether’s theorem informs us that the Euler-Lagrange derivative can always be cast as the divergence of a stress tensor, in all but the simplest case — an energy proportional to the area functional describing surface tension — the identification of this tensor is subtle: unlike the stress associated with area, the stress will depend not only on the intrinsic geometry but also on how the hypersurface bends. The tug on the hypersurface will possess a normal component.
A small deformation of the hypersurface is described by the infinitesimal change in the embedding functions $`𝐗`$
$$𝐗𝐗+\delta 𝐗.$$
(3)
Note the following points:
(1) as a consquence of the reparametrization invariance of $`H`$ in a closed geometry, the response of $`H`$ is independent of the tangential projection of $`\delta 𝐗`$; thus
$$\delta H=𝑑A𝐧\delta 𝐗$$
(4)
involves only the normal projection. The Euler-Lagrange derivative of $`H`$ with respect to $`𝐗`$ is denoted by $`𝐧`$.
(2) the translational invariance of $`H`$ implies that its Euler-Lagrange derivative is a divergence
$$𝐧=_a𝐟^a.$$
(5)
The hypersurface current $`𝐟^a𝐚`$ is associated with the invariance of $`H`$ under a translation $`\delta 𝐗=𝐚`$. When the Euler-Lagrange equation $`=0`$ is satisfied, this current is conserved. The closure of the geometry then permits $`\delta H`$ to be recast in the remarkably simple form
$$\delta H=𝑑A𝐟^a_a\delta 𝐗.$$
(6)
This expression involves one less derivative than Eq.(4). Note that one does not need to know how $`𝐟^a`$ itself transforms. Eq.(6) is valid whether or not the Euler-Lagrange equation are satisfied. This equation will be used to examine the response of $`H`$ to the deformation in the hypersurface induced by conformal transformations of space.
## 3 The stress
The stress $`𝐟^a`$ associated with $`H`$ is given by
$$𝐟^a=(T^{ab}^{ac}K_c{}_{}{}^{b})𝐞_b_b^{ab}𝐧,$$
(7)
where $`^{ab}`$ is the functional derivative of $`H`$ with respect to $`K_{ab}`$,
$$^{ab}=\frac{}{K_{ab}}_c\left(\frac{}{_cK_{ab}}\right)+\mathrm{}$$
(8)
and $`T^{ab}=(2/\sqrt{g})\delta H/\delta g_{ab}`$ is the intrinsic stress tensor associated with the metric $`g_{ab}`$. This construction involves treating $`g_{ab}`$ and $`K_{ab}`$ as independent variables in $``$; to do this consistently requires one to introduce a set of auxiliary variables to constrain $`g_{ab}`$ and $`K_{ab}`$ to satisfy the Gauss-Weingarten structural relationships. The ellipsis appearing on the r.h.s of Eq.(8) indicates terms which appear if $``$ depends on derivatives of $`K_{ab}`$ higher than first. A simple derivation of Eq.(7) is provided in . We note, in particular, that $`𝐟^a`$ decomposes into tangential and normal parts:
$$𝐟^a=f^{ab}𝐞_b+f^a𝐧.$$
(9)
This decomposition has the following properties which are relevant:
(1) The tangential projections of Eq.(5) provide a consistency condition on the components of the stress
$$_af^{ab}+K^{ab}f_a=0;$$
(10)
the normal component determines $``$:
$$=_af^aK_{ab}f^{ab}.$$
(11)
(2) In general, the normal stress $`f^a`$ is a divergence.
(3) Even though both $`^{ab}`$ and $`T^{ab}`$ are symmetric tensors, $`f^{ab}`$ will not generally be symmetric. On one hand, as Eq.(11) indicates clearly, only the symmetric part of $`f^{ab}`$ contributes to the Euler-Lagrange derivative. An anti-symmetric contribution, if present, will however show up in the consistency conditions (10) and so cannot be discarded naively.
Let us now consider specific forms for the function $``$ appearing in Eq.(1) which will be used in the construction of conformal invariants.
### 3.1 $`(g_{ab},K_{ab})`$
Suppose that $``$ does not involve derivatives of $`K_{ab}`$: $`=(g_{ab},K_{ab})`$. Then
$$^{ab}=/K_{ab},$$
(12)
and it is simple to show that
$$T^{ab}=2^{ac}K_c{}_{}{}^{b}g^{ab},$$
(13)
so that
$$𝐟^a=(^{ac}K_c{}_{}{}^{b}g^{ab})𝐞_b_b^{ab}𝐧.$$
(14)
Note that $`f^{ab}`$ is manifestly symmetric. A straighforward calculation gives
$$=_a_b^{ab}K_{ac}K^c{}_{b}{}^{}_{}^{ab}+K,$$
(15)
where $`K=g^{ab}K_{ab}`$ is the trace of $`K_{ab}`$ ($`D`$ times the mean curvature).
In particular, for the Canham-Helfrich or Willmore energy, one has
$$H_0=\frac{1}{2}𝑑AK^2,$$
(16)
$`=K^2/2`$ and $`^{ab}=Kg^{ab}`$, so that
$$𝐟^a=K(K^{ab}\frac{1}{2}g^{ab}K)𝐞_b_aK𝐧,$$
(17)
and (see also )
$$=^2K+\frac{1}{2}K(K^22K^{ab}K_{ab}),$$
(18)
If $`K=0`$, then $`=0`$. The relationship between $``$ and $`𝐟^a`$ for $`H_0`$ will play a role in the construction of a higher-derivative conformal invariant of a four dimensional hypersurface.
Note that it is unnecessary to admit an explicit intrinsic curvature dependence in $``$. This is because the Gauss-Codazzi equations
$$_{abcd}=K_{ac}K_{bd}K_{ad}K_{bc}$$
(19)
completely fix the Riemann tensor, as well as its contractions, the Ricci tensor $`_{ab}=g^{cd}R_{acbd}`$ and the scalar curvature $`=g^{ab}_{ab}`$, in terms of the extrinsic curvature.<sup>1</sup><sup>1</sup>1The Riemann tensor is defined intrinsically by the failure of the $`_a`$ to commute: for a space vector $`V_a`$, we have the Ricci identity $`(_a_b_b_a)V_c=_{abc}{}_{}{}^{d}V_{d}^{}`$. However, if one is interested explicitly in a functional of the intrinsic geometry, $`=(g_{ab},_{abcd},_e_{abcd},\mathrm{})`$, it may then be more appropriate to treat these tensors as functionals of $`g_{ab}`$ alone, and ignore the integrability conditions (19). If this is done, $`^{ab}=0`$ and $`T^{ab}`$ is the stress tensor of the (purely) metric theory defined by $``$. Now, $`f^{ab}=T^{ab}`$ and it is manifestly symmetric; furthermore $`f^a=0`$.<sup>2</sup><sup>2</sup>2Intrinsically defined invariants are not the only ones possessing this property. One can show that the geometrical invariants constructed out of the symmetric polynomials in the curvature, $`P_N(\sigma _1,\mathrm{},\sigma _D)`$, $`ND`$, where $`\{\sigma _i\}`$ are the principal curvatures, do also. For example, when $`N=1`$, $`P_1=K`$ and $`f^{ab}=K^{ab}g^{ab}K`$, which is conserved, so that $`f^a=0`$. The Euler-Lagrange derivative is then simply $`=K_{ab}T^{ab}`$; the consistency condition then reads $`_aT^{ab}=0`$ — the metric stress tensor is conserved. Clearly, it does not matter how one decides to split the burden on $`g_{ab}`$ and $`K_{ab}`$, so long as it is done consistently when performing the variations in the derivation of $`𝐟^a`$. As discussed in detail elsewhere, if $`𝐟^a`$ is treated as a differential form, the difference between its value in the two representations is an exact form.
### 3.2 $`(g_{ab},K_{ab},_cK_{ab})`$
If one extends the class of functionals to include a dependence on $`_cK_{ab}`$, there are few useful general statements concerning the structure of $`𝐟^a`$. Our limited goal, however, is to identify conformal invariants of closed hypersurfaces so we do not need to consider the most general form.
Consider candidate polynomials in $`_aK_{bc}`$ and $`K_{ab}`$ that are consistent with scale invariance. When $`D=2`$, there are none. When $`D=3`$, there is a Chern-Simons type topological energy; it vanishes on a closed geometry. When $`D=4`$, the quadratics in $`_cK_{ab}`$ are scale invariant. As shown in , however, any quadratic in derivatives of $`K_{ab}`$ is expressible, modulo a divergence, as a sum of the simple invariant
$$H_1=\frac{1}{2}𝑑A(K)^2$$
(20)
and an integral over some quartic in $`K_{ab}`$. The latter is of the form, $`(g_{ab},K_{ab})`$, already considered in section (3.1). So $`H_1`$ is the only invariant that needs to be considered.
The demonstration of this claim involves the Codazzi-Mainardi integrability conditions,
$$_aK_{bc}_bK_{ac}=0,$$
(21)
as well as the Ricci identity applied to $`K_{ab}`$
$$[_a,_b]K_{cd}=R_{abc}{}_{}{}^{f}K_{fd}^{}+R_{abd}{}_{}{}^{f}K_{cf}^{}.$$
(22)
Consider the energy constructed using the quadratic $`_aK_{bc}^aK^{bc}`$. One first uses Eq.(21) followed by an integration by parts to obtain
$$dA(_aK_{bc})(^aK^{bc})=dA(_aK_{bc})(^bK^{ac})=dAK^b{}_{c}{}^{}_{a}^{}_bK^{ac}.$$
(23)
One then makes use of Eq.(22) to switch derivatives so that
$`{\displaystyle 𝑑A(_aK_{bc})(^aK^{bc})}`$ $`=`$ $`{\displaystyle }dA(K_c{}_{}{}^{b}_{b}^{}_aK^{ac}_{abcd}K^{ac}K^{bd}+_{ab}K^{ac}K_c{}_{}{}^{b}).`$ (24)
The contracted Codazzi-Mainardi equations, $`_aK^{ab}^bK=0`$, and another integration by parts are applied to the first term to nudge it into the required form:
$$dAK_c{}_{}{}^{b}_{b}^{}_aK^{ac}=dA_bK^b{}_{c}{}^{}_{a}^{}K^{ac}=dA^cK_cK.$$
(25)
One concludes that
$`{\displaystyle 𝑑A(_aK_{bc})(^aK^{bc})}`$ $`=`$ $`{\displaystyle }dA(^cK_cK+_{abcd}K^{ac}K^{bd}_{ab}K^{ac}K_c{}_{}{}^{b})`$ (26)
$`=`$ $`{\displaystyle 𝑑A\left(^cK_cK+(\mathrm{tr}K^2)^2K\mathrm{tr}K^3\right)}.`$
The notation $`\mathrm{tr}K^n=K_{a_1}{}_{}{}^{a_2}\mathrm{}K_{a_n}^{a_1}`$ has been introduced. On the second line, the Gauss-Codazzi equations (19) have been used to eliminate the Riemann tensor in favour of a quadratic in extrinsic curvature. The energy $`H_1`$, given by Eq.(20), is reproduced modulos a quartic in extrinsic curvature.
Note that for $`H_1`$ one has (this is true for any dimension $`D`$),
$$^{ab}=_c(g^{ab}^cK)=g^{ab}^2K,$$
(27)
and
$$T^{ab}=^aK^bK\frac{1}{2}g^{ab}(K)^22K^{ab}^2K.$$
(28)
The second derivative term originates in the variation of $`_aK`$ with respect to $`g_{ab}`$. The correponding stress tensor is
$$𝐟^a=\left[^aK^bK\frac{1}{2}g^{ab}(K)^2K^{ab}^2K\right]𝐞_b+^a^2K𝐧.$$
(29)
Again $`f^{ab}`$ is symmetric.<sup>3</sup><sup>3</sup>3Evidently, one has to proceed to a relatively high order energy to produce an $`f^{ab}`$ which is not symmetric: for $`=K^{ab}_aK_bK`$, $`^{ab}`$ does not commute with $`K_{ab}`$ and thus $`f^{ab}`$ is not symmetric. In this case, it is simple to check that
$$=[^2+\mathrm{tr}K^2]^2KK^{ab}[_aK_bK\frac{1}{2}g_{ab}(K)^2].$$
(30)
Note that if $`K=0`$, then $`=0`$ so that minimal hypersurfaces also minimize $`H_1`$.
## 4 Scaling
On one hand, it is trivial to identify energies which are scale invariant. On the other, the imprint of scale invariance on the stress tensor is subtle and it will be relevant to our interpretation of the response, under special conformal transformations, of the energy in terms of the stress tensor.
Consider an energy with a fixed scaling dimension. Under a change of scale $`𝐗\mathrm{\Lambda }𝐗`$, where $`\mathrm{\Lambda }`$ is a positive constant, one has
$$H[\mathrm{\Lambda }𝐗]=\mathrm{\Lambda }^{D+d}H[𝐗]$$
(31)
for some $`d`$, or alternatively, in terms of the corresponding density, $`[\mathrm{\Lambda }𝐗]=\mathrm{\Lambda }^d[𝐗]`$. $`H`$ is scale invariant when $`d=D`$.
Consider now an infinitesimal change of scale, $`\mathrm{\Lambda }=1+\lambda `$; at first order in $`\lambda `$, Eq.(31) gives
$$\delta _\lambda H=(D+d)\lambda H.$$
(32)
On the other hand, the first order variation Eq.(6) with the substitution $`\delta 𝐗=\lambda 𝐗`$ expresses $`\delta _\lambda H`$ in terms of the trace of the tangential stress,
$$\delta _\lambda H=\lambda dAf^a{}_{a}{}^{},$$
(33)
where $`f^a{}_{a}{}^{}=g_{ab}f^{ab}`$. Comparison of Eq.(32) with Eq.(33) furnishes an identity,
$$(D+d)H=dAf^a{}_{a}{}^{}.$$
(34)
Only the trace of the tangential stress tensor contributes to the change of $`H`$ under scaling. Locally, this implies that
$$f^a{}_{a}{}^{}=(D+d)+_aG^a,$$
(35)
where $`G^a`$ is a hypersurface vector field. Modulo a divergence, the trace is proportional to the integrand.
For functionals of the form $`(g_{ab},K_{ab})`$ one can show that $`G^a=0`$. <sup>4</sup><sup>4</sup>4To see this, consider the response to a deformation of $`H`$ on a region with boundary: For functionals of the form $`(g_{ab},K_{ab})`$, Eq.(6) is replaced by
$$\delta H=𝑑A𝐟^a_a\delta 𝐗+𝑑A_a[^{ab}𝐞_b\delta 𝐧].$$
(36) In particular, under a change of scale, $`\delta 𝐧=0`$ and Eq.(36) reproduces Eq.(33) without any boundary term. In particular, a scale invariant functional of this form has vanishing tangential trace: $`f^a{}_{a}{}^{}=0`$. This is not true of higher derivative scale invariants: Eq.(34) does, however, imply that the trace is the divergence of a hypersurface vector field:
$$f^a{}_{a}{}^{}=_aG^a.$$
(37)
In particular, in the case of the functional $`H_1`$ defined by Eq.(20), inspection of Eq.(29) gives for the trace of $`f_{ab}`$,
$$f^a{}_{a}{}^{}=\frac{4D}{2}(K)^2^2K^2/2.$$
(38)
One identifies $`G^a=^aK^2/2`$.
## 5 Special conformal transformations
Infinitesimally, a special conformal transformation induces a change in $`𝐗`$ given by
$$\delta _𝐜𝐗=𝐗^2𝐜2(𝐜𝐗)𝐗.$$
(39)
The corresponding response of the energy is determined using Eq.(6):
$`\delta _𝐜H`$ $`=`$ $`{\displaystyle 𝑑A𝐟^a_a\delta _𝐜𝐗}`$ (40)
$`=`$ $`{\displaystyle 𝑑A\left[f^{ab}(𝐞_b_a\delta 𝐗)+f^a(𝐧_a\delta 𝐗)\right]}`$
$`=`$ $`2{\displaystyle }dA[f^a{}_{a}{}^{}(𝐜𝐗)f^a𝐜𝐟_{0a}]`$
$`+2{\displaystyle }dAf^{ab}[(𝐞_a𝐜)(𝐞_b𝐗)(ab)].`$
where
$$𝐟_0^a=(𝐞^a𝐗)𝐧(𝐧𝐗)𝐞^a.$$
(41)
The identities
$$𝐞_a_b\delta _𝐜𝐗=2[(𝐞_a𝐜)(𝐞_b𝐗)(ab)]2(𝐜𝐗)g_{ab},$$
(42)
and
$$𝐧_a\delta _𝐜𝐗=2\left[(𝐞_a𝐗)(𝐜𝐧)(𝐞_a𝐜)(𝐗𝐧)\right]=2𝐜𝐟_{0a}$$
(43)
have been used on the third line of Eq.(40). In the case of any energy we will consider $`f^{ab}`$ is symmetric so that the term appearing on the last line in Eq.(40) vanishes. The equation thus simplifies to
$$\delta _𝐜H=2dA[f^a{}_{a}{}^{}(𝐜𝐗)f^a𝐜𝐟_{0a}].$$
(44)
Further simplification is possible using the structure of $`𝐟^a`$. Using the fact that $`f^a=_b^{ab}`$, where $`^{ab}`$ is given by Eq.(8), the second term appearing on the r.h.s of Eq.(44) can be cast as
$$𝑑Af^a𝐜𝐟_{0a}=𝑑A^{ab}𝐜_b𝐟_{0a}.$$
(45)
However, the definition of $`𝐟_0^a`$ (41) gives
$$_b𝐟_{0a}=g_{ab}𝐧+K_b{}_{}{}^{c}((𝐗𝐞_a)𝐞_c(ac)).$$
(46)
The tangential projection is anti-symmetric in $`a`$ and $`b`$ and so does not contribute to the r.h.s of Eq.(45) if $`f^{ab}=T^{ab}^{ac}K_c^b`$ is symmetric. Even when it is not, it cancels against an identical term appearing on the last line in Eq.(40). There follows the identity
$$\delta _𝐜H=2dA[f^a{}_{a}{}^{}(𝐜𝐗)^a{}_{a}{}^{}(𝐜𝐧)],$$
(47)
where $`^a{}_{a}{}^{}=g_{ab}^{ab}`$. The response of $`H`$ to an infinitesinal special conformal transformation has been expressed as a difference of two terms. Each of therse terms involves a trace. Note that in the case of any intrinsic geometrical invariant, there exists a representation in which the second term vanishes.
The energy $`H`$ is conformally invariant if and only if Eqs.(37) and
$$f^a{}_{a}{}^{}(𝐜𝐗)^a{}_{a}{}^{}(𝐜𝐧)=_ah^a,$$
(48)
are satisfied. $`h^a`$ is a hypersurface vector field.
We now construct a four-dimensional energy, polynomial in the curvature and its derivatives, that is consistent with these constraints.
## 6 Conformally invariants polynomial in $`K_{ab}`$
A scale invariant energy with a density $``$ depending on $`g_{ab}`$ and $`K_{ab}`$ but not on their derivatives has traceless $`f^{ab}`$: $`f^a{}_{a}{}^{}=0`$. To be invariant under special conformal transformations, one also requires that
$$𝑑A^a{}_{a}{}^{}𝐧=0.$$
(49)
This is clearly satisfied if $`^{ab}`$ is also traceless: $`^a{}_{a}{}^{}=0`$. It is straightforward to construct polynomial functionals with this property. These are the well-known invariants involving the traceless part of the extrinsic curvature tensor $`\stackrel{~}{K}_{ab}`$ ($`\stackrel{~}{K}^a{}_{a}{}^{}=0`$)
$$\stackrel{~}{K}_{ab}=K_{ab}\frac{K}{D}g_{ab}.$$
(50)
Let $``$ be a product of terms $`_n`$, each of which is a trace over a product of $`n`$ $`\stackrel{~}{K}_{ab}`$’s (see definition below Eq.(26)); $`_n=\mathrm{tr}\stackrel{~}{K}^n`$. Note that
$$\mathrm{\Pi }_{ab}{}_{}{}^{cd}:=\frac{\stackrel{~}{K}_{ab}}{K_{cd}}=\frac{1}{2}(\delta _a{}_{}{}^{c}\delta _{b}^{}{}_{}{}^{d}+\delta _a{}_{}{}^{d}\delta _{b}^{}{}_{}{}^{c})\frac{1}{D}g_{ab}g^{cd}$$
(51)
projects out the trace ($`\mathrm{\Pi }_{ab}{}_{}{}^{cd}g_{cd}^{}=0`$). For each factor $`_n`$, we thus find that
$$g_{ab}\frac{_n}{K_{ab}}=0.$$
(52)
Consequently, $`^a{}_{a}{}^{}=0`$ (see also ). Two-dimensional surfaces are considered in an appendix.
In four dimensions, there are two polynomial conformal invariants constructed this way corresponding to the two independent quartics $`\mathrm{tr}\stackrel{~}{K}^4`$ and $`(\mathrm{tr}\stackrel{~}{K}^2)^2`$. One linear combination of the two is the Weyl invariant of the intrinsic geometry .
It is possible to satisfy Eq.(49) in a rather less trivial way. It was seen that the translation invariance of any functional of the form (1) implies the identity Eq.(5) between its Euler-Lagrange derivative and the hypersurface divergence of a stress tensor. If this identity is integrated over a closed hypersurface it follows immediately that the Euler-Lagrange derivative satisfies
$$𝑑A𝐧=0.$$
(53)
Thus, if it is possible to cast $`^a_a`$ as the Euler-Lagrange derivative of some translationally invariant functional of the form (1), then Eq.(49) will be satisfied. However, one can show that the only energy densities constructed using $`K_{ab}`$ consistent with this condition are proportional to a sum of the symmetric polynomials in the principal curvatures. Consisitently with scale invariance leaves only the determinant. Thus the only conformal invariant generated this way is the Gauss-Bonnet topological invariant .
Note that the Paneitz invariant, which has been the centre of recent research, is the difference between the Gauss-Bonnet and the Weyl invariants . As such, it is an invariant of the intrinsic geometry.
## 7 Conformal invariant quadratic in gradients of $`K_{ab}`$
In four dimensions, the functional $`H_1`$ defined by Eq.(20) is scale invariant but it is not conformally invariant. It is possible, however, to construct a conformal invariant by adding to $`H_1`$ the integral of an appropriate quartic in $`K_{ab}`$. One way to identify this quartic is as follows:
(i) First determine how $`H_1`$ transforms:
For $`H_1`$, neither $`f^a_a`$ nor $`^a_a`$ vanishes; one must contend with the two terms appearing in Eq.(47). For $`H_1`$, Eq.(38) gives $`f^a{}_{a}{}^{}=^2K^2/2`$. On performing two integrations by parts and using the Gauss-Weingarten equations, one finds that
$$dAf^a{}_{a}{}^{}(𝐜𝐗)=\frac{1}{2}dAK^3(𝐜𝐧).$$
(54)
In addition, using $`^{ab}=g^{ab}^2K`$,
$$dA^a{}_{a}{}^{}(𝐜𝐧)=4dA^2K(𝐧𝐜).$$
(55)
The identities (54) and (55) are now substituted into Eq.(47) to give for $`\delta _𝐜H_1`$
$$\delta _𝐜H_1=𝑑A(K^3+8^2K)(𝐧𝐜).$$
(56)
(ii) Next note that the r.h.s of Eq.(56) can be simplified by using an identity associated with the quadratic energy $`H_0`$ defined by Eq.(16). Using Eq.(18), Eq.(5) implies that
$$𝑑A\left[^2K+(\mathrm{tr}K^2\frac{1}{2}K^2)K\right](𝐧𝐜)=0,$$
(57)
for any $`𝐜`$. This identity allows $`\delta _𝐜H_1`$ given by Eq.(56) to be expresssed in terms of a cubic polynomial in $`K_{ab}`$:
$$\delta _𝐜H_1=8𝑑A(\mathrm{tr}K^2\frac{5}{8}K^2)K(𝐧𝐜).$$
(58)
Now let $`_2=(K)^2/2+^{}`$ where $`^{}`$ is quartic in $`K_{ab}`$. If
$$g_{ab}\frac{^{}}{K_{ab}}=8(\mathrm{tr}K^2\frac{5}{8}K^2)K,$$
(59)
then $`H=𝑑A`$ will be conformally invariant. The choice of $`^{}`$ is clearly not unique: however, it is modulo a linear combination of the two conformally covariant quartics, $`\mathrm{tr}\stackrel{~}{K}^4`$ and $`(\mathrm{tr}\stackrel{~}{K}^2)^2`$. The simplest choice is a linear combination of the two invariants, $`K^4`$ and $`K^2\mathrm{tr}K^2`$. A short calculation gives
$$^{}=K^2\mathrm{tr}K^2\frac{7}{16}K^4,$$
(60)
and one identified the following four-dimensional conformally invariant energy:
$$H_2=\frac{1}{2}𝑑A\left((K)^2\frac{7}{8}K^4+2K^2\mathrm{tr}K^2\right).$$
(61)
This identification appears to be new. Unlike the two quartic conformal invariants in four-dimensions, the conformal invariance of $`H_2`$ involves a delicate balance between gradients and quartics.
With our choice of quartic $`^{}`$, $`H_2`$ is not positive, unlike the invariants constructed using $`\mathrm{tr}\stackrel{~}{K}^4`$ and $`(\mathrm{tr}\stackrel{~}{K}^2)^2`$, which are. It is, however, possible to form a positive conformally invariant gradient energy by adding an appropriate linear combination of the other invariants. The most general four-dimensional conformally invariant energy polynomial in the extrinsic curvature will involve a linear combination of all three invariants.
Note that a physically realistic four-dimensional generalization of the Willmore energy will involve $`H_2`$. Consider the Monge description of the hypersurface in terms of a height function $`h`$ above a reference plane. With respect to Cartesian coordinates on this plane, the extrinsic curvature tensor takes the form
$$K_{ab}=\frac{_a_bh}{(1+(h)^2)^{1/2}},$$
(62)
where now $`_a`$ is the flat derivative on this plane. To lowest order in $`h`$, $`K_{ab}_a_bh+𝒪(h^3)`$. In the Gaussian approximation, quadratic in $`h`$, all quartics in $`K_{ab}`$ vanish; in particular, the conformal invariants constructed using $`\mathrm{tr}\stackrel{~}{K}^4`$ and $`(\mathrm{tr}\stackrel{~}{K}^2)^2`$ both vanish. For $`H_2`$ given by Eq.(61) only the gradient term survives and one is left with
$$H_2=\frac{1}{2}𝑑A_{}(\mathrm{\Delta }h)^2+𝒪(h^4),$$
(63)
where $`dA_{}`$ is the area element and $`\mathrm{\Delta }`$ is the flat Laplacian on this plane. The conformal invariance of $`H_2`$ is, of course, necessarily mutilated in the approximation process.
## 8 Generalization
The construction of the invariant in Section (7) suggests that it may be useful to substitute Eq.(37) for $`f^a_a`$ into Eq.(47). If $`f^a_a`$ is the Laplacian of some scalar, $`f^a{}_{a}{}^{}=^2G`$, then one can perform two integrations by parts to re-express
$$dAf^a{}_{a}{}^{}(𝐜𝐗)=dAGK(𝐜𝐧).$$
(64)
It is now possible to peel off the space vector in Eq.(47) so that Eq.(48) is replaced by
$$(GK^a{}_{a}{}^{})𝐧=_a𝐅^a.$$
(65)
A sufficient condition for conformal invariance is that $`f^a{}_{a}{}^{}=^2G`$, and $`GK^a_a`$ is an Euler-Lagrange derivative. $`𝐅^a`$ is the corresponding stress tensor. This is clearly not the only way that conformal invariants can occur. In fact, if the energy depends only on the intrinsic geometry, then $`^{ab}`$ vanishes and there are no solutions of this form. An interesting exercise would be to identify other energies with an $`f^a_a`$ which is the Laplacian of a scalar.
## 9 Discussion
In the statistical field theory of surfaces, conformally invariant Hamiltonians provide fixed points of the renormalization group flow. While the theory of two-dimensional surfaces is well understood , next to nothing is known about possible four-dimensional counterparts. The identification of the appropriate invariants is a small first step towards the formulation of such a theory. Before plunging into statistical field theory, however, there are questions of an elementary nature that should be addressed. What are the minima of the conformally invariant energy? Even without constraints, highly non-trivial vacua appear to be admitted. Is there is a useful analogue of Willmore’s conjecture ? The classification of solutions will involve topological selection rules beyond the scope of this paper. It would also to know if four-dimensional analogues of Goetz and Helfrich’s egg carton geometries exist .
In the same way that the two-dimensional Willmore functional finds an application in relativistic field theory with the replacement of a Euclidean signature metric by a Lorentzian one, it is possible that the four-dimensional conformally invariant bending energy will find a role in braneworld cosmology . In this context, the generalization of the conformally invariant energy (61) to accommodate a curved bulk should be straightforward.
Finally, it should be noted that our construction has been framed in the language of classical differential geometry. Its translation into the language of Cartan’s exterior differential systems should be straightforward but useful, especially for addressing issues of a topolgical nature (see also ).
Acknowledgments
I thank Riccardo Capovilla, Markus Deserno and Martin Müller for helpful comments. I am also grateful to Denjoe O’ Connor for hospitality at DIAS Dublin where this work was begun. Partial support from CONACyT grant 44974-F is acknowledged.
## Appendix A Two-dimensional surfaces
When $`D=2`$, the only polynomial in $`\stackrel{~}{K}_{ab}`$, with scaling dimension $`d=2`$ and vanishing $`^a_a`$ is $`\mathrm{tr}\stackrel{~}{K}^2=\stackrel{~}{K}^{ab}\stackrel{~}{K}_{ab}`$. The corresponding energy is the Willmore energy.
The Gauss-Bonnet topological invariant,
$$dA\mathrm{det}K^a_b$$
(66)
is, of course, also a conformal invariant. For a two-dimensional surface $`\mathrm{det}K^a{}_{b}{}^{}=/2`$. In fact, any quadratic invariant in extrinsic curvature is trivially also conformally invariant. This is because any scalar in the extrinsic curvature can be expressed as a linear combination of $`\mathrm{tr}\stackrel{~}{K}^2=\mathrm{tr}K^2K^2/D`$ and $`=K^2\mathrm{tr}K^2`$ both of which give rise to conformal invariants when $`D=2`$. Thus, in particular, the two quadratics in $`K_{ab}`$, $`K^2`$ and $`\mathrm{tr}K^2`$, also provide conformal invariants. This fact can be phrased in an alternative way. We note that both $`=K^2`$ and $`=\mathrm{tr}K^2`$, have $`^a{}_{a}{}^{}K`$ so that Eq.(49) is satisfied. |
warning/0507/hep-ph0507138.html | ar5iv | text | # Black Hole Remnants at the LHC
## 1 Introduction
High energetic particle collisions will eventually lead to strong gravitational interactions and result in the formation of a black hole’s horizon. In the presence of large additional compactified dimensions , it could be possible that the threshold for black hole production lies within the accessible range for future experiments (e.g. LHC, CLIC). In the context of models with such large extra dimensions, black hole production is predicted to drastically change high energy physics already at the LHC. These effective models with extra dimensions are string-inspired extensions to the Standard Model in the overlap region between ’top-down’ and ’bottom-up’ approaches.
The possible production of TeV-scale black holes at the LHC is surely one of the most exciting predictions of physics beyond the Standard Model and has received a great amount of interest during the last years . For reviews on the subject the interested reader is referred to .
Due to their Hawking-radiation , these small black holes will have a temperature of some $`100`$ GeV and will decay quickly into $`1025`$ thermally distributed particles of the Standard Model (before fragmentation of the emitted partons). This yields a signature unlike all other new predicted effects. The black hole’s evaporation process connects quantum gravity with quantum field theory and particle physics, and is a promising way towards the understanding of Planck scale physics.
Thus, black holes are a fascinating field of research which features an interplay between General Relativity, thermodynamics, quantum field theory, and particle physics. The investigation of black holes would allow to test Planck scale effects and the onset of quantum gravity. Therefore, the understanding of the black holes properties is a key knowledge to the phenomenology of physics beyond the Standard Model.
Recently, the production of black holes has been incorporated into detailed numerical simulations for black hole production and decay in ultra-high energetic hadron-hadron interactions .
So far the numerical simulation has assumed that the black hole decays in its final phase completely into some few particles of the Standard Model. However, from the theoretical point of view, there are strong indications that the black hole does not evaporate completely, but leaves a stable black hole remnant. In this work, we will include this possibility into the numerical simulation and examine the consequences for the observables of the black hole event. These investigations might allow to reconstruct initial parameters of the model from observed data and can shed light onto these important questions.
The aim of this investigation is not to derive the formation of a remnant from a theory of modified gravity but to incorporate the assumption of such a remnant into the possible signatures black hole events in high energetic particle interactions.
This paper is organized as follows: the next section briefly reviews basic facts about black holes in extra dimensions. Section 3 discusses the issue of black hole remnants and introduces a useful parametrization for the thermodynamical treatment. In section 4, we discuss the results of the numerical simulation. We conclude in section 5.
Throughout this paper we adopt the convention $`\mathrm{}=c=k_B=1`$.
## 2 Black Holes in Extra Dimensions
Arkani-Hamed, Dimopoulos and Dvali proposed a solution to the hierarchy problem by the introduction of $`d`$ additional compactified space-like dimensions in which only gravitons can propagate. The Standard Model (SM) particles are bound to our 4-dimensional sub-manifold, called our 3-brane.
Gauss’ law then relates the fundamental mass scale of the extended theory, $`M_\mathrm{f}`$, to the apparent Planck scale, $`m_\mathrm{p}10^{16}`$ TeV, by the volume of the extra dimensions. In the case of toroidal compactification on radii of equal size this yields
$`m_\mathrm{p}^2=M_\mathrm{f}^{d+2}R^d.`$ (1)
Thus, for large radii, the Planck scale can be lowered to a new fundamental scale, $`M_\mathrm{f}`$ which can lie close by the electroweak scale.
The radius $`R`$ of the extra dimensions is then in the range mm to $`10^3`$ fm for $`d`$ from $`2`$ to $`7`$, or the inverse radius $`1/R`$ lies in energy range eV to MeV, resp. Since this radii are large compared to the Planck scale, this setting is called a scenario with large extra dimensions (LXDs). For recent constraints on the parameter of the model see e.g. .
Using the higher dimensional Schwarzschild-metric , it can be derived that the horizon radius $`R_H`$ of a black hole is substantially increased in the presence of LXDs , reflecting the fact that gravity at small distances becomes stronger. For a black hole of mass $`M`$ one finds
$`R_H^{d+1}={\displaystyle \frac{1}{d+1}}{\displaystyle \frac{1}{M_\mathrm{f}^{d+1}}}{\displaystyle \frac{M}{M_\mathrm{f}}}.`$ (2)
The horizon radius for a black hole with mass $``$ TeV is then $`10^3`$ fm, and thus $`R_HR`$ for black holes which can possibly be produced at colliders or in ultra high energetic cosmic rays. Also in higher dimensions the entropy, $`S`$, of the black hole is proportional to its horizon surface which is given by
$`𝒜_{(d+3)}`$ $`=`$ $`\mathrm{\Omega }_{(d+3)}R_H^{d+2}`$ (3)
where $`\mathrm{\Omega }_{(d+3)}`$ is the surface of the $`d+3`$-dimensional unit sphere
$`\mathrm{\Omega }_{(d+3)}={\displaystyle \frac{2\pi ^{\frac{d+3}{2}}}{\mathrm{\Gamma }(\frac{d+3}{2})}}.`$ (4)
Black holes with masses in the range of the lowered Planck scale should be a subject of quantum gravity. Since there is yet no theory available to perform these calculations, the black holes are treated as semi classical objects which form intermediate meta-stable states. Thus, the black holes are produced and decay according to the semi classical formalism of black hole physics.
To compute the production probability, the cross-section of the black holes can be approximated by the classical geometric cross-section
$`\sigma (M)\pi R_H^2\mathrm{\Theta }(MM_{\mathrm{min}}),`$ (5)
an expression which contains only the fundamental Planck scale as coupling constant. $`M_{\mathrm{min}}`$ is the threshold above which the production can occur and expected to be a few $`\times M_\mathrm{f}`$. As has been shown recently , such a threshold arises naturally in certain types of higher order curvature gravity.
The semi classical black hole cross section has been under debate , but further investigations justify the use of the classical limit at least up to energies of $`10M_\mathrm{f}`$ . It has further been shown that the naively expected classical result remains valid also in string-theory . However, this interesting topic is still a matter of ongoing research, see e.g. the very recent contributions in Refs. .
A common approach to improve the naive picture of colliding point particles, is to treat the creation of the horizon as a collision of two shock fronts in an Aichelburg-Sexl geometry describing the fast moving particles . Due to the high velocity of the moving particles, space time before and after the shocks is almost flat and the geometry can be examined for the occurrence of trapped surfaces.
These semi classical considerations do also give rise to form factors which take into account that not the whole initial energy is captured behind the horizon. These factors have been calculated in and depend on the number of extra dimensions, however their numerical values are of order one. Setting $`M_\mathrm{f}1`$TeV and $`d=2`$ one finds $`\sigma 1`$ TeV$`{}_{}{}^{2}400`$ pb. With this cross section it is further found that these black holes will be produced at the LHC in huge numbers on the order of $`10^9`$ per year .
Once produced, the black holes will undergo an evaporation process whose thermal properties carry information about the parameters $`M_\mathrm{f}`$ and $`d`$. An analysis of the evaporation will therefore offer the possibility to extract knowledge about the topology of our space time and the underlying theory.
The evaporation process can be categorized in three characteristic stages :
1. Balding phase: In this phase the black hole radiates away the multi-pole moments it has inherited from the initial configuration, and settles down in a hairless state. During this stage, a certain fraction of the initial mass will be lost in gravitational radiation.
2. Evaporation phase: The evaporation phase starts with a spin down phase in which the Hawking radiation carries away the angular momentum, after which it proceeds with the emission of thermally distributed quanta until the black hole reaches Planck mass. The radiation spectrum contains all Standard Model particles, which are emitted on our brane, as well as gravitons, which are also emitted into the extra dimensions. It is expected that most of the initial energy is emitted during this phase in Standard Model particles .
3. Planck phase: Once the black hole has reached a mass close to the Planck mass, it falls into the regime of quantum gravity and predictions become increasingly difficult. It is generally assumed that the black hole will either completely decay in some last few Standard Model particles or a stable remnant will be left, which carries away the remaining energy.
The evaporation phase is expected to be the most important phase for high energy collisions. The characteristics of the black hole’s evaporation in this phase can be computed using the laws of black hole thermodynamics and are obtained by first solving the field equations for the metric of the black hole, then deriving the surface gravity, $`\kappa `$, from which the temperature of the black hole follows via
$`T={\displaystyle \frac{\kappa }{2\pi }}.`$ (6)
By identifying the total energy of the system with the mass of the black hole one then finds the entropy $`S`$ by integrating the thermodynamical identity
$`{\displaystyle \frac{S}{M}}={\displaystyle \frac{1}{T}},`$ (7)
which fixes the constant factor relating the entropy to the horizon surface. A possible additive constant is generally chosen such that the entropy is zero for vanishing horizon surface. In contrast to classical thermodynamical objects, this does not necessarily imply that the entropy vanishes at zero temperature, see e.g. .
By now, several experimental groups include black holes into their search for physics beyond the Standard Model. For detailed studies of the experimental signatures, PYTHIA 6.2 has been coupled to CHARYBDIS creating an event generator allowing for the simulation of black hole events and data reconstruction from the decay products. Previous analysis within this framework are summarized in Refs. . Ideally, the energy distribution of the decay products allows a determination of the temperature (by fitting the energy spectrum to the predicted shape) as well as of the total mass of the object (by summing up all energies). This then allows to reconstruct the scale $`M_\mathrm{f}`$ and the number of extra dimensions.
These analysis however, have so far omitted the possibility of a stable black hole remnant but assume instead a final decay into some few particles, whose number is treated as a free parameter ranging from $`25`$.
In the following we will examine the possibility that a stable black hole remnant of about Planck mass is left with use of the PYTHIA 6.2/CHARYBDIS event generator package.
## 3 Black Hole Remnants
The final fate of black holes is an unresolved subject of ongoing research. The last stages of the evaporation process are closely connected to the information loss puzzle. The black hole emits thermal radiation, whose sole property is the temperature, regardless of the initial state of the collapsing matter. So, if a black hole completely decays into statistically distributed particles, unitarity can be violated. This happens when the initial state is a pure quantum state and then evolves into a mixed state .
When one tries to avoid the information loss problem two possibilities are left. The information is regained by some unknown mechanism or a stable black hole remnant is formed which keeps the information. Besides the fact that it is unclear in which way the information should escape the horizon there are several other arguments for black hole remnants :
* The uncertainty relation: The Schwarzschild radius of a black hole with Planck mass is of the order of the Planck length. Since the Planck length is the wavelength corresponding to a particle of Planck mass, a problem arises when the mass of the black hole drops below Planck mass. Then one has trapped a mass inside a volume which is smaller than allowed by the uncertainty principle . To avoid this problem, Zel’dovich has proposed that black holes with masses below Planck mass should be associated with stable elementary particles . Also, the occurrence of black hole remnants within the framework of a generalized uncertainty principle has been investigated in .
* Corrections to the Lagrangian: The introduction of additional terms, which are quadratic in the curvature, yields a decrease of the evaporation temperature towards zero . This holds also for extra dimensional scenarios and is supported by calculations in the low energy limit of string theory . The production of TeV-scale black holes in the presence of Lovelock higher-curvature terms has been examined in and it was found that these black holes can become thermodynamically stable since their evaporation takes an infinite amount of time.
* Further reasons for the existence of remnants have been suggested to be black holes with axionic charge , the modification of the Hawking temperature due to quantum hair or magnetic monopoles . Coupling of a dilaton field to gravity also yields remnants, with detailed features depending on the dimension of space-time .
* One might also see the arising necessity for remnant formation by applying the geometrical analogy to the black hole and quantize the radiation into wavelengths that fit on the surface, i.e. the horizon . The smaller the size of the black hole, the smaller the largest possible wavelength and the larger the smallest possible energy quantum that can be emitted. Should the energy of the lowest energy level already exceed the total mass of the black hole, then no further emission is possible. Not surprisingly, this equality happens close to the Planck scale and results in the formation of a stable remnant.
Of course these remnants, which in various context have also been named Maximons, Friedmons, Cornucopions, Planckons or Informons, are not a miraculous remedy but bring some new problems along. Such is e.g. the necessity for an infinite number of states which allows the unbounded information content inherited from the initial state.
## 4 Signatures of Black Hole Remnants
We now attempt to construct a numerically applicable model for modifications of the black hole’s temperature in order to simulate the formation of a black hole remnant. Though the proposals of remnant formation in the literature are build on various different theoretical approaches, they have in common that the temperature of the black hole drops to zero already at a finite black hole mass. We will denote the mass associated with this finite remnant size with $`M_\mathrm{R}`$ and make the reasonable identification $`M_\mathrm{R}=M_{\mathrm{min}}`$. Instead of deriving such a minimal mass within the frame of a specific model, we aim in this work to parametrize its consequences for high energy collisions.
For our purposes, we will assume that we are dealing with a theory of modified gravity which results in a remnant mass and parametrize the deviations of the entropy $`S(M)`$. This entropy now might differ from the Hawking-entropy by correction terms in $`M_\mathrm{R}/M`$. For black hole masses $`M`$ much larger than $`M_\mathrm{R}`$ we require to reproduce the standard result. The expansion then reads
$$S(M)=𝒜_{(d+3)}M_f^{d+2}\left[a_0+a_1\left(\frac{M_\mathrm{R}}{M}\right)+a_2\left(\frac{M_\mathrm{R}}{M}\right)^2+\mathrm{}\right]$$
(8)
with dimensionless coefficients $`a_i`$ depending on the specific model (see e.g. ). As defined in Eq. (3), $`𝒜`$ is the surface of the black hole and a function of $`M`$. For the standard scenario one has
$$a_0=\frac{d+1}{d+2}\frac{2\pi }{\mathrm{\Omega }_{(d+3)}},a_{i>1}=0.$$
(9)
Note that in general
$$S_0=S(M=M_\mathrm{R})$$
(10)
will differ from the unmodified black hole entropy since the Schwarzschild-radius can be modified.
It should be understood that an underlying theory of modified gravity will allow to compute $`M_\mathrm{R}=M_\mathrm{R}(a_i)`$ explicitly from the initially present parameters. This specific form of these relations however, depends on the ansatz. We will instead treat $`M_\mathrm{R}`$ as the most important input parameter. Though the coefficients $`a_i`$ in principle modify the properties of the black hole’s evaporation, the dominating influence will come from the existence of a remnant mass itself, making the $`a_i`$ hard to extract from the observables.
To make this point clear, let us have a closer look at the evaporation rate of the black hole by assuming a remnant mass. Note, that the Hawking-evaporation law can not be applied towards masses that are comparable to the energy of the black hole because the emission of the particle will have a non-negligible back reaction. In this case, the black hole can no longer be treated in the micro canonical ensemble but instead, the emitted particles have to be added to the system, allowing for a loss of energy into the surrounding of the black hole. Otherwise, an application of the Hawking-evaporation down to small masses comparable to the temperature of the black hole, would yield the unphysical result that the evaporation rate diverges because one has neglected that the emitted quanta lower the mass of the black hole.
This problem can be appropriately addressed by including the back reaction of the emitted quanta as has been derived in . It is found that in the regime of interest here, when $`M`$ is of order $`M_\mathrm{f}`$, the emission rate for a single particle micro state is modified and given by the change of the black hole’s entropy
$$n(\omega )=\frac{\mathrm{exp}[S(M\omega )]}{\mathrm{exp}[S(M)]}.$$
(11)
If the average energy of the emitted particles is much smaller than $`M`$, as will be the case for $`MM_\mathrm{f}`$, one can make the approximation
$`S(M)S(M\omega ){\displaystyle \frac{S}{M}}\omega ={\displaystyle \frac{\omega }{T}}`$ (12)
which, inserted in Eq.(11) reproduces the familiar relation. The single particle distribution can be understood by interpreting the occupation of states as arising from a tunnelling probability . From the single particle number density (Eq. 11) we obtain the average particle density by counting the multi particle states according to their statistics
$$n(\omega )=\left(\mathrm{exp}[S(M)S(M\omega )]+s\right)^1,$$
(13)
where
$`s`$ $`=`$ $`1\text{for Fermi-Dirac statistic}`$
$`s`$ $`=`$ $`0\text{for Boltzmann statistic}`$
$`s`$ $`=`$ $`1\text{for Bose-Einstein statistic},`$ (14)
and $`\omega MM_\mathrm{R}`$, such that nothing can be emitted that lowers the energy below the remnant mass. Note, that this number density will assure that the remnant is formed even if the time variation of the black hole’s temperature (or its mass respectively) is not taken into account.
For the spectral energy density we then use this particle spectrum and integrate over the momentum space. Since we are concerned with particles of the Standard Model which are bound to the 3-brane, their momentum space is the usual 3-dimensional one. This yields
$`\epsilon ={\displaystyle \frac{\mathrm{\Omega }_{(3)}}{(2\pi )^3}}\zeta (4){\displaystyle _0^{MM_\mathrm{R}}}{\displaystyle \frac{\omega ^3\mathrm{d}\omega }{\mathrm{exp}[S(M)S(M\omega )]+s}}.`$ (15)
From this, we obtain the evaporation rate with the Stefan-Boltzmann law to
$`{\displaystyle \frac{\mathrm{d}M}{\mathrm{d}t}}={\displaystyle \frac{\mathrm{\Omega }_{(3)}^2}{(2\pi )^3}}R_H^2\zeta (4){\displaystyle _0^{MM_\mathrm{R}}}{\displaystyle \frac{\omega ^3\mathrm{d}\omega }{\mathrm{exp}[S(M)S(M\omega )]+s}}.`$ (16)
Since we are dealing with emitted particles bound to the brane, the surface through which the flux disperses is the $`2`$-dimensional intersection of the black hole’s horizon with the brane.
Inserting the modified entropy Eq. (8) into the derived expression Eq. (16), one sees that the evaporation rate depends not only on $`M_\mathrm{R}`$ but in addition on the free parameters $`a_i`$. However, for large $`M`$ the standard scenario is reproduced and we can apply the canonical ensemble. E.g. for the Fermi-Dirac statistic one obtains
$`{\displaystyle \frac{\mathrm{d}M}{\mathrm{d}t}}={\displaystyle \frac{\mathrm{\Omega }_{(3)}^2}{(2\pi )^3}}R_H^2\zeta (4)\mathrm{\Gamma }(4)T^4\text{for}MM_\mathrm{R}.`$ (17)
Whereas for $`M/M_\mathrm{R}1`$, the dominant contribution from the integrand in Eq. (16) comes from the factor $`\omega ^3`$ and the evaporation rate will increase with a power law. The slope of this increase will depend on $`S_0`$. From this qualitative analysis, we can already conclude that the coefficients $`a_i`$ will influence the black hole’s evaporation only in the intermediate mass range noticeably. If we assume the coefficients to be in a reasonable range – i.e. each $`a_i`$ is of order $`1`$ or less and the coefficient $`a_{i+1}`$ is smaller<sup>1</sup><sup>1</sup>1From naturalness, one would expect the coefficients to become smaller with increasing $`i`$ by at least one order of magnitude see e.g. . than the coefficient $`a_i`$ and the series breaks off at a finite $`i`$ – then the deviations from the standard evaporation are negligible as is demonstrated in Figs. 1, 2 and 3.
Figure 1 shows the evaporation rate Eq. (16) for various $`d`$ with the standard parameters (9). Figure 2 and 3 show various choices of parameters for $`d=3`$ and $`d=5`$ as examples. Note, that setting $`a_3`$ to $`1`$ is already in a very extreme range since a natural value was several orders of magnitude smaller: $`a_310^3`$ (in this case the deviations would not be visible in the plot). For our further numerical treatment, we have included the possibility to vary the $`a_i`$ but one might already at this point expect them not to have any influence on the characteristics of the black hole’s evaporation except for a slight change in the temperature-mass relation.
From the evaporation rate Eq.(16) one obtains by integration the mass evolution $`M(t)`$ of the black hole. This is shown for the continuous mass case in Figure 4. For a realistic scenario one has to take into account that the mass loss will proceed by steps by radiation into the various particles of the Standard Model.
## Results
We have included the evaporation rate, parametrized according to the previous section, into the black hole event generator CHARYBDIS and examined the occurring observables within the PYTHIA environment. Since these black hole remnants are stable, they are of special interest as they are available for close investigations. Especially those remnants carrying an electric charge offer exciting possibilities as investigated in .
It has also been shown in that no naked singularities have to be expected for reasonably charged black holes and that the modification of the Hawking radiation due to the electric charge can be neglected for the parameter ranges one expects at the LHC. This means in particular that the interaction of emitted charged particles with the black hole does not noticeably modify the emission probability. Although there might be uncertainties in the low energy limit where QED or QCD interactions might have unknown consequences for the processes at the horizon.
The formation of a remnant indeed solves a (technical) problem occuring within the treatment of a final decay: it might in principle have happened that during its evaporation process, the black hole has emitted mostly electrically charged particles and ended up with an electric charge of order ten. In such a state, it would then be impossible for the black hole to decay into less than ten particles of the SM, whereas the standard implementation allows only a decay into a maximum of 5 particles.
Therefore, in the original numerical treatment, the process of Hawking radiation has before been assumed to minimize the charge of the evaporating hole in each emission step. In such a way, it was assured that the object always had a small enough charge to enable the final decay in $`5`$ particles without any violation of conservation laws. This situation changes if the remnant is allowed to keep the electric charge. In the here presented analysis, the assumption of charge minimization has therefore been dropped as it is no longer necessary. However, we want to stress, that the in- or exclusion of charge minimization does not modify the observables investigated<sup>2</sup><sup>2</sup>2The differences in the finally observable charged particle distributions from the black hole decay are changed by less than 5% compared to the charge minimization setting..
When attempting to investigate slowly decaying objects, one might be concerned whether these decay in the collision region or might be able to leave the detector, thereby still emitting radiation. As shown for the continuous case in Fig. 4, the average energy of the emitted particles drops below an observable range within a $`10`$ fm radius. Even if one takes into account the large $`\gamma `$-factor, the black hole will have shrunken to remnant-mass safely in the detector region. This is shown for a sample of simulated events in Fig. 4 (symbols) which displays the mass evolution of these collider produced black holes. Here, we estimated the time, $`t`$, for the stochastic emission of a quanta of energy $`E`$ to be $`1/E`$. This numerical result agrees very well with the expectations from the continuous case.
To understand the fast convergence of the black hole mass, recall the spectral energy density which enters in Eq.(16) and which dictates the distribution of the emitted particles. Even though the spectrum is no longer an exactly Planckian, it still retains a maximum at energies $`1/T`$. If the black hole’s mass decreases, the emission of the high energetic end of the spectrum is no longer possible. For masses close to the Planck scale, the spectrum has a maximum at the largest possible energies that can be emitted. Thus, the black hole has a high probability to emit its remaining energy in the next emission process. However, theoretically, the equilibrium time goes to infinity (because the evaporation rate falls to zero, see Fig. 1) and the black hole will emit an arbitrary amount of very soft photons. For practical purposes, we cut off<sup>3</sup><sup>3</sup>3The emission of objects carrying color charge is disabled after the maximally possible energy drops below the mass of the lightest meson, i.e. the pion. the evaporation when the black hole reached the mass $`M_\mathrm{R}+0.1`$ GeV.
Figure 5 shows the rapidity of the produced black hole remnants in a proton-proton collision at $`\sqrt{s}=14`$ TeV. All plots are for $`d=2`$ since a higher number of extra dimensions leads to variations of less than 5%. The reader should be aware that the present numerical studies assume the production of one black hole in every event. To obtain the absolute cross sections the calculated yields have to be multiplied by the black hole production cross section $`\sigma (\mathrm{pp}BH)`$. Due to the uncertainties in the absolute production cross section of black holes we have taken this factor explicitely out. For the present examination we have initialized a sample of 50,000 events. The black hole remnants are strongly peaked around central rapidities, making them potentially accessible to the CMS and ATLAS experiments. In Figure 6 we show the distribution of the produced black hole remnants as a function of the transverse momentum.
Figure 7 shows the transverse momentum, $`p_T`$, of the decay products as it results from the modified multi particle number density Eq. (13) before fragmentation. Figure 8 shows the $`p_T`$-spectrum after fragmentation. In both cases, one clearly sees the additional contribution from the final decay which causes a bump in the spectrum which is absent in the case of a remnant formation. After fragmentation, this bump is slightly washed out but still present. However, from the rapidity distribution and the fact that the black hole event is spherical, a part of the high $`p_T`$-particles will be at large $`y`$ and thus be not available in the detector. We therefore want to mention that one has to include the experimental acceptance in detail if one wants to compare to experimental observables.
Figure 9 shows the total multiplicities of the event. When a black hole remnant is formed, the multiplicity is increased due to the additional low energetic particles that are emitted in the late stages instead of a final decay with $`25`$ particles. Note that this multiplicity increase is not an effect of the remnant formation itself, but stems from the treatment of the decay in the micro-canonical ensemble used in the present calculation. I.e. the black hole evaporates a larger amount of particles with lower average energy.
Figure 10 shows the sum over the transverse momenta of the black holes’ decay products. To interpret this observable one might think of the black hole event as a multi-jet with total $`\mathrm{\Sigma }p_T`$. As is evident, the formation of a remnant lowers the total $`\mathrm{\Sigma }p_T`$ by about $`M_\mathrm{R}`$. This also means, that the signatures of the black hole as previously analyzed are dominated by the daubtful final decay and not by the Hawking phase. It is interesting to note that the dependence on $`M_\mathrm{f}`$ is dominated by the dependence on $`M_\mathrm{R}`$, making the remnant mass the primary observable, leading to an increase in the missing energy.
## 5 Conclusion
We have parametrized the modifications to the black hole evaporation arising from the presence of a remnant mass. The modified spectral density is included in the numerical simulation for black hole events. To give a specific example, we have examined the formation of black hole remnants in proton proton collision at $`\sqrt{s}=14`$ TeV and set it in contrast to a final decay of the black hole.
We predict a significant decrease of the total transverse momentum of the black hole remnant events due to the absence of the final decay particles. Even more, the multiplicity of the event is increased by a factor $`3/2`$ arising from the micro-canonical treatment of the evaporation process.
The formation of the black hole remnant results in a strong modification of most predicted black hole signatures. However, remnant formation itself leads to prominent experimental signatures (see e.g. $`\mathrm{\Sigma }p_T`$). This makes the search for black hole remnants promising and experimentally accessible for the CMS and ATLAS experiments.
## Acknowledgements
We thank Horst Stöcker for helpful discussions. This work was supported by NSF PHY/0301998 and DFG. SH thanks the FIAS for kind hospitality. |
warning/0507/physics0507057.html | ar5iv | text | # Numerical integration of the discrete-ordinate radiative transfer equation in strongly non homogeneous media.
## I Introduction
The mathematical modeling of radiative transfer in which the phenomena of absorption, emission and scattering are taken into account is usually made using the linearized Boltzmann equation, also known as radiation transfer equation. This equation describes the transfer of radiation, with a given wavelength, through a medium with certain absorbing and scattering properties. A particular application of this equation is the study of radiation transfer in the atmosphere, where the medium properties vary strongly with height. Getting accurate solutions of this problem is important for evaluating energy balance on planetary atmospheres as well as for the so called ”inverse problem” where the boundary conditions (in particular the ground albedo properties) are deduced from measurements of radiation and knowledge of the medium characteristics.
In this work we consider the time-independent, monochromatic radiative transfer equation using the well-tested and widely used discrete-ordinate method of Stamnes et al. Stamnes and the plane parallel approach where the optical properties depend only on the vertical coordinate $`z`$. The procedure requires the solution of a system of $`n`$ coupled linear ordinary differential equations (one for each stream or discrete-ordinate component of the intensity). This set of equations is subject to a two-point boundary condition at the top and bottom of the medium. In the general case no analytic solutions exist for this problem since the medium optical properties (phase function, absorption and scattering coefficients) depend on the position $`z`$ in the vertically inhomogeneous medium.
To obtain a formal solution, the medium is generally assumed to be layered with piecewise constant optical properties, i.e. it is divided into $`N`$ adjacent homogeneous layers where the eigenvalues are computed. The coefficients of the solution are determined by imposing the continuity condition at the boundaries between adjacent layers and the two-point boundary conditions at the top and bottom of the medium Liou ; Kylling .
Unfortunately in many cases the optical properties of the medium are not homogeneous, in fact they show strong variations along the vertical axis and have therefore different characteristic length scales. In these cases, as we will see below, the homogeneous layers assumption may lead to errors of $`10\%`$ in the estimation of the scattered radiation. Furthermore in most practical applications the medium characteristics are only known, based on measurements or calculations, at a discrete number of points Carmen&Ana . It is therefore required to implement a method that is able to cope both with the strong variations on the medium properties and with the discrete character of the information available.
Here we solve the discrete-ordinate approximation to the radiative transfer equation with a different approach. Since the equations and the boundary conditions are linear in the intensity, the two-point boundary problem can be solved with a shooting method. The problem is then reduced to the solution of an initial value problem which can be solved numerically. We propose a numerical integration of the initial value problem that applies an adaptive step method (such as step doubling), interpolates the optical properties from the discrete set of available data (for example using a cubic spline) and uses a weighted evaluation of the integrand along the interval (in our case a 5th order Runge-Kutta scheme). This allows the adaptation of the numerical integration to the local characteristic length-scales of the problem under consideration. The assumption of homogeneity is thus not required.
As an example we apply this method to solve the two-stream discrete ordinate version of the, non-emitting, radiative transfer equation in one spatial dimension (these results can be equivalently extended to multi-stream and emitting versions of this equation). We present two benchmark cases with known analytical solutions and compare the performance obtained when the step doubling interpolating method is used and when the piecewise homogeneity is imposed. We will show that using the same available discrete information, our method can significantly improve the accuracy of the solution. In section II we will first introduce the radiative transfer equation and the equations for the direct and diffuse intensity components in the two-stream discrete ordinate approach. In section III we will describe two benchmark cases characterized by linear and exponential height dependence of the optical coefficients in the equations. In section IV we will explain our numerical procedure and will validate our results and those obtained under the assumption of piecewise homogeneous layers against the exact analytical solutions of the benchmark problems. Finally conclusions are presented in section V.
## II Multiple Scattering and the Radiative Transfer Equation
Our aim is to solve the radiative transfer problem in the plane parallel approach. As represented in Figure 1 this corresponds to the case in which optical properties only depend on altitude $`z`$ in a plane parallel geometry. Following for instance Stamnes the radiative transfer equation in this approach is written as,
$`\mathrm{cos}\theta {\displaystyle \frac{dI(z,\theta ,\varphi ,\lambda )}{dz}}=\beta _T(z,\lambda )I(z,\theta ,\varphi ,\lambda )+{\displaystyle \frac{\beta _{sca}(z,\lambda )}{4\pi }}{\displaystyle _0^{2\pi }}𝑑\varphi ^{}`$ (1)
$`{\displaystyle _0^\pi }\mathrm{sin}\theta ^{}d\theta ^{}p(z,\theta ,\varphi ,z^{},\theta ^{},\varphi ^{},\lambda )I(z,\theta ^{},\varphi ^{})`$
where $`\theta `$ and $`\varphi `$ are respectively the polar and azimuthal angles. Here $`\beta _{sca}`$ represents the attenuation due to scattering effects and $`\beta _T(z,\lambda )`$ is the extinction coefficient defined as $`\beta _T(z,\lambda )=\beta _{sca}(z,\lambda )+_i\beta _i(z,\lambda )`$ where the summation in $`i`$ extents to all the molecular components considered. Here $`\beta _i(z,\lambda )=n_i(z)\sigma _i(\lambda )`$ is the absorption coefficient of a given component, $`\sigma _i(\lambda )`$ is the attenuation cross section due to absorption and $`n_i(z)`$ the atmospheric number density for species $`i`$ (which generally depends exponentially on the height $`z`$). In Eq. (1), $`p(z,\theta ,\varphi ,z^{},\theta ^{},\varphi ^{},\lambda )`$ is the phase function of the scattering particles which is normalized as follows: $`\frac{1}{4\pi }_0^{2\pi }𝑑\varphi _0^\pi p(\theta ,\varphi ,z^{},\theta ^{},\varphi ^{},\lambda )\mathrm{sin}(\theta )𝑑\theta =1`$. This function gives the probability for a photon of wavelength $`\lambda `$, incident on the scattering particle with angles $`(\theta ^{},\varphi ^{})`$ to be scattered in the direction $`(\theta ,\varphi )`$ and satisfies $`p(\theta ,\varphi ,\theta ^{},\varphi ^{},\lambda )=p(\mathrm{cos}\mathrm{\Theta },\lambda )`$ where $`\mathrm{\Theta }`$ is the scattering angle, related to the polar and azimuthal angles by $`\mathrm{cos}\mathrm{\Theta }=\mathrm{cos}\theta ^{}\mathrm{cos}\theta +\mathrm{sin}\theta ^{}\mathrm{sin}\theta \mathrm{cos}(\varphi \varphi ^{})`$.
The solution of Eq. (1) is splitted into two terms
$`I(z,\theta ,\varphi ,\lambda )=I^{dir}(z,\lambda )+I^{dif}(z,\mu ,\varphi ,\lambda ).`$ (2)
$`I^{dir}(z,\lambda )`$, the direct intensity, is the solution of Eq. (1) when there is no multiple scattering (no integral term):
$$\mu _0\frac{dI^{dir}(z)}{dz}=\beta _T(z)I^{dir}(z)$$
(3)
where $`\mu _0`$ is the cosine of the polar angle for the incident radiation. The diffuse intensity is the solution of
$`\mu {\displaystyle \frac{dI^{dif}(z,\mu ,\varphi ,\lambda )}{dz}}=\beta _T(z,\lambda )I^{dif}(z,\mu ,\varphi ,\lambda )+{\displaystyle \frac{\beta _{sca}(z,\lambda )}{4\pi }}{\displaystyle _0^{2\pi }}𝑑\varphi ^{}`$ (4)
$`{\displaystyle _0^\pi }𝑑\mu ^{}p(z,\mu ,\varphi ,z^{},\mu ^{},\varphi ^{},\lambda )I^{dif}(z,\mu ^{},\varphi ^{})`$
$`{\displaystyle \frac{\beta _{sca}(z,\lambda )I^{dir}(z,\lambda )}{4\pi }}p(z,\mu ,\varphi ,z^{},|\mu _0|,\varphi _0,\lambda )^{},\lambda ).`$
with $`\mu =\mathrm{cos}\theta `$. In this equation the variable $`\mu `$ takes values in the range $`1<\mu <1`$. Negative $`\mu `$ corresponds to radiation going downwards, whereas positive $`\mu `$ describes radiation going upwards. In order to solve Eq. (4) the diffuse intensity is expanded in a $`2n`$ Fourier cosine series (from now on we drop the superscript $`dif`$): $`I(z,\mu ,\varphi ,\lambda )=_{m=0}^{2n1}I^m(z,\mu )\mathrm{cos}m(\varphi _0\varphi )`$. The phase function is expanded in a basis of $`2n`$ Legendre polynomials $`p(z,\mu ,\varphi ,z^{},\mu ^{},\varphi ^{},\lambda )=p(z,\mathrm{cos}\mathrm{\Theta },\lambda )=_{l=0}^{2n1}(2l+1)g_lP_l(\mathrm{cos}\mathrm{\Theta })`$. With these transformations and the theorem of addition of spherical harmonics Eq. (4) becomes a set of integro-differential equations depending only on the $`z`$ and $`\mu `$ coordinates:
$`\mu {\displaystyle \frac{dI^m(z,\mu )}{dz}}`$ $`=`$ $`\beta _T(z,\lambda )I^m(z,\mu )+J^m(z,\mu ),`$ (5)
where
$`J^m(z,\mu )`$ $`=`$ $`{\displaystyle \frac{\beta _{sca}(z,\lambda )}{2}}{\displaystyle \underset{l=0}{\overset{2n1}{}}}(2l+1)g_l^mP_l^m(\mu )`$
$`\left({\displaystyle _1^1}P_l^m(\mu ^{})I^m(z,\mu ^{})𝑑\mu ^{}+{\displaystyle \frac{I^{dir}(z,\lambda )}{2\pi }}(2\delta _{0,m})(1)^{(l+m)}P_l^m(|\mu _0|)\right)`$
being $`P_l^m(\mu )`$ the associated Legendre polynomial, $`g_l^m=g_l\frac{(lm)!}{(l+m)!}`$, and $`g_l=\frac{1}{2}_1^1p(\mathrm{cos}\mathrm{\Theta })P_l(\mathrm{cos}\mathrm{\Theta })d(\mathrm{cos}\mathrm{\Theta })`$.
In the discrete ordinate approximation the angular integral term in Eq. (II) is represented by a summation over $`n`$ Gaussian quadrature points $`\mu _s`$ (fixed angles) also called ”streams”. The intensity given by Eq. (5) must satisfy boundary conditions at the top ($`z=z_{max}`$) and bottom ($`z=0`$) of the medium. Therefore we end up with a system of $`n`$ coupled ordinary differential equations of the type (5), one for each stream $`\mu _s`$, and subject to a two-point boundary condition. Solving these equations, we will get a discrete approximation to the angular distribution of $`I^m(z,\mu )`$ from the top of the atmosphere to the surface level.
In this work, for simplicity, we will consider the two-stream approximation to this problem and obtain the two streams $`I^m(z,\mu _1)`$ (downwards) and $`I^m(z,\mu _2)`$ (upwards) fulfilling the boundary conditions that impose, first, no diffuse radiation incident at the top, and second, no radiation reflected back at the surface. This is expressed as follows,
$`I^m(z=z_{max},\mu _1)`$ $`=`$ $`0`$ (7)
$`I^m(z=0,\mu _2)`$ $`=`$ $`0.`$ (8)
In some cases it is preferred to use the optical depth of the medium $`\tau `$ instead of the vertical geometric distance $`z`$. This is a non-dimensional variable which is defined as $`\tau (z)=_z^{\mathrm{}}\beta _T(z)𝑑z`$. It describes the attenuation within the medium of an incident beam of radiation when emission and multiple scattering are ignored, i.e. it is a measurement of the direct component attenuation.
## III The benchmark problems
In this section we define three benchmark problems with an exact solution where the optical properties $`\beta _T(z,\lambda )`$ and $`\beta _{sca}(z,\lambda )`$ in Eq. (1) have different altitude ($`z`$) dependences. For the sake of simplicity we consider that particles scatter radiation uniformly, i.e. the scattering phase function is $`p(\mathrm{cos}\mathrm{\Theta })=1`$ and that $`z_{max}=1`$.
For the integral part in Eq. (II) and following the two-stream approximation we consider a weighted summation over the streams $`\mu _1=1/\sqrt{(}3)`$ (down) and $`\mu _2=1/\sqrt{(}3)`$ (up) which are the zeroes of the second order Legendre polynomial, i.e. $`P_2(\mu _i)=0`$, $`i=1,2`$. Then, Eq. (5) can be transformed in two coupled ordinary differential equations for the down ($`I_1`$) and up ($`I_2`$) radiation intensity:
$`\mu _1{\displaystyle \frac{dI_1}{dz}}`$ $`=`$ $`\beta _T(z)I_1+{\displaystyle \frac{0.5\beta _{sca}(z)}{2}}\left(I_1+I_2+{\displaystyle \frac{I^{dir}(z)}{2\pi }}\right),`$ (9)
$`\mu _2{\displaystyle \frac{dI_2}{dz}}`$ $`=`$ $`\beta _T(z)I_2+{\displaystyle \frac{0.5\beta _{sca}(z)}{2}}\left(I_1+I_2+{\displaystyle \frac{I^{dir}(z)}{2\pi }}\right).`$ (10)
After a change of variables $`z^{}=1z`$ (from now on dropping the prime in $`z`$) and taking into account that $`\mu _1=\mu _2`$ the above equations may be rewritten as,
$$\mu _0\frac{dI^{dir}}{dz}=\left[A(z)B(z)\right]\mu _2I^{dir}$$
(11)
$`{\displaystyle \frac{dI_1}{dz}}`$ $`=`$ $`A(z)I_1+B(z)I_2+I^{dir}(z){\displaystyle \frac{B(z)}{2\pi }}`$ (12)
$`{\displaystyle \frac{dI_2}{dz}}`$ $`=`$ $`A(z)I_2B(z)I_1I^{dir}(z){\displaystyle \frac{B(z)}{2\pi }}`$ (13)
where
$`A(z)`$ $`=`$ $`{\displaystyle \frac{\beta _T(z)+0.25\beta _{sca}(z)}{\mu _2}}`$ (14)
$`B(z)`$ $`=`$ $`{\displaystyle \frac{0.25\beta _{sca}(z)}{\mu _2}}`$ (15)
$`I^{dir}(z)`$ $`=`$ $`I_0e^{{\scriptscriptstyle }\frac{A(z)B(z)}{\mu _0}\mu _2dz}`$ (16)
and $`I_0`$ is fixed by the condition at the top $`I^{dir}(z=0)`$. In the following examples we will solve this problem with the two-point boundary conditions given in Eqs. (7)-(8), which are now explicitely written as $`I_1(z=0)=0`$ (at the top) and $`I_2(z=1)=0`$ (at the bottom). The incident intensity on the top is characterized by $`I_0=100`$ and $`\mu _0=0.788`$. We will consider two different height dependences of the optical properties.
### III.1 Linear dependence
We propose as study-case a linear dependence on the extinction coefficients. In particular we have considered the functions $`A(z)`$ and $`B(z)`$,
$`A(z)`$ $`=`$ $`{\displaystyle \frac{z(4a^2+c^2)}{2c}}`$ (17)
$`B(z)`$ $`=`$ $`{\displaystyle \frac{z\left(4a^2c^2\right)}{2c}}.`$ (18)
With this choice the general solutions of Eqs. (12)-(13) for the diffuse components are,
$$I_1(z)=e^{(z^2a)}C_1+e^{(z^2a)}C_2\frac{I_0(\mu _0^2\mu _0\mu _2)(4a^2c^2)}{16\pi (c^2\mu _{2}^{}{}_{}{}^{2}+4a^2\mu _{0}^{}{}_{}{}^{2})}e^{\left(\frac{cz^2\mu _2}{2\mu _0}\right)}$$
(19)
for the downwards radiation, and
$$I_2(z)=\frac{2a+c}{2ac}e^{(z^2a)}C_\mathit{1}\frac{2ac}{2a+c}e^{(z^2a)}C_\mathit{2}\frac{I_0(\mu _{\mathit{0}}^{}{}_{}{}^{2}+\mu _\mathit{0}\mu _\mathit{2})(4a^2c^2)}{16\pi (c^2\mu _{\mathit{2}}^{}{}_{}{}^{2}+4a^2\mu _{\mathit{0}}^{}{}_{}{}^{2})}e^{\left(\frac{cz^2\mu _\mathit{2}}{2\mu _\mathit{0}}\right)}$$
(20)
for the upwards component. Here $`C_1`$ and $`C_2`$ are arbitrary constants fixed with the boundary conditions in Eqs. (7)-(8). The direct component of the intensity is now
$$I^{dir}(z)=I_0e^{\left(\frac{cz^2\mu _2}{2\mu _\mathit{0}}\right)}.$$
(21)
For the specific choice $`c=10.5,a=4`$ the constants $`C_1`$ and $`C_2`$ that satisfy the boundary conditions are $`C_1=33.0927`$ and $`C_2=6.71346e05`$. For this particular set of parameters the total optical path is $`\tau 3`$. The optical functions $`\beta _T(z)`$ and $`\beta _{sca}(z)`$ are shown in Fig. 3-I and in Fig. 3-II the attenuation of the direct intensity as a function of height. The two-stream components of the diffused intensity are shown in Fig. 3-I.
### III.2 Exponential dependence
In atmospheric layers the gas density typically depends on the altitude as a growing exponential from $`z=0`$ (the top) to $`z=1`$ (the bottom). For this reason the choice of exponential dependence in the extinction coefficients is very convenient as it allows us to validate our method with solutions similar to those appearing in atmospheres, where our method has been applied Carmen&Ana . In particular we have chosen,
$`A(z)`$ $`=`$ $`{\displaystyle \frac{e^{(az)}(a^2b^2+c^2)}{2c}}`$ (22)
$`B(z)`$ $`=`$ $`{\displaystyle \frac{e^{(az)}(a^2b^2c^2)}{2c}}`$ (23)
The general solution for the downwards intensity is,
$$I_1(z)=e^{(e^{(az)}b)}C_\mathit{2}+e^{(e^{(az)}b)}C_\mathit{1}\frac{(c^2+a^2b^2)(\mu _{\mathit{0}}^{}{}_{}{}^{2}\mu _0\mu _2)}{16\pi (c^2\mu _{\mathit{2}}^{}{}_{}{}^{2}+a^2b^2\mu _{\mathit{0}}^{}{}_{}{}^{2})}I_0e^{\left(\frac{c\mu _\mathit{2}e^{(az)}}{\mu _\mathit{0}a}\right)}$$
(24)
and for the upwards intensity is,
$$I_2(z)=\frac{cab}{c+ab}e^{(e^{(az)}b)}C_\mathit{2}\frac{c+ab}{c+ab}e^{(e^{(az)}b)}C_\mathit{1}\frac{(a^2b^2c^2)(\mu _{\mathit{0}}^{}{}_{}{}^{2}+\mu _0\mu _2)}{16\pi (a^2b^2\mu _{\mathit{0}}^{}{}_{}{}^{2}c^2\mu _{\mathit{2}}^{}{}_{}{}^{2})}I_0e^{\left(\frac{c\mu _\mathit{2}e^{(az)}}{\mu _\mathit{0}a}\right)}$$
(25)
where $`C_1`$ and $`C_2`$ are arbitrary constants fixed with the boundary conditions. The direct component of the intensity is,
$$I^{dir}(z)=I_0e^{\left(\frac{ce^{az}\mu _2}{a\mu _\mathit{0}}\right)}.$$
(26)
For the specific choice $`c=0.8,a=3,b=0.2`$, the constants $`C_1`$ and $`C_2`$ that satisfy the boundary conditions are $`C_1=71.6787`$ and $`C_2=8.51812e05`$. The total optical path, for this particular set of parameters, is $`\tau 3`$. The optical functions $`\beta _T(z)`$ and $`\beta _T(z)`$ are shown in Fig. 3-I and the direct intensity in Fig. 3-II. The two-stream components of the diffused intensity are shown in Fig. 3-II.
## IV The numerical method
In contrast to the examples discussed above, where the optical properties of the media are known functions, in the general case the medium optical properties are only known, based on measurements or models, at a discrete set of points. In order to solve the radiative transfer equation, it is generally assumed that these optical properties are piecewise constant in a layer around the point, i.e. that the plane parallel layers are homogeneous piecewise. As we shall see later this should not be assumed in the case of atmospheric layers, where the optical properties show a strong (exponential) altitude dependence. We therefore propose to solve the problem numerically, using an adaptative step procedure and interpolating the optical values where necessary.
Equation (3) is an initial value problem and it is easily integrated knowing the incident intensity on the top $`I_0`$. Equations (12)-(13) and their boundary conditions constitute a two-point boundary problem which can be solved for instance with a shooting method as those explained in nr ; henar . For problems with linear boundary conditions, shooting methods are able to get the exact solution in few steps. We choose values for all the variables at one boundary, which must be consistent with any boundary condition there, but are otherwise arranged to depend on arbitrary free parameters whose initial values are guessed. We then integrate the ordinary differential equations with initial value methods, arriving at the other boundary, and adjust the free parameters at the starting point that zeros the discrepancies at the other boundary. Thanks to this procedure the problem is reduced to the solution of an initial value problem where no assumption of piecewise homogeneity is required.
The numerical integration of the equation for the direct component (3) and of the set of $`n`$ ordinary differential equations (5) for the diffused streams is done with a fifth-order Runge-Kutta (RG) scheme with variable step RG . Given a equation $`\frac{dI_i}{dz}=f(z,I_i)`$ where $`I_i(z)`$ is known, the RG scheme uses a weighted average of approximated values of $`f(z,I_i)`$ at several points within the interval $`(z,z+dz)`$ to obtain $`I_i(z+dz)`$. In contrast to other methods where piecewise homogeneity is required, this method takes into account the variation along the layer of integration of the diffused intensity streams, the direct component and the medium characteristics. For the adaptative step control we use a ”step doubling” technique: the local error is estimated by comparing a solution obtained with a fifth-order scheme and the one obtained with a fourth-order method. The integrating step is halved if this error is above a desired tolerance. Next, given the tabulated function $`\beta _{T,i}=\beta _T(z_i)`$ and $`\beta _{sca,i}=\beta _{sca}(z_i)`$ with $`i=1\mathrm{}N`$, we interpolate to obtain the new values at the locations required by the RG and the adaptative step doubling integration scheme. We have chosen a cubic interpolating spline which allows the interpolating formula to be smooth in the first derivative, and continuous in the second derivative, both within the interval and at its boundaries nr .
As an example, next we solve for the direct and diffused intensity in the two benchmark problems explained in section III. For problems with optical depths of the order of $`\tau 1`$, standard calculations based on the homogeneous layer approximation divide the medium into $`N=10`$ homogeneous layers Liou . We will analyze problems of $`\tau 3`$ and keep a similar ratio for our comparisons, we will then increase the number of layers. The tabulated optical properties $`\beta _T(z_i),\beta _{sca}(z_i)`$ at $`N`$ equidistant $`z_i`$ points are given as input.
A summary of the maximal relative error $`E(I)=|I_{dif}^{num}I_{dif}^{exact}|/I_{dif}^{exact}`$ for the benchmark problems as a function of the method and number of initial layers is given in table 1. We compare the relative error in the solution where piecewise homogeneity is assumed with the one obtained with our RG, step doubling, interpolating method. Figure 4 shows the relative error as a function of height $`z`$ for the step doubling interpolating technique with $`N=30`$ layers, for the linear (I) and exponential case (II). For a ratio of layers $`N/\tau 10`$ the piecewise homogeneous approximation leads to errors of the order of $`10\%`$ which can be reduced to $`1\%`$ by increasing this ratio to $`N/\tau 80`$. The step doubling interpolating RG technique permits the adaption of the numerical integration to the local characteristic length-scales of the medium and, for the same input, increases the accuracy of the solution between one and two orders of magnitude allowing for fast and accurate radiative transfer calculations.
## V Conclusions and discussion
We have considered the discrete-ordinate approach to the radiation transfer equation. In view of our calculations the error in the plane parallel approach, assuming $`N`$ piecewise homogeneous layers, with $`N10\tau `$ can go up to $`10\%`$ of the diffused intensity and thus the inhomogeneities in the atmosphere cannot be ignored. We have validated a general purpose numerical method that, based on an adaptative step integration and an interpolation of the local optical properties, can significantly improve (up to two orders of magnitude) the accuracy of the solution. This is furthermore of interest for practical applications, such as atmospheric radiation transfer, where the scattering and absorbing properties of the media are only known, based on experiments or theoretical models, at certain discrete points. Furthermore, this numerical method can be straightforward extended to multiple streams, non trivial boundary conditions and non uniform scattering function allowing for fast and accurate solutions of the radiative transfer equation.
## VI Acknowledgments.
M.-P. Z. is supported by the Instituto Nacional de T cnica Aerospacial (Spain). A.M.M. thanks to the Spanish Government for a Ramón y Cajal Research Fellowship. |
warning/0507/math0507026.html | ar5iv | text | # RSK Insertion for Set Partitions and Diagram Algebras
## 1 Introduction
Two fundamental identities in the representation theory of the symmetric group $`S_k`$ are
$$(\mathrm{a})n^k=\underset{\genfrac{}{}{0pt}{}{\lambda k}{\mathrm{}(\lambda )n}}{}f^\lambda d_\lambda ,\text{ and }(\mathrm{b})k!=\underset{\lambda k}{}(f^\lambda )^2,$$
(1.1)
where $`\lambda `$ varies over partitions of the integer $`k`$ of length $`\mathrm{}(\lambda )n`$, $`f^\lambda `$ is the number of standard Young tableaux of shape $`\lambda `$, and $`d_\lambda `$ is the number of column strict tableaux of shape $`\lambda `$ with entries from $`\{1,\mathrm{},n\}`$. The Robinson-Schensted-Knuth (RSK) insertion algorithm provides a a bijection between sequences $`(i_1,\mathrm{},i_k),1i_jn,`$ and pairs $`(P_\lambda ,Q_\lambda )`$ consisting of a standard Young tableau $`P_\lambda `$ of shape $`\lambda `$ and a column strict tableau $`Q_\lambda `$ of shape $`\lambda `$, thus providing a combinatorial proof of (1.1.a). If we restrict $`i_1,\mathrm{},i_k`$ to be a permutation of $`1,\mathrm{},k`$, then $`Q_\lambda `$ is a standard tableau and we have a proof of (1.1.b).
Identity (1.1.b) comes from the decomposition of the group algebra $`[S_k]`$ into irreducible $`S_k`$-modules $`V^\lambda ,\lambda k`$, where $`dim(V^\lambda )=f^\lambda `$ and the multiplicity of $`V^\lambda `$ in $`[S_k]`$ is also $`f^\lambda `$. Identity (1.1.a) comes from the Schur-Weyl duality between $`S_k`$ and the general linear group $`GL_n()`$ on the $`k`$-fold tensor product $`V^k`$ of the fundamental representation $`V`$ of $`GL_n()`$. There is an action of $`S_k`$ on $`V^k`$ by tensor place permutations, and via this action, $`[S_k]`$ is isomorphic to the centralizer algebra $`End_{GL_n()}(V^k)`$. As a bimodule for $`S_k\times GL_n()`$,
$$V^k\underset{\lambda k}{}S^\lambda V^\lambda ,$$
(1.2)
where $`S^\lambda `$ is an irreducible $`S_k`$-module of dimension $`f^\lambda `$ and $`V^\lambda `$ is an irreducible $`GL_r()`$-module of dimension $`d_\lambda `$. We get (1.1.a) by computing dimensions on each side of (1.2).
R. Brauer \[Br\] defined an algebra $`B_k(n)`$, which is isomorphic to the centralizer algebra of the orthogonal group $`O_n()GL_n()`$, when $`n2k`$, i.e., $`B_k(n)End_{O_n()}(V^k).`$ The dimension of the Brauer algebra is $`(2k1)!!=(2k1)(2k3)\mathrm{}31`$. Since $`O_n()GL_n()`$, their centralizers satisfy $`B_k(n)[S_k]`$. Berele \[Be\] generalized the RSK correspondence to give combinatorial proof of the $`B_k(n)`$-analog of (1.1.a). Sundaram \[Sun\] (see also \[Ter\]) gave a combinatorial proof of the $`B_k(n)`$-analog of (1.1.a).
We now take this restriction further to $`S_{n1}S_nO_n()GL_n()`$, where $`S_n`$ is viewed as the subgroup of permutation matrices in $`GL_n()`$ and $`S_{n1}S_n`$ corresponds to the permutations that fix $`n`$. Under this restriction, $`V`$ is the permutation representation of $`S_n`$, and when $`n2k`$, the centralizer algebras are the partition algebras,
$$A_k(n)End_{S_n}(V^k)\text{and}A_{k+\frac{1}{2}}(n)End_{S_{n1}}(V^k).$$
The partition algebra $`A_k(n)`$ first appeared independently in the work of Martin \[Mar1, Mar2, Mar3\] and Jones \[Jo\] arising from applications in statistical mechanics. See \[HR2\] for a survey paper on partition algebras.
For $`k\frac{1}{2}_{>0}`$ and $`n2k`$, $`A_k(n)`$ is a semisimple algebra over $``$ with a basis indexed by the set partitions of $`\{1,\mathrm{},2k\}`$. Thus, the dimension of $`A_k(n)`$ is the $`2k`$th Bell number $`B(2k)`$. Define
$$\mathrm{\Lambda }_n^k=\{\lambda n||\lambda |\lambda _1k\}$$
(1.3)
(these are partitions of $`n`$ with at most $`k`$ boxes below the first row of their Young diagram). Then for $`k_{>0}`$, the irreducible representations of $`A_k(n)`$ are indexed by partitions in the set $`\mathrm{\Lambda }_n^k`$, and the irreducible representations of $`A_{k+\frac{1}{2}}(n)`$ are indexed by partitions in the set $`\mathrm{\Lambda }_{n1}^k`$.
Using the Schur-Weyl duality between $`S_n`$ and $`A_k(n)`$ we get the identity
$$n^k=\underset{\lambda \mathrm{\Lambda }_n^k}{}f^\lambda m_k^\lambda ,$$
(1.4)
where $`m_k^\lambda `$ is the number of vacillating tableaux of shape $`\lambda `$ and length $`2k`$ (defined in Section 2.2), which are sequences integer partitions in the Bratteli diagram of $`A_k(n)`$. Identity (1.4) is the partition algebra analog of (1.1.a). In Section 3, we prove (1.4) using RSK column insertion and jeu de taquin. Decomposing $`A_k(n)`$ as a bimodule for $`A_k(n)A_k(n)`$ gives
$$B(2k)=\underset{\lambda }{}(m_k^\lambda )^2,$$
(1.5)
which is the $`A_k(n)`$-analog of (1.1.b). In Section 4, we give a bijective proof of (1.5) that contains as a special cases the RSK algorithms for $`S_k`$ and $`B_k(n)`$.
Martin and Rollet \[MR\] have given a different combinatorial proof of the second identity (1.5). Their bijection has the elegant property that pairs of paths in the difference between the Bratteli diagrams of $`A_k(\mathrm{})`$ and $`A_k(\mathrm{}1)`$ are in exact correspondence with the set partitions of $`\{1,\mathrm{},2k\}`$ into $`\mathrm{}`$ parts. The advantages of the correspondence in this paper are:
1. Our algorithm for (1.5) contains, as special cases, the known RSK correspondences for a number of diagram algebras which appear as subalgebras of $`A_k(n)`$:
1. The group algebra of the symmetric group $`S_k`$,
2. The Brauer algebra $`B_k(n)`$,
3. The Temperley-Lieb algebra $`T_k(n)`$,
4. The planar partition algebra $`P_k(n)`$.
5. The rook monoid algebra $`R_n`$ and the planar rook monoid algebra $`PR_n`$.
Thus we obtain combinatorial proofs of the analog of (1.5) for each of these algebras (see equations (5.1), (5.2), (5.3), and (5.4)).
2. Using Fomin growth diagrams we show that our algorithm for (1.5) is symmetric in the sense that if $`d(P,Q)`$ then $`\mathrm{flip}(d)(Q,P)`$ where $`\mathrm{flip}(d)`$ is the diagram $`d`$ flipped over its horizontal axis. This is the generalization of the property for the symmetric group that if $`\pi (P,Q)`$ then $`\pi ^1(Q,P)`$. As a consequence, we show that the number of symmetric diagrams equals the sum of the dimensions of the irreducible representations (for each of the diagram algebras mentioned in item 1 above).
3. Our algorithms for the bijections in (1.4) and (1.5) each use iterations of RSK insertion and jeu de taquin (see equations (3.4) and (4.4)).
The main idea for the algorithm in this paper came from a bijection of R. Stanley between fixed point free involutions in the symmetric group $`S_{2k}`$ and Brauer diagrams. This led to the insertion scheme of Sundaram in \[Sun\] for the Brauer algebra. After we distributed a preliminary version of this paper, R. Stanley and colleagues independently came out with the paper \[CDDSY\], which studies the crossing and nesting properties of an extended version of the insertion used in this paper for the bijection in (1.5). We have adopted the term “vacillating tableaux” from \[CDDSY\], and we use the crossing property of \[CDDSY\] to show that our algorithm restricts appropriately to the planar partition algebra.
## 2 The Partition Algebra and Vacillating Tableaux
For $`k_{>0}`$, let
$`A_k`$ $`=`$ $`\left\{\text{set partitions of }\{1,2,\mathrm{},k,1^{},2^{},\mathrm{},k^{}\}\right\},\text{and}`$ (2.1)
$`A_{k+\frac{1}{2}}`$ $`=`$ $`\left\{d\mathrm{\Pi }_{k+1}\right|(k+1)\text{ and }(k+1)^{}\text{ are in the same block}\}.`$ (2.2)
The propagating number of $`dA_k`$ is
$$\mathrm{pn}(d)=\left(\begin{array}{c}\text{the number of blocks in }d\text{ that contain both an element}\hfill \\ \text{of }\{1,2,\mathrm{},k\}\text{ and an element of }\{1^{},2^{},\mathrm{},k^{}\}\hfill \end{array}\right).$$
(2.3)
For convenience, represent a set partition $`dA_k`$ by a graph with $`k`$ vertices in the top row, labeled $`1,\mathrm{},k`$, and $`k`$ vertices in the bottom row, labeled $`1^{},\mathrm{},k^{}`$, with vertex $`i`$ and vertex $`j`$ connected by a path if $`i`$ and $`j`$ are in the same block of the set partition $`d`$. For example,
$$\text{1}\text{2}\text{3}\text{4}\text{5}\text{6}\text{7}\text{8}1^{}2^{}3^{}4^{}5^{}6^{}7^{}8^{}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{represents}\{\{1,2,4,2^{},5^{}\},\{3\},\{5,6,7,3^{},4^{},6^{},7^{}\},\{8,8^{}\},\{1^{}\}\},$$
and has propagating number 3. The graph representing $`d`$ is not unique.
Define the composition $`d_1d_2`$ of partition diagrams $`d_1,d_2A_k`$ to be the set partition $`d_1d_2A_k`$ obtained by placing $`d_1`$ above $`d_2`$, identifying the bottom dots of $`d_1`$ with the top dots of $`d_2`$, and removing any connected components that live entirely in the middle row. For example,
$$d_1d_2=\text{ }\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.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.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}.$$
Diagram multiplication makes $`A_k`$ into an associative monoid with identity, $`1\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\mathrm{}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.},`$ and the propagating number satisfies $`\mathrm{pn}(d_1d_2)\mathrm{min}(\mathrm{pn}(d_1),\mathrm{pn}(d_2)).`$
For $`k\frac{1}{2}_{>0}`$ and $`n`$, the partition algebra $`A_k(n)=\text{span-}\{dA_k\}`$ is an associative algebra over $``$ with basis $`A_k`$. Multiplication in $`A_k(n)`$ is defined by
$$d_1d_2=n^{\mathrm{}}(d_1d_2),$$
where $`\mathrm{}`$ is the number of blocks removed from the the middle row when constructing the composition $`d_1d_2`$. In the example above $`d_1d_2=n^2d_1d_2`$.
For each $`k\frac{1}{2}_{>0}`$, the following are submonoids of the partition monoid $`A_k`$:
$`S_k`$ $`=`$ $`\{dA_k|\mathrm{pn}(d)=k\},I_t=\{dA_k|\mathrm{pn}(d)t\},0<tk,`$
$`B_k`$ $`=`$ $`\{dA_k|\text{all blocks of }d\text{ have size 2}\},`$
$`R_k`$ $`=`$ $`\{dA_k|\begin{array}{c}\text{all blocks of }d\text{ have at most one vertex in }\{1,\mathrm{}k\}\hfill \\ \text{and at most one vertex in }\{1^{},\mathrm{}k^{}\}\hfill \end{array}\}.`$
Here, $`S_k`$ is the symmetric group, $`B_k`$ is the Brauer monoid, and $`R_k`$ is the rook monoid. A set partition is planar \[Jo\] if it can be represented as a graph without edge crossings inside of the rectangle formed by its vertices. The following are planar submonoids,
$$P_k=\{dA_k|d\text{ is planar}\},T_k=B_kP_k,PR_k=R_kP_k,1=S_kP_k.$$
Here, $`T_k`$ is the Temperley-Lieb monoid. Examples of diagrams in the various submonoids are:
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$$\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}R_7.\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}PR_7.$$
For each monoid, we make an associative algebra in the same way that we construct the partition algebra $`A_k(n)`$ from the partition monoid $`A_k(n)`$. For example, we obtain the the Brauer algebra $`B_k(n)`$, the Temperley-Lieb algebra $`T_k(n)`$, the group algebra of the symmetric group $`S_k`$ in this way. Multiplication in the rook monoid algebra $`R_k`$ is done without the coefficient $`n^{\mathrm{}}`$ (see \[Ha\]).
For $`\mathrm{}_{>0}`$, the Bell number $`B(\mathrm{})`$ is the number of set partitions of $`\{1,2,\mathrm{},\mathrm{}\}`$, the Catalan number is $`C(\mathrm{})=\frac{1}{\mathrm{}+1}\left(\genfrac{}{}{0pt}{}{2\mathrm{}}{\mathrm{}}\right)=\left(\genfrac{}{}{0pt}{}{2\mathrm{}}{\mathrm{}}\right)\left(\genfrac{}{}{0pt}{}{2\mathrm{}}{\mathrm{}+1}\right),`$ and $`(2\mathrm{})!!=(2\mathrm{}1)(2\mathrm{}3)\mathrm{}531.`$ These have generating functions (see \[Sta, 1.24f, and 6.2\]),
$$\underset{\mathrm{}0}{}B(\mathrm{})\frac{z^{\mathrm{}}}{\mathrm{}!}=\mathrm{exp}(e^z1),\underset{\mathrm{}0}{}C(\mathrm{}1)z^{\mathrm{}}=\frac{1\sqrt{14z}}{2z},$$
and
$$\underset{\mathrm{}0}{}(2(\mathrm{}1))!!\frac{z^{\mathrm{}}}{\mathrm{}!}=\frac{1\sqrt{12z}}{z}.$$
For $`k\frac{1}{2}_{>0}`$, the monoids have cardinality
$$\mathrm{Card}(A_k)=B(2k),\mathrm{Card}(P_k)=\mathrm{Card}(T_{2k})=C(2k),\text{for }k\frac{1}{2}_{>0}$$
and
$$\begin{array}{ccc}\mathrm{Card}(S_k)=k!,\hfill & & \mathrm{Card}(B_k)=(2k)!!,\hfill \\ \mathrm{Card}(R_k)=\underset{\mathrm{}=0}{\overset{k}{}}\left(\genfrac{}{}{0pt}{}{k}{\mathrm{}}\right)^2\mathrm{}!\hfill & & \mathrm{Card}(PR_k)=\left(\genfrac{}{}{0pt}{}{2k}{k}\right),\hfill \end{array}\text{for }k_{>0}\text{.}$$
### 2.1 Schur-Weyl Duality Between $`S_n`$ and $`A_k(n)`$
The irreducible representations of $`S_n`$ are indexed by integer partitions of $`n`$. If $`\lambda (\lambda _1,\mathrm{},\lambda _{\mathrm{}})(^0)^{\mathrm{}}`$ with $`\lambda _1\mathrm{}\lambda _{\mathrm{}}`$ and $`\lambda _1+\mathrm{}+\lambda _{\mathrm{}}=n`$, then $`\lambda `$ is a partition of $`n`$, denoted $`\lambda n`$. If $`\lambda n`$, then we write $`|\lambda |=n`$. If $`\lambda =(\lambda _1,\mathrm{},\lambda _{\mathrm{}})`$ and $`\mu =(\mu _1,\mathrm{},\mu _{\mathrm{}})`$ are partitions such that $`\mu _i\lambda _i`$ for each $`i`$, then we say that $`\mu \lambda `$, and $`\lambda /\mu `$ is the skew shape given by deleting the boxes of $`\mu `$ from the Young diagram of $`\lambda `$.
Let $`V`$ be the $`n`$-dimensional permutation representation of the symmetric group $`S_n`$. If we view $`S_{n1}S_n`$ as the subgroup of permutations that fix $`n`$, then $`V`$ is isomorphic to the left coset representation $`[S_n/S_{n1}]`$. Let $`V^k`$ be the $`k`$-fold tensor product representation of $`V`$, and let $`V^0=`$. From the “tensor identity” (see for example \[HR2\]), we have the following restriction-induction rule rule for $`i0`$,
$$V^{(i+1)}V^i[S_n/S_{n1}]\mathrm{Ind}_{S_{n1}}^{S_n}(\mathrm{Res}_{S_{n1}}^{S_n}(V^i)).$$
Thus, $`V^k`$ is obtained from $`k`$ iterations of restricting to $`S_{n1}`$ and inducing back to $`S_n`$.
Let $`V^\lambda `$ denote the irreducible representations of $`S_n`$ indexed by $`\lambda n`$. The restriction and induction rules for $`S_{n1}S_n`$ are given by
$`\mathrm{Res}_{S_{n1}}^{S_n}(V^\lambda )`$ $``$ $`{\displaystyle \underset{\mu (n1),\mu \lambda }{}}V^\mu ,\text{for }\lambda n`$ (2.5)
$`\mathrm{Ind}_{S_{n1}}^{S_n}(V^\mu )`$ $``$ $`{\displaystyle \underset{\lambda n,\mu \lambda ,}{}}V^\lambda ,\text{for }\mu (n1)\text{.}`$ (2.6)
In each case $`\lambda /\mu `$ consists of a single box. Starting with the trivial representation $`V^{(n)}`$ and iterating the restriction (2.5) and induction (2.6) rules, we see that the irreducible $`S_n`$-representations that appear in $`V^k`$ are labeled by the partitions in
$$\mathrm{\Lambda }_n^k=\{\lambda n||\lambda |\lambda _1k\},$$
(2.7)
and the irreducible $`S_{n1}`$-representations that appear in $`V^k`$ are labeled by the partitions in $`\mathrm{\Lambda }_{n1}^k`$.
There is an action of $`A_k(n)`$ on $`V^k`$ (see \[Jo, MR, HR2\]) that commutes with $`S_n`$ and maps $`A_k(n)`$ surjectively onto the centralizer $`End_{S_n}(V^k)`$. This generalizes the actions of $`[S_k]`$, by place permutations, and $`B_k(n)`$ on $`V^k`$. Furthermore, when $`n2k`$ we have
$$A_k(n)End_{S_n}(V^k)\text{and}A_{k+\frac{1}{2}}(n)End_{S_{n1}}(V^k).$$
(2.8)
By convention, we let $`A_0(n)=A_{\frac{1}{2}}(n)=`$. Since anything that commutes with $`S_n`$ on $`V^k`$ will also commute with $`S_{n1}`$, we have $`A_k(n)A_{k+\frac{1}{2}}(n)`$, and thus
$$A_0(n)A_{\frac{1}{2}}(n)A_1(n)A_{1\frac{1}{2}}(n)\mathrm{}A_{(k\frac{1}{2}}(n)A_k(n).$$
(2.9)
The Bratteli diagram for $`A_k(n)`$ consists of rows of vertices, with the rows labeled by $`0,\frac{1}{2},1,1\frac{1}{2},\mathrm{},k`$, such that the vertices in row $`i`$ are $`\mathrm{\Lambda }_n^i`$ and the vertices in row $`i+\frac{1}{2}`$ are $`\mathrm{\Lambda }_{n1}^i`$. Two vertices are connected by an edge if they are in consecutive rows and they differ by exactly one box. Figure 1 shows the Bratteli diagram for $`A_k(6)`$. Let $`n2k`$. By double centralizer theory (see, for example, \[HR2\]), we know that
* The irreducible representations of $`A_k(n)`$ can be indexed by $`\mathrm{\Lambda }_n^k`$, so we let $`M_k^\lambda `$ denote the irreducible $`A_k(n)`$ representation indexed by $`\lambda \mathrm{\Lambda }_n^k`$.
* The decomposition of $`V^k`$ as an $`S_n\times A_k(n)`$-bimodule is given by
$$V^k\underset{\lambda \mathrm{\Lambda }_n^k}{}V^\lambda M_k^\lambda .$$
(2.10)
* The dimension of $`M_k^\lambda `$ equals the multiplicity of $`V^\lambda `$ in $`V^k`$. The edges in the Bratteli diagram exactly follow the restriction and induction rules in (2.5), and (2.6). and so
$$m_k^\lambda =dim(M_k^\lambda )=\left\{\begin{array}{c}\text{ the number of paths from the top }\hfill \\ \text{ of the Bratteli diagram to }\lambda \text{ }\hfill \end{array}\right\}.$$
### 2.2 Vacillating Tableaux
The dimension of the irreducible $`S_n`$-module $`V^\lambda `$ equals the number $`f^\lambda `$ of standard tableaux of shape $`\lambda `$. A standard tableau of shape $`\lambda `$ is a filling of the Young diagram of $`\lambda `$ with the numbers $`1,2,\mathrm{},n`$ in such a way that each number appears exactly once, the rows increase from left to right, and the columns increase from top to bottom. We can identify a standard tableaux $`T_\lambda `$ of shape $`\lambda `$ with a sequence $`(\mathrm{}=\lambda ^{(0)},\lambda ^{(1)},\mathrm{},\lambda ^{(n)}=\lambda )`$ such that $`|\lambda ^{(i)}|=i`$, $`\lambda ^{(i)}\lambda ^{(i+1)}`$, and such that $`\lambda ^{(i)}/\lambda ^{(i1)}`$ is the box containing $`i`$ in $`T_\lambda `$. For example,
$$\text{1}\text{2}\text{3}\text{4}\text{5}\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }=(\mathrm{},\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }).$$
The sequence $`(\mathrm{}=\lambda ^{(0)},\lambda ^{(1)},\mathrm{},\lambda ^{(n)}=\lambda )`$ is a path in Young’s lattice, which is the Bratteli diagram for $`S_n`$. The number of standard tableaux $`f^\lambda `$ can be computed using the hook formula (see \[Sag, §3.10\]). We let $`𝒮𝒴𝒯(\lambda )`$ denote the set of standard tableaux of shape $`\lambda `$.
Let $`\lambda \mathrm{\Lambda }_n^k`$. A vacillating tableaux of shape $`\lambda `$ and length $`2k`$ is a sequence of partitions,
$$\left((n)=\lambda ^{(0)},\lambda ^{(\frac{1}{2})},\lambda ^{(1)},\lambda ^{(1\frac{1}{2})},\mathrm{},\lambda ^{(k\frac{1}{2})},\lambda ^{(k)}=\lambda \right),$$
satisfying, for each $`i`$,
* $`\lambda ^{(i)}\mathrm{\Lambda }_n^i`$ and $`\lambda ^{(i+\frac{1}{2})}\mathrm{\Lambda }_{n1}^i,`$
* $`\lambda ^{(i)}\lambda ^{(i+\frac{1}{2})}`$ and $`|\lambda ^{(i)}/\lambda ^{(i+\frac{1}{2})}|=1,`$
* $`\lambda ^{(i+\frac{1}{2})}\lambda ^{(i+1)}`$ and $`|\lambda ^{(i+1)}/\lambda ^{(i+\frac{1}{2})}|=1.`$
The vacillating tableaux of shape $`\lambda `$ correspond exactly with paths from the top of the Bratteli diagram to $`\lambda `$. Thus, if we let $`𝒱𝒯_k(\lambda )`$ denote the set of vacillating tableaux of shape $`\lambda `$ and length $`k`$, then
$$m_k^\lambda =dim(M_k^\lambda )=|𝒱𝒯_k(\lambda )|.$$
(2.11)
Let $`n2k`$, and for a partition $`\lambda `$, define $`\lambda ^{}`$ and $`\overline{\lambda }`$ as follows,
$$\begin{array}{cccc}\text{if}\hfill & \lambda =(\lambda _1,\mathrm{},\lambda _{\mathrm{}})n,\hfill & \text{then}\hfill & \lambda ^{}=(\lambda _2,\mathrm{},\lambda _{\mathrm{}})(n\lambda _1),\hfill \\ \text{if}\hfill & \lambda =(\lambda _1,\mathrm{},\lambda _{\mathrm{}})sn,\hfill & \text{then}\hfill & \overline{\lambda }=(ns,\lambda _1,\mathrm{},\lambda _{\mathrm{}})n.\hfill \end{array}$$
Since $`n2k`$ we are guaranteed that $`\overline{\lambda }`$ is a partition and that $`0|\lambda ^{}|k`$. The sets
$$\mathrm{\Lambda }_n^k=\{\lambda n||\lambda ^{}|k\}\text{and}\mathrm{\Gamma }_k=\{\lambda t|\mathrm{\hspace{0.17em}0}tk\}$$
(2.12)
are in bijection with one another using the maps,
$$\begin{array}{ccc}\mathrm{\Lambda }_n^k& & \mathrm{\Gamma }_k\\ \lambda & & \lambda ^{}\end{array}\text{and}\begin{array}{ccc}\mathrm{\Gamma }_k& & \mathrm{\Lambda }_n^k\\ \lambda & & \overline{\lambda }\end{array}.$$
(2.13)
Via these bijections, we can use either $`\mathrm{\Gamma }_k`$ or $`\mathrm{\Lambda }_n^k`$ to index the irreducible representations of $`A_k(n)`$. For example, the following sequences represent the same vacillating tableau $`P_\lambda `$, the first using diagrams from $`\mathrm{\Lambda }_n^k`$ and the second from $`\mathrm{\Gamma }_k`$,
$`P_\lambda `$ $`=`$ $`(\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }),`$
$`=`$ $`(\mathrm{},\mathrm{},\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },).`$
For our bijection in Section 3 we will use $`\mathrm{\Lambda }_n^k`$, and for our bijection in Section 4 we will use $`\mathrm{\Gamma }_k`$. The Bratteli diagram for $`A_k(n),n2k`$, is given in Figure 2, using labels from $`\mathrm{\Gamma }_k`$, along with the number of vacillating tableaux for each shape $`\lambda `$.
## 3 A Bijective Proof of $`n^k={\displaystyle \underset{\lambda \mathrm{\Lambda }_n^k}{}}f^\lambda m_k^\lambda `$
Comparing dimensions on both sides of the identity (2.10) gives
$$n^k=\underset{\lambda \mathrm{\Lambda }_n^k}{}f^\lambda m_k^\lambda ,\text{ for }n2k\text{.}$$
(3.1)
For example, when $`n=6`$ and $`k=3`$, the $`f^\lambda `$ in the bottom row of Figure 1 (see also Figure 2) are 1, 5, 9, 10, 5, 16, 10, the corresponding $`m_3^\lambda `$ are 5, 10, 6, 6, 1, 2, 1, and we have $`6^3=216=15+510+96+106+51+162+101.`$
To give a combinatorial proof of (3.1), we need to find a bijection of the form
$$\left\{(i_1,\mathrm{},i_k)\right|\mathrm{\hspace{0.17em}1}i_jn\}\underset{\lambda \mathrm{\Lambda }_n^k}{}𝒮𝒴𝒯(\lambda )\times 𝒱𝒯_k(\lambda ).$$
(3.2)
To do so, we construct an invertible function that turns a sequence $`(i_1,\mathrm{},i_k)`$ of numbers in the range $`1i_jn`$ into a pair $`(T_\lambda ,P_\lambda )`$ consisting of a standard tableaux $`T_\lambda `$ of shape $`\lambda `$ and vacillating tableaux $`P_\lambda `$ of shape $`\lambda `$ and length $`2k`$ for some $`\lambda \mathrm{\Lambda }_n^k`$. Our bijection uses jeu de taquin and RSK column insertion.
If $`T`$ is a standard tableau of shape $`\lambda n`$, then Schützenberger’s \[Scü\] jeu de taquin provides an algorithm for removing the box containing $`x`$ from $`T`$ and producing a standard tableau $`S`$ of shape $`\mu (n1)`$ with $`\mu \lambda `$ and entries $`\{1,\mathrm{},n\}\{x\}.`$ We only need a special case of jeu de taquin for our purposes. See \[Sag, §3.7\] or \[Sta, §A1.2\] for the full-strength version and its applications.
If $`S`$ is a standard tableau, let $`S_{i,j}`$ denote the entry of $`S`$ in row $`i`$ (numbered left-to-right) and column $`j`$ (numbered top-to-bottom). We say that a corner of $`S`$ is a box whose removal leaves the Young diagram of a partition. Thus the corners of $`S`$ are the boxes that are both at the end of a row and the end of a column. The following algorithm will delete $`x`$ from $`T`$ leaving a standard tableau $`S`$ with $`x`$ removed. We denote this process by $`x\stackrel{\mathrm{jdt}}{}T`$.
| 1. Let $`c=S_{i,j}`$ be the box containing $`x`$. |
| --- |
| 2. While $`c`$ is not a corner, do |
| | A. Let $`c^{}`$ be the box containing $`\mathrm{min}\{S_{i+1,j},S_{i,j+1}\}`$; |
| | B. Exchange the positions of $`c`$ and $`c^{}`$. |
| 3. Delete $`c`$. |
If only one of $`S_{i+1,j},S_{i,j+1}`$ exists at step 2.A, then the minimum is taken to be that single value. Below is an example of $`2\stackrel{\mathrm{jdt}}{}T`$,
$$\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{1}\text{2}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{3}\text{5}\text{4}\text{6}\text{7}\text{10}\text{8}\text{9}\text{12}\text{11}\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{1}\text{4}\text{3}\text{5}\text{2}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{6}\text{7}\text{10}\text{8}\text{9}\text{12}\text{11}\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{1}\text{4}\text{3}\text{5}\text{8}\text{6}\text{7}\text{10}\text{2}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{9}\text{12}\text{11}\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{1}\text{4}\text{3}\text{5}\text{8}\text{6}\text{7}\text{10}\text{11}\text{9}\text{12}\text{2}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{1}\text{4}\text{3}\text{5}\text{8}\text{6}\text{7}\text{10}\text{11}\text{9}\text{12}.$$
A bijective proof of the $`S_n`$ identity $`k!=_{\lambda k}(f^\lambda )^2`$ was originally found by Robinson \[Rob\] and later found, independently and in the form we present here, by Schensted \[Sch\]. Knuth \[Kn\] analyzed this algorithm and extended it to prove the identity (1.1.a). See \[Sta, §7 Notes\] for a nice history of the RSK algorithm. Let $`S`$ be a tableau of partition shape $`\mu `$, with $`|\mu |<n`$, with increasing rows and columns, and with distinct entries from $`\{1,\mathrm{},n\}`$. Let $`x`$ be a positive integer that is not in $`S`$. The following algorithm inserts $`x`$ into $`S`$ producing a standard tableau $`T`$ of shape $`\lambda `$ with $`\mu \lambda `$, $`|\lambda /\mu |=1`$, whose entries are the union of those from $`S`$ and $`\{x\}`$. We denote this process by $`x\stackrel{\mathrm{RSK}}{}S`$.
| 1. Let $`R`$ be the first row of $`S`$. |
| --- |
| 2. While $`x`$ is less than some element in $`R`$, do |
| | A. Let $`y`$ be the smallest element of $`R`$ greater than $`x`$; |
| | B. Replace $`yR`$ with $`x`$; |
| | C. Let $`x:=y`$ and let $`R`$ be the next row. |
| 3. Place $`x`$ at the end of $`R`$ (which is possibly empty). |
For example, here is the insertion of $`2`$ into the output of the jeu de taquin example above
$$\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{1}\text{4}\text{3}\text{5}\text{8}\text{6}\text{7}\text{10}\text{11}\text{9}\text{12}\text{2}\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{1}\text{2}\text{3}\text{5}\text{8}\text{6}\text{7}\text{10}\text{11}\text{9}\text{12}\text{4}\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{1}\text{2}\text{3}\text{5}\text{4}\text{6}\text{7}\text{10}\text{11}\text{9}\text{12}\text{8}\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{1}\text{2}\text{3}\text{5}\text{4}\text{6}\text{7}\text{8}\text{11}\text{9}\text{12}\text{10}\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{1}\text{2}\text{3}\text{5}\text{4}\text{7}\text{8}\text{11}\text{9}\text{10}\text{12}\text{6}$$
It is possible to invert the process of row insertion using row uninsertion. See \[Sag\] for details. In the example above, the number 2 and the leftmost tableau are the result of uninserting 10 from the rightmost tableau.
Given $`i_1,\mathrm{},i_k`$, with $`1i_jn`$, we will produce a pair $`(T_\lambda ,P_\lambda )`$, $`\lambda \mathrm{\Lambda }_n^k`$, consisting of a standard tableau $`T_\lambda `$ and a vacillating tableau $`P_\lambda `$. First, initialize the 0th tableau to be the standard tableaux of shape $`(n)`$, namely,
$$T^{(0)}\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{1}\text{2}\mathrm{}n.$$
(3.3)
Then, recursively define standard tableaux $`T^{(j+\frac{1}{2})}`$ and $`T^{(j+1)}`$ by
$$\begin{array}{ccc}\hfill T^{(j+\frac{1}{2})}& =& \left(i_{j+1}\stackrel{\mathrm{jdt}}{}T^{(j)}\right),\hfill \\ \hfill T^{(j+1)}& =& \left(i_{j+1}\stackrel{\mathrm{RSK}}{}T^{(j+\frac{1}{2})}\right),\hfill \end{array}0jk1.$$
(3.4)
Let $`\lambda ^{(j)}\mathrm{\Lambda }_n^j`$ be the shape of $`T^{(j)}`$, and let $`\lambda ^{(j+\frac{1}{2})}\mathrm{\Lambda }_{n1}^j`$ be the shape of $`T^{(j+\frac{1}{2})}`$. Then let
$$P_\lambda =(\lambda ^{(0)},\lambda ^{(\frac{1}{2})},\lambda ^{(1)},\lambda ^{(1\frac{1}{2})},\mathrm{},\lambda ^{(k)})\text{and}T_\lambda =T^{(k)},$$
(3.5)
so that $`P_\lambda `$ is a vacillating tableau of shape $`\lambda =\lambda ^{(k)}\mathrm{\Lambda }_n^k`$, and $`T_\lambda `$ is a standard tableaux of the same shape $`\lambda `$. We denote this iterative “delete-insert” process that associates the pair $`(T_\lambda ,P_\lambda )`$ to the sequence $`(i_1,\mathrm{},i_k)`$ by
$$(i_1,\mathrm{},i_k)\stackrel{\mathrm{DI}}{}(T_\lambda ,P_\lambda ).$$
(3.6)
Figure 3 shows the calculations which give
$$(2,4,3)\stackrel{\mathrm{DI}}{}(\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{1}\text{2}\text{3}\text{6}\text{4}\text{5},(\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ })).$$
###### Theorem 3.1.
The function $`(i_1,\mathrm{},i_k)\stackrel{\mathrm{DI}}{}(T_\lambda ,P_\lambda )`$ provides a bijection between sequences in $`\left\{(i_1,\mathrm{},i_k)\right|\mathrm{\hspace{0.17em}1}i_jn\}`$ and $`_{\lambda \mathrm{\Lambda }_n^k}𝒮𝒴𝒯(\lambda )\times 𝒱𝒯_k(\lambda )`$ and thus gives a combinatorial proof of (3.1).
###### Proof.
We prove the theorem by constructing the inverse of $`\stackrel{\mathrm{DI}}{}`$.
Let $`\lambda ^{(j+\frac{1}{2})}\lambda ^{(j+1)}`$ with $`\lambda ^{(j+1)}\mathrm{\Lambda }_n^{(j+1)}`$ and $`\lambda ^{(j+\frac{1}{2})}\mathrm{\Lambda }_{n1}^j`$, and let $`T^{(j+1)}`$ be a standard tableau of shape $`\lambda ^{(j+1)}`$. We can uniquely determine $`i_{j+1}`$ and a tableau $`T^{(j+\frac{1}{2})}`$ of shape $`\lambda ^{(j+\frac{1}{2})}`$ such that $`T^{(j+1)}=(i_{j+1}\stackrel{\mathrm{RSK}}{}T^{(j+\frac{1}{2})}).`$ To do this let $`b`$ be the box in $`\lambda ^{(j+1)}/\lambda ^{(j+\frac{1}{2})}`$. Then let $`i_{j+1}`$ and $`T^{(j+\frac{1}{2})}`$ be the result of uninserting the number in box $`b`$ of $`T^{(j+1)}`$ (using the fact that RSK insertion is invertible).
Now let $`T^{(j+\frac{1}{2})}`$ be a tableau of shape $`\lambda ^{(j+\frac{1}{2})}\mathrm{\Lambda }_{n1}^{(j)}`$ with increasing rows and columns and entries $`\{1,\mathrm{},n\}\{i_{j+1}\}`$ and let $`\lambda ^{(j)}\lambda ^{(j+\frac{1}{2})}`$ with $`\lambda ^{(j)}\mathrm{\Lambda }_n^j`$. We can uniquely produce a standard tableau $`T^{(j)}`$ such that $`T^{(j+\frac{1}{2})}=(i_{j+1}\stackrel{\mathrm{jdt}}{}T^{(j)})`$. To do this, let $`b`$ be the box in $`\lambda ^{(j)}/\lambda ^{(j+\frac{1}{2})}`$, put $`i_{j+1}`$ in position $`b`$ of $`T^{(j+\frac{1}{2})}`$, and perform the inverse of jeu de taquin to produce $`T^{(j)}`$. That is, move $`i_{j+1}`$ into a standard position by iteratively swapping it with the larger of the numbers just above it or just to its left.
Given $`\lambda \mathrm{\Lambda }_n^k`$ and $`(P_\lambda ,T_\lambda )𝒮𝒴𝒯(\lambda )\times 𝒱𝒯_k(\lambda )`$, we apply the process above to $`\lambda ^{(k\frac{1}{2})}\lambda ^{(k)}`$ and $`T^{(k)}=T_\lambda `$ producing $`i_k`$ and $`T^{(k1)}`$. Continuing this way, we can produce $`i_k,i_{k1},\mathrm{},i_1`$ and $`T^{(k)},T^{(k1)},\mathrm{},T^{(1)}`$ such that $`(i_1,\mathrm{},i_k)\stackrel{\mathrm{DI}}{}(T_\lambda ,P_\lambda ).`$
## 4 A Bijective Proof of $`B(2k)={\displaystyle \underset{\lambda \mathrm{\Gamma }_k}{}}(m_k^\lambda )^2`$
If we view $`A_k(n)`$ as a bimodule for $`A_k(n)A_k(n)`$, with the first tensor factor acting by left multiplication on $`A_k(n)`$ and the second tensor factor acting by right multiplication on $`A_k(n)`$, then the decomposition into irreducibles for $`A_k(n)A_k(n)`$ is $`A_k(n)_{\lambda \mathrm{\Gamma }_k}M_k^\lambda \overline{M}_k^\lambda ,`$ where $`\overline{M}_k^\lambda `$ is the irreducible right $`A_k(n)`$-module indexed by $`\lambda `$ (so $`M_k^\lambda \overline{M}_k^\lambda `$). This is a generalization of the decomposition of the group algebra of a finite group. Comparing dimensions on both sides gives,
$$B(2k)=\underset{\lambda \mathrm{\Gamma }_k}{}(m_k^\lambda )^2.$$
(4.1)
For example using the numbers $`m_3^\lambda `$ from the bottom row of Figure 2, we have $`B(6)=203=5^2+10^2+6^2+6^2+1^2+2^2+1^2.`$
To give a combinatorial proof of (4.1), we find a bijection of the form
$$A_k\underset{\lambda \mathrm{\Gamma }_k}{}𝒱𝒯_k(\lambda )\times 𝒱𝒯_k(\lambda ),$$
(4.2)
by constructing a function that takes a set partition $`dA_k`$ and produces a pair $`(P_\lambda ,Q_\lambda )`$ of vacillating tableaux. In Section 5 we show that our bijection restricts to work for all the diagram subalgebras described in Section 2. In particular, if $`d`$ corresponds to a permutation in $`S_k`$ then our bijection is the usual RSK algorithm described in Section 3.2.
We will draw diagrams $`dA_k`$ using a standard representation as single row with the vertices in order $`1,\mathrm{},2k`$, where we relabel vertex $`j^{}`$ with with the label $`2kj+1`$. We draw the edges of the standard representation of $`dA_k`$ in a specific way: connect vertices $`i`$ and $`j`$, with $`ij`$, if and only if $`i`$ and $`j`$ are related in $`d`$ and there does not exist $`k`$ related to $`i`$ and $`j`$ with $`i<k<j`$. In this way, each vertex is connected only to its nearest neighbors in its block. For example,
$$\text{1}\text{2}\text{3}\text{4}\text{5}\text{6}1^{}2^{}3^{}4^{}5^{}6^{}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{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We label each edge $`e`$ of the diagram $`d`$ with $`2k+1\mathrm{}`$, where $`\mathrm{}`$ is the right vertex of $`e`$. Define the insertion sequence of a diagram to be the sequence $`E=(E_j)`$ indexed by $`j`$ in the sequence $`\frac{1}{2},1,1\frac{1}{2},\mathrm{},2k1,2k\frac{1}{2},2k`$, where
$`E_j`$ $`=`$ $`\{\begin{array}{cc}a,\hfill & \text{if vertex }j\text{ is the left endpoint of edge }a,\hfill \\ \mathrm{},\hfill & \text{if vertex }j\text{ is not a left endpoint},\hfill \end{array}`$
$`E_{j\frac{1}{2}}`$ $`=`$ $`\{\begin{array}{cc}a,\hfill & \text{if vertex }j\text{ is the right endpoint of edge }a,\hfill \\ \mathrm{},\hfill & \text{if vertex }j\text{ is not a right endpoint}.\hfill \end{array}`$
###### Example 4.1.
The edge labeling for the diagram of the set partition $`\{\{1,3,4^{}\},\{2,1^{}\},`$ $`\{4,3^{},2^{}\}\}`$ is given by
$$\text{1}\text{2}\text{3}\text{4}1^{}2^{}3^{}4^{}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{1}\text{2}\text{3}\text{4}5678\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\text{.}\t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and the insertion sequence is
$$\begin{array}{ccccccccccccccccc}j& {\scriptscriptstyle \frac{1}{2}}& 1& 1{\scriptscriptstyle \frac{1}{2}}& 2& 2{\scriptscriptstyle \frac{1}{2}}& 3& 3{\scriptscriptstyle \frac{1}{2}}& 4& 4{\scriptscriptstyle \frac{1}{2}}& 5& 5{\scriptscriptstyle \frac{1}{2}}& 6& 6{\scriptscriptstyle \frac{1}{2}}& 7& 7{\scriptscriptstyle \frac{1}{2}}& 8\\ E_j& \mathrm{}& \text{6}& \mathrm{}& \text{1}& \text{6}& \text{4}& \mathrm{}& \text{3}& \text{4}& \mathrm{}& \text{3}& \text{2}& \text{2}& \mathrm{}& \text{1}& \mathrm{}\end{array}.$$
The insertion sequence of a standard diagram completely determines the edges, and thus the connected components, of the diagram, so the following proposition follows immediately.
###### Proposition 4.2.
$`dA_k`$ is completely determined by its insertion sequence.
For $`dA_k`$ with insertion sequence $`E=(E_j)`$, we will produce a pair $`(P_\lambda ,Q_\lambda )`$ of vacillating tableaux. Begin with the empty tableaux,
$$T^{(0)}=\mathrm{}.$$
(4.3)
Then recursively define standard tableaux $`T^{(j+\frac{1}{2})}`$ and $`T^{(j+1)}`$ by
$$\begin{array}{ccc}\hfill T^{(j+\frac{1}{2})}& =& \{\begin{array}{cc}E_{j+\frac{1}{2}}\stackrel{\mathrm{jdt}}{}T^{(j)},\hfill & \text{if }E_{j+\frac{1}{2}}\mathrm{},\hfill \\ T^{(j)},\hfill & \text{if }E_{j+\frac{1}{2}}=\mathrm{}\text{,}\hfill \end{array}\hfill \\ \hfill T^{(j+1)}& =& \{\begin{array}{cc}E_{j+1}\stackrel{\mathrm{RSK}}{}T^{(j+\frac{1}{2})},\hfill & \text{if }E_{j+1}\mathrm{},\hfill \\ T^{(j+\frac{1}{2})},\hfill & \text{if }E_{j+1}=\mathrm{}\text{,}\hfill \end{array}\hfill \end{array}0j2k1.$$
(4.4)
Let $`\lambda ^{(i)}`$ be the shape of $`T^{(i)}`$, let $`\lambda ^{(i+\frac{1}{2})}`$ be the shape of $`T^{(i+\frac{1}{2})}`$, and let $`\lambda =\lambda ^{(k)}`$. Define
$`Q_\lambda `$ $`=`$ $`(\mathrm{},\lambda ^{(\frac{1}{2})},\lambda ^{(1)},\mathrm{},\lambda ^{(k\frac{1}{2})},\lambda ^{(k)})𝒱𝒯(\lambda ),`$
$`P_\lambda `$ $`=`$ $`(\lambda ^{(2k)},\lambda ^{(2k\frac{1}{2})},\mathrm{},\lambda ^{(k+\frac{1}{2})},\lambda ^{(k)})𝒱𝒯(\lambda ).`$
In this way, we associate a pair of vacillating tableaux $`(P_\lambda ,Q_\lambda )`$ to a set partition $`dA_k`$, which we denote by
$$d(P_\lambda ,Q_\lambda ).$$
In Figure 4, we see the insertion of the diagram in Example 4.1. The result of this is to pair the diagram in Example 4.1 with the following pair of vacillating tableaux,
$`Q_\lambda `$ $`=`$ $`(\mathrm{},\mathrm{},\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }),`$
$`P_\lambda `$ $`=`$ $`(\mathrm{},\mathrm{},\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }).`$
###### Remark 4.3.
We have numbered the edges of the standard diagram of $`d`$ in increasing order from right to left, so if $`E_{j+\frac{1}{2}}\mathrm{}`$, then $`E_{j+\frac{1}{2}}`$ will be the largest element of $`T^{(j)}`$. Thus, in $`T^{(j+\frac{1}{2})}=(E_{j+\frac{1}{2}}\stackrel{\mathrm{jdt}}{}T^{(j)})`$ we know that $`E_{j+\frac{1}{2}}`$ is in a corner box, and jeu de taquin simply deletes that box.
###### Theorem 4.4.
The function $`d(P_\lambda ,Q_\lambda )`$ provides a bijection between the set partitions $`A_k`$ and pairs of vacillating tableaux $`_{\lambda \mathrm{\Gamma }_k}𝒱𝒯(\lambda )\times 𝒱𝒯_k(\lambda )`$ and thus gives a combinatorial proof of (4.1).
###### Proof.
We prove the theorem by constructing the inverse of $`d(P_\lambda ,Q_\lambda )`$. First, we use $`Q_\lambda `$ followed by $`P_\lambda `$ in reverse order to construct the sequence $`\lambda ^{(\frac{1}{2})},\lambda ^{(1)},\mathrm{},\lambda ^{(2k\frac{1}{2})},\lambda ^{(2k)}`$. We initialize $`T^{(2k)}=\mathrm{}`$.
We now show how to construct $`T^{(i+\frac{1}{2})}`$ and $`E_{i+1}`$ so that $`T^{(i+1)}=(E_{i+1}\stackrel{\mathrm{RSK}}{}T^{(i+\frac{1}{2})})`$. If $`\lambda ^{(i+\frac{1}{2})}=\lambda ^{(i+1)}`$, then let $`T^{(i+\frac{1}{2})}=T^{(i+1)}`$ and $`E_{i+1}=\mathrm{}`$. Otherwise, $`\lambda ^{(i+1)}/\lambda ^{(i+\frac{1}{2})}`$ is a box $`b`$, and we use uninsertion on the value in $`b`$ to produce $`T^{(i+\frac{1}{2})}`$ and $`E_{i+1}`$ such that $`T^{(i+1)}=(E_{i+1}\stackrel{\mathrm{RSK}}{}T^{(i+\frac{1}{2})})`$. Since we uninserted the value in position $`b`$, we know that $`T^{(i+\frac{1}{2})}`$ has shape $`\lambda ^{(i+\frac{1}{2})}`$.
We then show how to construct $`T^{(i)}`$ and $`E_{i+\frac{1}{2}}`$ so that $`T^{(i+\frac{1}{2})}=(E_{i+\frac{1}{2}}\stackrel{\mathrm{jdt}}{}T^{(i)})`$. If $`\lambda ^{(i)}\lambda ^{(i+\frac{1}{2})}`$, then let $`T^{(i)}=T^{(i+\frac{1}{2})}`$ and $`E_{i+\frac{1}{2}}=\mathrm{}`$. Otherwise, $`\lambda ^{(i)}/\lambda ^{(i+\frac{1}{2})}`$ is a box $`b`$. Let $`T^{(i)}`$ be the tableau of shape $`\lambda ^{(i)}`$ with the same entries as $`T^{(i+1)}`$ and having the entry $`2ki`$ in box $`b`$. Let $`E_{i+\frac{1}{2}}=2ki`$. At any given step $`i`$ (as we work our way down from step $`2k`$ to step $`\frac{1}{2}`$), $`2ki`$ is the largest value added to the tableau thus far, so $`T^{(i)}`$ is standard. Furthermore, $`T^{(i+\frac{1}{2})}=(E_{i+\frac{1}{2}}\stackrel{\mathrm{jdt}}{}T^{(i)})`$, since $`E_{i+\frac{1}{2}}=2ki`$ is already in a corner and thus jeu de taquin simply deletes it.
This iterative process will produce $`E_{2k},E_{2k1},\mathrm{},E_{\frac{1}{2}}`$ which completely determines the set partition $`d`$. By the way we have constructed $`d`$, we have $`d(P_\lambda ,Q_\lambda )`$. ∎
## 5 Properties of the Insertion Algorithm
In this section we investigate a number of the nice properties of the insertion algorithm from Section 4. First, we note that our algorithm restricts to the insertion algorithms for the subalgebras discussed in Section 1:
###### Remark 5.1.
* If $`dS_k`$ is a permutation diagram, then the standard representation of $`d`$ pairs vertices $`1,\mathrm{}k`$ with vertices $`k+1,\mathrm{},2k`$. Our algorithm has us insert on steps $`j`$ with $`1jk`$, delete on steps $`j+\frac{1}{2}`$ with $`kj2k1`$, and do nothing on all other steps. If we ignore these “do-nothing” steps, then our algorithm becomes the usual RSK insertion and $`P_\lambda `$ and $`Q_\lambda `$ are the insertion and recording tableaux, respectively.
* If $`dB_k`$ is a Brauer diagram, then each vertex in the standard representation of $`d`$ is incident to exactly one edge, so exactly one of $`E_{(j\frac{1}{2})}`$ or $`E_j`$ will be nonempty for each $`1j2k`$. If we ignore the steps where we do nothing, then for each vertex we either insert or delete the edge number associated to that vertex. This is precisely the insertion scheme of Sundaram \[Sun\] (see also \[Ter\]). This algorithm produces “oscillating tableaux,” which either increase or decrease by one box on each step. If we insert all the Brauer diagrams, we produce the Bratteli diagram for $`B_k(n)`$ (see for example \[HR1\]).
* If $`dR_k`$ is a rook monoid diagram, then our algorithm has us insert or do nothing on steps $`0<jk`$ and delete or do nothing on steps $`j+\frac{1}{2}`$ with $`kj2k1`$. If we insert all rook monoid diagrams, we obtain the Bratteli diagram for $`R_k`$ (see \[Ha\]) .
* If $`dA_{k\frac{1}{2}}`$ then the standard representation of $`d`$ has an edge, labeled $`k`$, connecting vertices $`k`$ and $`k+1`$. At step $`k`$ in the insertion process, we will insert this edge into $`T^{(k\frac{1}{2})}`$. The label $`k`$ is larger than all the entries in $`T^{(k\frac{1}{2})}`$, since the entries that are in $`T^{(k\frac{1}{2})}`$ correspond to edges whose right endpoint is larger than $`k+1`$. Thus we insert this entry at the end of the first row and in the next step, $`k+\frac{1}{2}`$, we immediately delete it, so that $`\lambda ^{(k\frac{1}{2})}=\lambda ^{(k+\frac{1}{2})}`$. If we remove $`\lambda ^{(k)}`$ from our pair of paths, we obtain a pair of vacillating tableaux $`(P_{\lambda ^{(k{\scriptscriptstyle \frac{1}{2}})}},Q_{\lambda ^{(k{\scriptscriptstyle \frac{1}{2}})}})`$ of length $`2k1`$, and we get a combinatorial formula for the following dimension identity for $`A_{k\frac{1}{2}}(n)`$ involving the odd Bell numbers,
$$B(2k1)=\underset{\lambda \mathrm{\Gamma }_{k{\scriptscriptstyle \frac{1}{2}}}}{}(m_{k\frac{1}{2}}^\lambda )^2.$$
(5.1)
###### Proposition 5.2.
If $`dA_k`$ such that $`d(P_\lambda ,Q_\lambda )`$, then $`|\lambda |=\mathrm{pn}(d)`$, where $`\mathrm{pn}(d)`$ is the propagating number defined in (2.3).
###### Proof.
If we draw a vertical line between vertex $`k`$ and $`k+1`$ in the standard representation of $`d`$, then $`\mathrm{pn}(d)`$ is the number of edges that cross this line. The labels of these edges are exactly the numbers which have been inserted but not deleted at step $`k`$ of our algorithm. Thus the tableau $`T^{(k)}`$ has exactly these numbers in it and so $`\lambda ^{(k)}`$ consists of $`\mathrm{pn}(d)`$ boxes. ∎
Recall that $`I_t`$ is spanned by the diagrams whose propagating number satisfies $`\mathrm{pn}(d)t`$. The chain of ideals $`I_0(n)I_1(n)\mathrm{}I_k(n)=A_k(n)`$ is a key feature of the “basic construction” of $`A_k(n)`$ (see \[HR2, Mar3\]). The irreducible components of $`I_t(n)/(I_{t1}(n))`$ are those indexed by $`\lambda \mathrm{\Gamma }_k`$ with $`|\lambda |=t`$. The previous proposition tells us that our insertion algorithm respects this basic construction and gives a combinatorial proof of the following dimension identity for $`I_t(n)/(I_{t1}(n))`$,
$$\mathrm{Card}\left(\{dA_k|\mathrm{pn}(d)=t\}\right)=\underset{\genfrac{}{}{0pt}{}{\lambda \mathrm{\Gamma }_k}{\mathrm{pn}(d)=t}}{}(m_k^\lambda )^2.$$
(5.2)
###### Proposition 5.3.
Let $`d(P,Q)=(\lambda ^{(0)},\mathrm{},\lambda ^{(2k)}).`$ Then $`d`$ is planar if and only if the length of every partition in the sequence satisfies $`\mathrm{}(\lambda ^{(i)})1`$.
###### Proof.
This is a special case of the crossing and nesting theorem in \[CDDSY\]. If $`d`$ is planar, then the standard representation of $`d`$ is also planar. For example,
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Let $`a<b`$ be the labels of two edges in a planar diagram $`d`$ and suppose that $`b`$ is inserted before $`a`$ in the insertion sequence of $`d`$. Since $`a<b`$ the right vertex of $`b`$ is left of the right vertex of $`a`$. Since $`b`$ is inserted before $`a`$, the left vertex of $`b`$ is left of the left vertex of $`a`$. Since $`d`$ is planar, the right vertex of $`b`$ must then be left of the left vertex of $`a`$. Thus, $`b`$ is inserted and deleted before $`a`$ is inserted. This means that when we are inserting an edge label $`a`$ into $`T^{(i)}`$, we know that $`a`$ is larger than all the entries of $`T^{(i)}`$, and so $`a`$ placed at the end of the first row. Thus at each step of the insertion, we only ever obtain the empty tableau or a tableau consisting of 1 row. Conversely, to obtain a partition that has more than one column, there must be a pair $`a<b`$ such that $`a`$ is inserted after $`b`$ is inserted but before $`b`$ is deleted. This is necessary for $`a`$ to bump $`b`$ into the next row. By the reasoning above, this forces an edge crossing in the partition diagram. ∎
###### Remark 5.4.
* The planar partition diagrams $`dP_k`$ are exactly paired with the paths $`(P,Q)`$ whose vertices contain only length 0 or 1 partitions. The planar partition algebra $`P_k(n)`$ is isomorphic to the Temperley-Lieb algebra $`T_{2k}(n)`$, $`k\frac{1}{2}_{>0}`$, (see \[Jo\] or \[HR2\]) and has the Bratteli diagram shown in Figure 5. Our algorithm gives a combinatorial proof of the following dimension identity for $`T_{2k}(n)`$,
$$C(2k)=\underset{\genfrac{}{}{0pt}{}{\lambda \mathrm{\Lambda }_k}{\mathrm{}(\lambda )1}}{}\left(\left(\genfrac{}{}{0pt}{}{2k}{k|\lambda |}\right)\left(\genfrac{}{}{0pt}{}{2k}{k|\lambda |1}\right)\right)^2,k\frac{1}{2}_{>0},$$
(5.3)
where $``$ is the floor function and $`\left(\genfrac{}{}{0pt}{}{2k}{k|\lambda |}\right)\left(\genfrac{}{}{0pt}{}{2k}{k|\lambda |1}\right)`$ is the dimension of the irreducible $`T_{2k}(n)`$-module labeled by $`\lambda `$ (this is also the number of paths to $`\lambda `$ in the Bratteli diagram of $`T_{2k}(n)`$).
* The planar Brauer diagrams $`T_k`$ span another copy of the Temperley-Lieb algebra $`T_k(n)`$ inside $`A_k(n)`$, and our insertion algorithm when applied to these diagrams gives a different bijective proof of the identity in (b).
* In the insertion for planar rook monoid diagrams in $`PR_k`$, we insert or do nothing for steps $`j`$ with $`0<jk`$, we delete or do nothing for steps $`j+{\scriptscriptstyle \frac{1}{2}}`$ with $`kj<2k`$, and we do nothing on all other steps. Thus, the Bratteli diagram, which is shown in Figure 5, has levels $`j`$ and $`j+{\scriptscriptstyle \frac{1}{2}}`$ merged and is isomorphic to Pascal’s triangle. The irreducible representation associated to the partition of shape $`\lambda =(\mathrm{})`$ has dimension $`\left(\genfrac{}{}{0pt}{}{k}{\mathrm{}}\right)`$ (see \[HH\] for the representation theory of $`PR_k`$). Our algorithm gives an RSK insertion proof of following the well-known binomial identity,
$$\left(\genfrac{}{}{0pt}{}{2k}{k}\right)=\underset{\mathrm{}=0}{\overset{k}{}}\left(\genfrac{}{}{0pt}{}{k}{\mathrm{}}\right)^2.$$
(5.4)
Finally, we show that the algorithm has the symmetry property of the usual RSK algorithm for the symmetric group. Namely, if $`\pi (P,Q)`$ for $`\pi S_k`$, then $`\pi ^1(Q,P)`$. Using diagrams, the inverse of $`\pi `$ is simply the diagram of $`\pi `$ reflected over its horizontal axis. We define the flip of any partition diagram to be this reflection, so for example,
$$\begin{array}{ccccc}\hfill d& =& 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& & & & \\ \hfill \mathrm{flip}(d)& =& 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=& 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Notice that in the standard representation of $`d`$, a flip corresponds to a reflection over the vertical line between vertices $`k`$ and $`k+1`$. Our goal is to show that if $`d(P,Q)`$, then $`\mathrm{flip}(d)(Q,P)`$.
Given a partition $`dA_k`$, construct a triangular grid in the integer lattice $`\times `$ that contains the points in the triangle whose vertices are $`(0,0),(2k,0)`$ and $`(0,2k)`$. Number the columns $`1,\mathrm{},2k`$ from left to right and the rows $`1,\mathrm{},2k`$ from bottom to top. Place an X in the box in column $`i`$ and row $`j`$ if and only if, in the one-row diagram of $`d`$, vertex $`i`$ is the left endpoint of edge $`j`$. We then label the vertices of the diagram on the bottom row and left column with the empty partition $`\mathrm{}`$. For example, in our diagram $`d`$ from Example 4.1 we have
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Note that the triangular array completely determines the partition diagram and vice versa.
Now we inductively label the remaining vertices using the local rules of S. Fomin (see \[Ry1, Ry2\] or \[Sta, 7.13\] and the references there). If a box is labeled with $`\mu ,\nu ,`$ and $`\lambda `$ as shown below. Then we add the label $`\rho `$ according the rules that follow:
................................................................................................................................................................................................................................................................................................$`\lambda `$$`\mu `$$`\nu `$$`\rho `$
1. If $`\mu \nu `$, let $`\rho =\mu \nu `$, i.e., $`\rho _i=\mathrm{max}(\mu _i,\nu _i)`$ .
2. If $`\mu =\nu `$, $`\lambda \mu `$, and $`\lambda \mu `$, then this will automatically imply that $`\mu `$ can be obtained from $`\lambda `$ by adding a box to $`\lambda _i`$. Let $`\rho `$ be obtained from $`\mu `$ by adding a box to $`\mu _{i+1}`$.
3. If $`\mu =\nu =\lambda `$, then if the square does not contain an X, let $`\rho =\lambda `$, and if the square does contain an X, let $`\rho `$ be obtained from $`\lambda `$ by adding 1 to $`\lambda _1`$.
Using these rules, we can uniquely label every corner, one sep at a time. The resulting diagram is called the growth diagram $`G_d`$ for $`d`$. The diagram corresponding to the above example is
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Let $`P_d`$ denote the chain of partitions that follow the staircase path on the diagonal of $`G_d`$ from $`(0,2k)`$ to $`(k,k)`$ and let $`Q_d`$ denote the chain of partitions that follow the staircase path on the diagonal of $`G_d`$ from $`(2k,0)`$ to $`(k,k)`$. The pair $`(P_d,Q_d)`$ represents a pair of vacillating tableaux whose shape is the partition at $`(k,k)`$. The staircase path in our growth diagram above is the same as the path we get from insertion in Figure 4.
###### Theorem 5.5.
Let $`dA_k`$ with $`d(P,Q)`$. Then $`P_d=P`$ and $`Q_d=Q`$.
###### Proof.
Turn the diagram $`dA_k`$ into a diagram $`d^{}`$ on $`4k`$ vertices by splitting each vertex $`i`$ into two vertices labeled by $`i\frac{1}{2}`$ and $`i`$. If there is an edge from vertex $`j`$ to vertex $`i`$ in $`d`$, with $`j<i`$, let $`j`$ be adjacent to $`i\frac{1}{2}`$ in $`d^{}`$. If there is an edge from vertex $`j`$ to vertex $`i`$ in $`d`$, with $`j>i`$, let $`i`$ be adjacent to $`j\frac{1}{2}`$ in $`d^{}`$. Thus, in our example we have,
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In this way we turn each diagram in $`A_k`$ into a partial matching in $`A_{2k}`$. In his dissertation, T. Roby \[Ry1\] shows that the up-down tableaux coming from growth diagrams are equivalent to those from RSK insertion of Brauer diagrams (matchings). Dulucq and Sagan \[DS\] extend this insertion to work for partial matchings (and other generalizations for skew oscillating tableaux) and Roby \[Ry2\] shows that the growth diagrams also work for the Duluc-Sagan insertion of partial matchings. Furthermore, the insertion sequences that come from our diagram $`d`$ are the same as the insertion sequences that come from $`d^{}`$ using the methods of \[Sun, DS\]. ∎
A key advantage of the use of growth diagrams is that the symmetry of the algorithm is nearly obvious. We have that $`i`$ is the left endpoint of the edge labeled $`j`$ in $`d`$ if and only if $`j`$ is the left endpoint of the edge labeled $`i`$ in $`\mathrm{flip}(d)`$. Thus the growth diagram of $`G_d`$ is the reflection over the line $`y=x`$ of the growth diagram of $`G_{\mathrm{flip}(d)}`$, and so $`P_d=Q_{\mathrm{flip}(d)}`$ and $`Q_d=P_{\mathrm{flip}(d)}`$. Combining this point with the Theorem 5.5 gives
###### Corollary 5.6.
If $`d(P,Q)`$ then $`\mathrm{flip}(d)(Q,P)`$.
We say that $`dA_k`$ is symmetric if $`d=\mathrm{flip}(d)`$. The following is an example of a symmetric diagram in $`A_6`$,
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In $`S_k`$ the symmetric diagrams are involutions.
###### Corollary 5.7.
A diagram $`dA_k`$ is symmetric if and only if $`d(P,P)`$.
###### Proof.
If $`d`$ is symmetric, then by Corollary 5.6 we must have $`P=Q`$. To prove the converse, let $`P=Q`$ and place these vacillating tableaux on the staircase border of a growth diagram. The local rules above are invertible: given $`\mu ,\nu `$ and $`\rho `$ one can the rules backwards to uniquely find $`\lambda `$ and determine whether there is an X in the box. Thus, the interior of the growth diagram is uniquely determined. By the symmetry of having $`P=Q`$ along the staircase, the growth diagram must have a symmetric interior and a symmetric placement of the Xs. This forces $`d`$ to be symmetric. ∎
This corollary tells us that the number of symmetric diagrams in $`A_k`$ is equal to the number of vacillating tableaux of length $`2k`$, or the number of paths to level $`k`$ in the Bratteli diagram of $`A_k(n)`$. Thus,
$$\mathrm{Card}(\{dA_k|d\text{ is symmetric }\})=\underset{\lambda \mathrm{\Lambda }_k}{}m_k^\lambda .$$
(5.5)
Furthermore, since our insertion restricts to the diagram subalgebras, (5.5) is true if we replace $`A_k`$ with any of the monoids $`S_k,B_k,T_k,I_t,R_k,PR_k`$ and replace $`m_k^\lambda `$ with the dimension of the appropriate irreducible representation. In the case of the symmetric group, (5.5) reduces to the fact that the number of involutions in $`S_k`$ equals $`_{\lambda k}f^\lambda `$. |
warning/0507/cond-mat0507070.html | ar5iv | text | # Entropically driven transition to a liquid-crystalline polymer globule
## Abstract
A self-consistent-field theory (SCFT) in the grand canonical ensemble formulation is used to study transitions in a helix-coil multiblock copolymer globule. The helices are modeled as stiff rods. In addition to the established coil-globule transition we show for the first time that, even without explicit rod-rod alignment interaction, the system undergoes a transition to a nematic liquid-crystalline (LC) globular state. The LC-globule formation is driven by the hydrophobic helical segment attraction and the anisotropy of the globule surface energy. The full phase diagram of the copolymer was calculated. It discriminates between an open chain, amorphous globule and LC-globule. This model provides a relatively simple example of the interplay between secondary and tertiary structures in homopolypeptides. Moreover, it gives a simple explanation for the formation of helix bundles in certain globular proteins.
The formation of secondary structure in proteins and in particular the $`\alpha `$-helix-coil transition in homopolypeptide chains is one of the well-investigated conformational transitions in biomolecules. In the well-known Zimm-Bragg (ZB)-theory of this phenomenon Zimm ; Poland the polypeptide molecule is considered as a one-dimensional cooperative system and the problem can be solved exactly. However, ZB-theory is only valid in a fully denaturated state when all three dimensional interactions are negligible. In real systems the helical parts are often hydrophobic, such that this hydrophobicity drives the helix-coil copolymer into a globular phase. Furthermore, helical parts can be seen as rigid rods and hence a rod-rod alignment interaction could be taken into account. Considering all these facts, it is natural to pose the question: how does the chain compaction affect the secondary and tertiary structure? This problem has been partially discussed within computer simulations of globular proteins Dill ; Socci . Some preliminary theoretical results concerning the formation of a LC-globule can be found in the review Grosberg . For a system of fixed composition of stiff parts (helices) and flexible parts, micelle formation has been studied using scaling considerations Nowak1 .
In this paper the phase behavior of a helix-coil copolymer is discussed by means of a self-consistent field theory (SCFT). In addition to two- and three-body contact interactions between the segments, each segment is also allowed to undergo a microscopic transition from flexible to stiff (i.e. being part of a helix). To describe this local helix-coil transition a grand canonical formalism was used. A formal derivation of the grand canonical SCFT free energy functional $`F[\phi (𝐫),\psi (𝐫,𝐮);\mu ,ϵ,\sigma ]`$ is given in Nowak2 . We introduced a field $`\phi (𝐫)`$ associated with a flexible segment at position $`𝐫`$ and a field $`\psi (𝐫,𝐮)`$ associated with a helical (rod-like) segment at position $`𝐫`$ with orientation $`𝐮`$. The local densities of flexible and helical segments are given by $`\rho _\mathrm{c}(𝐫)=1/2\phi (𝐫)^2`$ and $`\rho _\mathrm{h}(𝐫)=1/2d^2u\psi (𝐫,𝐮)^2`$ respectively. It is very difficult to deal with the full orientation dependence of $`\psi (𝐫,𝐮)`$. However, the helices are modeled as stiff rods without chirality, so we expect the solution for $`\phi (𝐫)`$ and $`\psi (𝐫,𝐮)`$ to have an uniaxial symmetry. Therefore we choose, without loss of generality, the $`z`$-axis as the preferred orientation and expand the orientation dependence in terms of Legendre polynomials. To lowest non-trivial order this expansion reads
$`\psi (𝐫,𝐮)\left({\displaystyle \frac{1}{4\pi }}\right)^{1/2}\psi _0(𝐫)+\left({\displaystyle \frac{5}{4\pi }}\right)^{1/2}\psi _2(𝐫)P_2(𝐮𝐧_z).`$ (1)
Using this expansion $`𝐮`$ can be integrated out and the SCFT free energy functional of our model is then given by the following form
$`F[\phi (𝐫),\psi _0(𝐫),\psi _2(𝐫);\mu ,ϵ,\sigma ]={\displaystyle \frac{\mu ϵ}{2}}{\displaystyle }d^3r[\psi _0^2(𝐫)+\psi _2^2(𝐫)]{\displaystyle \frac{1}{210\left(\mu ϵ\right)}}{\displaystyle }d^3r\{35\psi _0(𝐫)_r^2\psi _0(𝐫)`$
$`+14\sqrt{5}\psi _0(𝐫)[2_z^2_x^2_y^2]\psi _2(𝐫)+\psi _2(𝐫)[25_x^2+25_y^2+55_z^2]\psi _2(𝐫)\}`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle d^3r\phi (𝐫)\left[\mu \frac{a^2}{6}_r^2\right]\phi (𝐫)}+{\displaystyle \frac{\chi }{4}}{\displaystyle d^3r\left[\psi _0^2(𝐫)+\psi _2^2(𝐫)\right]^2}+{\displaystyle \frac{v}{8}}{\displaystyle d^3r\left[\phi ^2(𝐫)+\psi _0^2(𝐫)+\psi _2^2(𝐫)\right]^2}`$
$`+{\displaystyle \frac{w}{48}}{\displaystyle d^3r\left[\phi ^2(𝐫)+\psi _0^2(𝐫)+\psi _2^2(𝐫)\right]^3}2\sqrt{\pi \sigma }{\displaystyle d^3r\phi (𝐫)\psi _0(𝐫)}.`$ (2)
The numerical prefactors in Eq.(2) are due to the expansion of the orientation dependence - see Eq.(1) - and the subsequent integration over $`𝐮`$. The chemical potential $`\mu `$ is used to control the total number of segments of the polymer $`N`$. To model the energy gain due to the formation of hydrogen bonds in helices we introduced an energy gain per helical segment $`ϵ`$. The cooperativity effect in the formation of helices is taken into account by the cooperativity parameter $`\sigma `$. It can be regarded as a fugacity of the interfaces (or junction points) between a helix and a flexible part. $`\sigma =1`$ means no cooperativity effect and $`\sigma =0`$ reflects total cooperativity (i.e. the system can only form a fully flexible chain or one long helix). The interaction parameters $`v`$ and $`w`$ are global two- and three-body interaction constants between all segments, whilst $`\chi `$ controls the strength of a selective two-body interaction (due to hydrophobicity) between the helical segments only. Further technical details can be found in Nowak2 .
At this point it is important to emphasize that all interactions in this model are point contact interactions. We do not take into account any angle dependent interactions between the helical segments which explicitly favor alignment of two helices. The total number of segments $`N`$ is given by the sum of all segments in the helical and coil state: $`N=N_\mathrm{h}+N_\mathrm{c}`$, where
$`N_\mathrm{h}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^3r\left[\psi _0^2(𝐫)+\psi _2^2(𝐫)\right]}`$
$`N_\mathrm{c}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^3r\phi ^2(𝐫)}.`$ (3)
Minimization of the functional in Eq.(2) with respect to $`\phi `$, $`\psi _0`$ and $`\psi _2`$ yields a set of three coupled differential equations for the three self-consistent fields. These differential equations are highly non-linear and can only be solved numerically. This was done using the finite element toolkit Gascoigne gascoigne .
The system shows a collapse transition from an open chain to a dense globule. In this letter we always consider finite systems, therefore this transition is a crossover transition with a finite broadness. Hence we have to define a point during the crossover as the transition point. Before we explain how we do this for the coil-globule transition it is necessary to clarify how we generally deal with the chemical potential. In the grand canonical ensemble the number of particles, here the total number of segments of the polymer $`N`$, is not fixed but its mean value is determined by equilibrium conditions. However, in a real experiment the helix-coil copolymer has a fixed length. In order to ensure this fixed length $`N`$ we tune the chemical potential $`\mu `$ for each set of physical parameters ($`v,w,\chi ,ϵ,\sigma `$) such that the equilibrium value of $`N`$, calculated by performing the integrations in Eq.(3) is equal to the desired one. For a given set of parameters $`N(\mu )`$ can be computed and a typical example of this is shown in Fig.1.
For $`\mu 0`$ the total number of segments $`N`$ diverges. This corresponds to the $`N\mu ^1`$ behavior of a $`\mathrm{\Theta }`$-solvent chain. The divergence of $`N`$ at a specific value $`\mu `$ on the right hand side of the minimum corresponds to a fully collapsed infinite globule. The minimum of $`N(\nu )`$ is naturally associated with the transition point. We always fix $`N`$ for a given set of parameters by tuning $`\mu `$, so from a plot like the one shown in Fig.1, we can distinguish whether the system is left of the transition point (i.e. in the open chain regime) or right of the transition point (i.e. in the globular regime). The transition point is now defined as the set of parameters at which the chosen fixed $`N`$ is equal to the minimum of $`N(\mu )`$. Both curves in Fig.(1) have their minimum at $`N=10^5`$, the value which we have chosen for the results presented below. The minimum of the dotted curve corresponds to the first triangle ($`ϵ=0.08`$, $`\chi =0`$) in the phase diagram Fig.(5) and the minimum of the continuous curve to the second triangle ($`ϵ=0.1`$, $`\chi =0.0138`$). It is pertinent to note that the characteristic dependence of $`N`$ on $`\mu `$ for a simple homopolymer globule has first been discussed by Kholodenko and Freed Kholod . For all the results presented below, the global three-body interaction constant is set to $`w=1.0`$, the global two-body interaction constant to $`v=0.2`$ and the cooperativity parameter to $`\sigma =10^4`$. The SCFT formalism is a reasonable approximation only for systems in a globular state or close to a globular state. Therefore $`v`$ has to be negative in order to ensure that this is always the case. Although we do not use the SCFT treatment of our model to describe an open chain of the helix-coil copolymer, it is perfectly possible with the definition given above to calculate the transition line between the globular state and the open chain state.
The really astonishing feature of this model is that it also shows a crossover transition from a disordered amorphous globule with low or mid fraction of helical segments $`\mathrm{\Theta }_\mathrm{h}=N_\mathrm{h}/N`$ to an ordered liquid-crystalline globule with very high fraction of helical segments. This transition occurs despite the fact that there is no explicit angle dependent alignment interactions. The transition is triggered by a subtle interplay of the entropy contribution (surface energy), represented by the derivative terms in Eq.(2), and bulk interaction energy, represented by the $`\chi `$-term. It is well known Grosberg1 that in a simple homopolymer globule the surface energy has an entropic nature (since the conformational set of surface segments is constrained) and is isotropic. As one can see from Eq. (2) in our case the surface energy is anisotropic, so that after a proper inspection Subseq one can ensure that the surface tension in the $`xy`$-direction is smaller than the one in the $`z`$-direction. That is why the system tries to maximize its lateral surface in $`xy`$-directions and minimize it in $`z`$-direction, i.e. a nematic, cigar shaped, LC-globule occurs.
To measure orientational order in the system we define the nematic order parameter $`S`$ as follows (see e.g. Prost )
$`S`$ $``$ $`{\displaystyle \frac{1}{3N}}{\displaystyle d^3rd^2uP_2(\mathrm{cos}\theta )\psi ^2(𝐫,𝐮)}`$
$`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle d^3r\left(\frac{2}{\sqrt{5}}\psi _2(𝐫)\left[\psi _0(𝐫)+\sqrt{5}\psi _2(𝐫)\right]\right)}.`$
Fig.(2) shows the increase in fraction of helical segments $`\mathrm{\Theta }_\mathrm{h}`$ with $`|\chi |`$ during the transition and Fig.(3) shows the simultaneous onset of a finite nematic order parameter $`S`$.
The onset of the transition is shifted to lower values of $`|\chi |`$ with increasing energy gain per helical segments $`ϵ`$. This is due to an increase of bulk interaction energy for fixed $`\chi `$ with increasing number of helical segments. The curves for $`ϵ=0.2`$ and $`ϵ=0.15`$ in Figs.(2,3) start at non-zero values of $`\chi `$. These values of $`\chi `$ correspond to the transition point between open chain state and globular state, i.e. the triangles in Fig.(5). The transition point for the crossover transition from an amorphous to a nematic LC-globule has to be defined in a reasonable way. We define these transition points as the inflection points of the $`S(\chi )`$-curves in Fig.(3).
To demonstrate how the shape of the globule changes during the transition we show in Fig.(4) a color-coded plot of the local density in $`\varrho `$-$`z`$ space, where $`\varrho =\sqrt{x^2+y^2}`$. The center of the globule is in the bottom left corner. $`\varrho `$ is increasing from left to right and $`z`$ is increasing from bottom to top. Red means high density and dark blue zero density.
The top two pictures show the density profile at the transition point between open chain and amorphous globule. At this point the system is spherical and has a very broad surface layer of decaying density resembling the characteristics of a random walk. Although the density of the helical segments shown on the right is higher than the density of the flexible segments shown on the left, their distribution and the shape of the profile is very similar. The middle two pictures show the density profile at the transition point between amorphous globule and LC-globule. The system adopts a slightly cylindrical shape indicating the onset of nematic order. It can also be seen that the density maximum of the flexible segments is not in the center of the globule denoting a repulsion of flexible segments from the center to the surface layer. The surface layer is now much narrower. The bottom two pictures show the density profile deep in the nematic LC-globule phase. The globule has developed a strongly asymmetric cylindrical shape indicating strong nematic order. The repulsion of flexible segments from the center towards the surface layer can be seen clearly and the surface layer is now very narrow.
The results shown in Fig.(3) allow us to compute a complete phase diagram of a helix-coil copolymer in $`ϵ\chi `$ space, see Fig.(5).
The triangles are the transition points between open chain and globule. The squares are the transition points between amorphous globule and LC-globule. Note, the points plotted in the phase diagram, Fig.(5), are what we defined above as the transition points of rather broad crossover transitions. Therefore the boundaries in the phase diagram have to be understood as ”center lines” of broader regions in which the crossover from one phase to the other occurs. Although the qualitative shape of the phase diagram stays the same, the position of the transition lines changes with $`N`$, $`\sigma `$ and $`v`$. If $`|v|`$ is increased the open chain region of the phase diagram becomes smaller (and disappears eventually). With decreasing $`\sigma `$ the transition from a disordered to an ordered globule becomes sharper and the transition line between these two regions is shifted to smaller values of $`|\chi |`$. Since this transition occurs due to an interplay between surface energy and bulk interaction energy, the transition also becomes sharper for decreasing system size $`N`$, which, at first sight, is a rather unusual and surprising behavior. For smaller systems the surface energy plays a bigger role and therefore leads to a sharper transition. For $`N\mathrm{}`$ the ordered globule phase finally disappears, since the surface contributions to the free energy vanish for infinite systems. A detailed discussion of the modification of the phase diagram will be given in a subsequent publication Subseq
In summary, we presented a SCFT for a multiblock helix-coil copolymer chain based on grand canonical ensemble considerations. The minimization of the corresponding free energy functional - Eq.(2) - has been done numerically and it has been shown that three phase states can be clearly seen: open helix-coil chain, amorphous globule and nematic LC-globule. It is a novel result that the formation of a LC-globule occurs without explicit alignment interactions between the helical parts. It is the entropical surface tension anisotropy which drives the globule in the nematic LC-state in order to maximize the density of the hydrophobic helical segments. In the presence of an explicit attractive alignment interaction (which can be taken into account by a Maier-Saupe term) one only sees an enhancement of this effect Subseq . We believe that this model provides a relatively simple example of the interplay between secondary and tertiary structure in homopolypeptides. It can also give a simple explanation for the formation of helix bundles in certain globular proteins. Both, simulations Zhou and experiments Ptitsyn show that proteins can adopt not only the native state and completely denaturated state (open chain) but also so-called premolten and molten globular states. For helix-bundle proteins the premolten globule, which does not show any order of the helices, corresponds to our amorphous globule. The molten globule with ordered helices but without native contacts corresponds to our LC-globule. The chain length of proteins is typically of the order of $`N10^210^3`$ instead of $`N10^5`$. However, as mentioned above the crossover from an amorphous globule to a nematic LC-globule with aligned helices becomes even sharper for shorter chains Subseq . This indicates that our model indeed provides a simple explanation for the formation of helix bundles.
C. Nowak would like to thank T. Richter for invaluable help on how to use Gascoigne and D. Andrienko for fruitful discussions. V.G. Rostiashvili and T.A. Vilgis appreciate the financial support of the German Science Foundation (SFB 625). |
warning/0507/hep-ph0507192.html | ar5iv | text | # CERN-PH-TH/2005-131UAB-FT-584 Split extended supersymmetry
## Abstract
We show how splitting supersymmetry reconciles a class of intersecting brane models with unification. The gauge sector in these models arises in multiplets of extended supersymmetry while matter states are in $`N=1`$ representations. A deformation of the angles between the branes gives large masses to squarks and sleptons, as well as supersymmetry breaking contributions to other string states. The latter generate at one-loop heavy Dirac masses for Winos and gluinos and can induce a mass term for the Higgsino doublets. We find that this scenario is compatible with gauge coupling unification at high scale for both cases where the gauge sector is $`N=2`$ and $`N=4`$ supersymmetric. Moreover a neutralino, combination of neutral Higgsinos and Binos, is a natural candidate for dark matter.
INTRODUCTION The necessity of a Dark Matter (DM) candidate and the fact that LEP data favor the unification of the three Standard Model (SM) gauge couplings are smoking guns for the presence of new physics at high energies. The latter can take the form of supersymmetry which, if broken at low energies, offers a framework for solving the gauge hierarchy problem. Supersymmetry is also welcome as it naturally arises in string theory, which provides a framework for incorporating the gravitational interaction in our quantum picture of the universe. Recently strong evidence has been accumulated for the presence of a tiny dark energy in the universe. This raises another hierarchy problem which is not solved by any known symmetry. It leads to reconsider our notion of naturalness, possibly also affecting our view of mass hierarchy. It was then proposed to consider that supersymmetry might be broken at high energies without solving the gauge hierarchy problem. More precisely, making squarks and sleptons heavy does not spoil unification and the existence of a DM candidate while at the same time it gets rid of all unwanted features of the supersymmetric SM related to its complicated scalar sector. On the other hand, experimental hints to the existence of supersymmetry persist since there are still gauginos and Higgsinos at the electroweak (EW) scale. This is the so-called split supersymmetry framework . Implementing this idea in string theory is straightforward . However one faces a generic problem: in simple brane constructions the gauge sector comes in multiplets of extended supersymmetry . In this work we show that these economical string-inspired brane models allow for unification of gauge couplings at scales safe from proton decay problems. Moreover they provide us with a natural DM candidate.
BRANE MODELS The simplest supersymmetric brane models are obtained as compactification on a six-dimensional torus which is a product of three factorized tori. The states in these models can be assembled into three sets. The first set is made of strings with both ends on the same stack of branes, leading to $`N=4`$ vector multiplets. In the following we will also consider a departure from this minimal case, when two chiral adjoint $`N=1`$ multiplets are projected out to remain with $`N=2`$ vector supermultiplets<sup>2</sup><sup>2</sup>2Note that there are also brane constructions with no extended supersymmetry in the gauge sector based on non-toroidal compactifications. The second set contains strings which stretch between two stacks of branes that intersect only in two out of the three internal tori, giving rise to $`N=2`$ hypermultiplets. The last set contains strings which are localized at intersection points of branes in all the three tori and the associated states form $`N=1`$ supermultiplets. In this setup of supersymmetric limit, the SM states are identified as follows:
* Gauge bosons emerge as massless modes of open strings with both ends on the same stack of coincident branes. They arise in $`N=2`$ or $`N=4`$ supermultiplets which are decomposed, for each gauge group factor $`G_a`$, into one $`N=1`$ vector superfield $`W_a`$ and one or three chiral adjoint superfields $`A_a`$, respectively.
* Quarks and leptons are identified with massless modes of open strings localized at point-like brane intersections and belong to $`N=1`$ chiral multiplets.
* Pairs of Higgs doublets originate from $`N=2`$ supersymmetry preserving intersections. They are localized in two tori where branes intersect, while they propagate freely in the third torus where the two brane stacks are parallel. Here we assume that all possible additional non-chiral states that may appear in generic string constructions can be made superheavy.
Supersymmetry breaking is then achieved by deforming brane intersections with a small angle $`\mathrm{\Theta }`$. As a result, a $`D`$-term with $`D=\mathrm{\Theta }M_S^2`$ appears, associated to a corresponding magnetized $`U(1)`$ factor with superfield strength $`𝒲`$. Here, $`M_S`$ is the string scale. Supersymmetry is then broken and soft masses are induced:
* A tree-level mass $`m_0\sqrt{\mathrm{\Theta }}M_S`$ for squarks and sleptons localized at the deformed intersections. All other scalars acquire in general high masses of order $`m_0`$ by one loop radiative corrections. Appropriate fine-tuning is needed in the Higgs sector to keep light $`n_H`$ doublets.
* A Dirac mass is induced through the dimension-five operator
$$\frac{a}{M_S}d^2\theta 𝒲W^aA_am_{1/2}^Da\frac{m_0^2}{M_S},$$
(1)
where $`a`$ accounts for a possible loop factor. Actually, this operator arises quite generally at one-loop level in intersecting D-brane models with a coupling that depends only on the massless (topological) sector of the theory . Note that this mass does not break $`R`$-symmetry and provides an answer to a problem of split supersymmetry related to the mechanism of generating gaugino masses.
UNIFICATION We now study the compatibility of this framework with one-loop unification. In the energy regime between the unification scale $`M_{GUT}`$ and the EW scale $`M_W`$, the renormalization group equations meet three thresholds. From $`M_{GUT}`$ to the common scalar mass $`m_0`$ all charged states contribute. Below $`m_0`$ squarks and sleptons (which do not affect unification), adjoint scalars and $`2n_H`$ Higgses decouple, while below $`m_{1/2}^D`$ the $`N=2`$ or $`N=4`$ gluinos and Winos drop out. Finally, at TeV energies Higgsinos (and maybe the Binos) decouple and we are left at low energies with the Standard Model with $`n_H`$ Higgs doublets. For the purpose of this computation we use $`M_SM_{GUT}`$ and we vary $`a`$ between $`a=1`$ and $`a=1/100`$. Realistic values for $`M_{GUT}`$ and $`m_0`$ are obtained in both $`N=4`$ and $`N=2`$ cases. The results are summarized in Table 1.
Notice that we have imposed perfect unification at one-loop. Although this is not necessary from the string theory point of view, it is one of the main motivations of split supersymmetry and can be imposed along the lines of . In this case, for $`N=4`$ there is *no* solution with $`n_H=2`$ Higgses at low energies, whereas for $`N=2`$ we find a solution with either $`n_H=1`$ or $`n_H=2`$. In all cases the unification scale is high enough to avoid problems with proton decay. For the two possible cases with one light Higgs ($`N=2`$ or $`N=4`$), $`M_{GUT}`$ is very close to the Planck scale so that there should be no need to explain the usual mismatch between these two scales. Varying the loop factor $`a`$ from 1 to $`1/100`$ amounts to an increase by one order of magnitude in the value of $`m_0`$, but $`M_{GUT}`$ and $`m_{1/2}^D`$ remain stable within $`𝒪(1)`$ factors.
The low energy sector of these models contains, besides the SM, just some fermion doublets (Higgsinos) and eventually two singlets (the Binos from the discussion below). It therefore illustrates the fact that only these states are needed for a minimal extension of the SM consistent with unification and DM candidates, and not the full fermion spectrum of split supersymmetry. A similar observation was recently done in the literature , without the presence of Dirac gauginos associated to the intermediate scale $`m_{1/2}^D`$. Here, we do not find these solutions because they do not unify at one-loop.
DARK MATTER Another constraint on the models is that they must provide a DM candidate. As usually in supersymmetric theories this should be the lightest neutralino. Pure Higgsinos $`(\stackrel{~}{H}_1,\overline{\stackrel{~}{H}}_2)^T`$ cannot be DM candidates because their mass is of Dirac type: $`\mu \stackrel{~}{H}_1\stackrel{~}{H}_2+h.c.`$ Quasi-Dirac Higgsinos would interact inelastically with matter via vector-like couplings as $`i(\overline{\stackrel{~}{H}}_{}\overline{\sigma }_\mu \stackrel{~}{H}_+\overline{\stackrel{~}{H}}_+\overline{\sigma }_\mu \stackrel{~}{H}_{})`$ where $`\stackrel{~}{H}_\pm \stackrel{~}{H}_1\pm \stackrel{~}{H}_2`$ are the mass eigenstates with mass eigenvalues $`\mu \pm ϵ`$ respectively. In Dark Matter direct detection experiments $`\stackrel{~}{H}_{}`$ can only scatter inelastically off of a nucleus of mass $`m_N`$ by transitioning to $`\stackrel{~}{H}_+`$ if $`ϵ<ϵ_0\frac{1}{2}\beta ^2m_N`$ . For Ge experiments with $`m_N=73`$ GeV and taking a typical escape Dark Matter particle velocity $`\beta c600`$ km/s one obtains $`ϵ_0146`$ keV. Since direct detection experiments have ruled out Dirac fermions up to masses of order $`50`$ TeV some mixing coming from the Binos is required in order to break the degeneracies of the two lightest neutralinos and provide a mass difference $`ϵ>ϵ_0`$ preventing inelastic scattering off the nucleus. The mass difference $`ϵ>ϵ_0`$ translates into an upper bound on the Dirac gaugino mass of about $`10^5`$ GeV, for the required Higgsino mass splitting to be generated through the EW symmetry breaking mixing (of order $`m_W^2/m_{1/2}^D`$) described below. This value compared to the values in Table 1 leads to the $`N=2`$, $`n_H=1`$ case as the only possibility to accommodate it <sup>3</sup><sup>3</sup>3In this case and to destroy the Dirac nature of mass eigenstates the Higgs sector should be arranged (as quarks and leptons do) in $`N=1`$ multiplets.. In fact one could have an order of magnitude suppression of the induced Dirac mass for Binos relative to the other gauginos, which is not unreasonable to assume in brane constructions.
In the other two models, the required suppression factor is much higher and the above mechanism would be very unnatural. However, since Binos play no rôle for unification as they carry no SM charge, we could imagine a scenario where $`m_{1/2}^D`$ vanishes identically for Binos, but not for the other gauginos. For instance consider the case where Dirac masses from the operator (1) are generated by loop diagrams involving $`N=2`$ hypermultiplets with supersymmetric masses of order $`M_{GUT}`$ and a supersymmetry breaking splitting of order $`\mathrm{\Theta }`$. It is then possible to choose these massive states such that they carry no hypercharge. In that case Binos can only have Majorana masses. Such masses are in general induced by a dimension-seven effective operator, generated at two-loop level :
$$\frac{b^2}{M_S^3}d^2\theta 𝒲^2\mathrm{Tr}W^2m_{1/2}^Mb^2\frac{m_0^4}{M_S^3},$$
(2)
where $`b`$ is another loop factor. Putting numbers in the above formula, we get for the $`N=4`$ $`n_H=1`$ model $`m_{1/2}^M5\times 10^6`$ GeV, which is the upper bound for DM with Majorana Bino mass <sup>4</sup><sup>4</sup>4This higher value for the upper bound on Majorana ($`5\times 10^6`$ GeV) versus Dirac Bino mass ($`10^5`$ GeV) we used before is due to the fact that the induced Higgsino mass splitting through EW symmetry breaking mixing is further suppressed by the weak angle in the Dirac case.. On the other hand, for the $`N=2`$ $`n_H=2`$ model, we get $`m_{1/2}^M100`$ GeV which is obviously acceptable. Finally, for the $`N=2`$ $`n_H=1`$ case, $`m_{1/2}^M10`$ keV which does not play any rôle if there is also a Dirac mass, as we assume and discussed above. Thus, the constraint of a viable DM candidate leaves us with two possibilities: (a) $`N=2`$ with $`n_H=1`$ and Dirac masses for all gauginos and (b) $`N=4`$ with $`n_H=1`$, or $`N=2`$ with $`n_H=2`$ and Majorana mass Binos.
HIGGSINO MASSES The Higgsinos themselves must acquire a mass of order the EW scale. This is induced by the following dimension-seven operator, generated at one loop level :
$$\frac{c}{M_S^3}d^2\theta 𝒲^2\overline{D}^2\overline{H}_1\overline{H}_2\mu c\frac{m_0^4}{M_S^3},$$
(3)
where $`c`$ is again a loop factor. The resulting numerical value is of the same order as $`m_{1/2}^M`$ of Eq. (2). Thus, such an operator can only give a sensible value of $`\mu `$ for the $`N=2`$ $`n_H=2`$ model. In the other two cases, $`N=4`$ or $`N=2`$ with $`n_H=1`$, $`\mu `$ remains an independent parameter.
Another constraint on these models may come from the life-time of the extra states. We first consider the case of $`N=2`$. Scalars can decay into gauginos, Dirac gluinos decay through squark loops sufficiently fast and Dirac Winos and Binos decay into Higgses and Higgsinos. On the other hand, for $`n_H=2`$, there are two Majorana Binos at low energies. However, as mentioned above, in the models we consider the two Higgs doublets form an $`N=2`$ hypermultiplet. It follows that there is an $`N=2`$ coupling between both Binos with Higgses and Higgsinos, implying that the only stable particle is the usual lightest sparticle (LSP). Finally, in the $`N=4`$ model, scalars still decay into gauginos, but we have now *two* Dirac gluinos, Winos and Binos; half of them decay as before, either through scalar loops or into Higgs-Higgsinos, while the other half can only decay through string massive states. Their lifetime is then estimated by $`\tau \left(M_S/10^{13}\text{GeV}\right)^4\left(10^2\text{GeV}/m_{\stackrel{~}{g}}\right)^5\tau _U`$, where $`m_{\stackrel{~}{g}}`$ is the gaugino mass and $`\tau _U`$ is the lifetime of the universe. For gluinos and Winos there is no problem, but Binos are very long lived although still safe, with a life-time of order $`\tau _U/10`$.
To summarize, at low energies we end up with two distinct scenarios after all massive particles are decoupled: i) $`n_H=1`$ with light Higgsinos (models with $`N=2`$ and $`N=4`$ gauge sector and $`n_H=1`$), and; ii) $`n_H=2`$ with light Higgsinos and Binos (model with $`N=2`$ gauge sector and $`n_H=2`$). In the $`n_H=1`$ scenario the DM candidate is mainly Higgsino, although the much heavier Bino is light enough to forbid any vector couplings. The relic density can be estimated (using for example the DarkSUSY program ) and reproduces the actual WMAP results for $`\mu 1.1`$ TeV.
The $`n_H=2`$ scenario is more interesting since there are more particles at low energies. The $`N=2`$ coupling between Binos-Higgs-Higgsinos leads, after EW symmetry breaking to the neutralino mass matrix,
$$\left(\begin{array}{cccc}M& 0& m_zs_wc_\beta & m_zs_ws_\beta \\ 0& M& m_zs_ws_\beta & m_zs_wc_\beta \\ m_zs_wc_\beta & m_zs_ws_\beta & 0& \mu \\ m_zs_ws_\beta & m_zs_wc_\beta & \mu & 0\end{array}\right)$$
in the basis $`(\stackrel{~}{B}_1,\stackrel{~}{B}_2,\stackrel{~}{H}_1,\stackrel{~}{H}_2)`$ and where $`M=m_{1/2}^M`$ stands for the Bino Majorana masses. The mass matrix can be diagonalized to obtain:
$$m_\chi =1/2\left[(M+ϵ_1\mu )ϵ_2\sqrt{(Mϵ_1\mu )^2+4m_z^2s_w^2}\right]$$
(4)
where the four different mass eigenvalues are labeled by $`ϵ_{1,2}=\pm 1`$.
The values of $`\mu `$ and $`M`$ can then be chosen in such a way as to reproduce the WMAP observed DM density, as previously studied . There are then three cases:
* $`M\mu `$: the Bino does not interact strongly enough to annihilate and will in general overclose the universe.
* $`M\mu `$: this model converges to the $`n_H=1`$ scenario and WMAP results require $`\mu 1.1`$ TeV.
* $`M\mu `$: the lightest neutralino ($`\chi `$) is in general a mixture of Higgsinos and Binos and is a natural candidate for DM. Low values of $`\mu `$ are now possible.
COLLIDER PHENOMENOLOGY In the $`n_H=1`$ scenario with $`\mu 1.1`$ TeV, the heavy spectrum will be hardly observable at LHC while a Linear Collider with center of mass energy of around 2.5 TeV will be needed to detect a possible signature.
In the $`n_H=2`$ scenario, the main collider signature is through the production of charginos. Their mass is given by $`m_{\chi ^\pm }=\mu +\delta \mu `$, where $`\delta \mu `$ is due to electromagnetic contributions and is of order $`300`$ to $`400`$ MeV. The produced charginos will decay into the neutralino, mainly through emission of a virtual $`W^\pm `$ which gives rise to lepton pairs or pions depending on its energy. This decay is governed by the mass difference $`\mathrm{\Delta }m_\chi =m_{\chi ^\pm }m_{\chi ^0}`$, which is a function of the two parameters $`M`$ and $`\mu `$. Because charginos are produced through EW processes, LHC will mainly be able to explore the case of very light charginos, which exist only in the limited area of the parameter space with $`M\mu `$. Unlike in low energy supersymmetry, the absence of cascade decays in this case will make it difficult to separate the signal from similar events produced by Standard Model $`W^\pm `$ production processes.
The search at future $`e^+e^{}`$ colliders is more promising, and can be discussed either as a function of the model parameters $`(M,\mu )`$, or as a function of the low energy observables $`(m_{\chi ^\pm },\mathrm{\Delta }m_\chi )`$. For most of the $`(M,\mu )`$ parameter range, $`\mathrm{\Delta }m_\chi `$ is small, at most of order a few GeV. Because the value of $`\delta \mu `$ is not small enough to make the chargino long-lived as to produce visible tracks in the vertex detectors, we have to rely on its decay products. This degeneracy implies that the produced leptons or pions are very soft and it would typically be difficult to disentangle them from the background due to emission of photons from the beam. The strategy is then to look for $`e^+e^{}\gamma +E_T/`$. A proper cut on the transverse momentum of the photon allows to eliminate the background of missing energy due to emission of $`e^+e^{}`$ pairs along the beam, as the conservation of transverse momentum implies now a simultaneous detection of electrons or positrons . The best possible scenario is when $`M`$ and $`\mu `$ are of the same order since, as soon as $`M`$ starts to be greater than $`\mu `$, the Binos quickly decouple and this model converges to the $`n_H=1`$ scenario with $`\mu 1.1`$ TeV. The Higgs sector of these models will also give a signal at LHC. The case with just one Higgs at low energies predicts a mass $`160`$ GeV depending on the exact value of $`m_0`$. In the case of two Higgses, one of them will have this bound whereas the mass of the other doublet is controlled by $`m_A`$ as in the MSSM. The result is that one Higgs will always be discovered at LHC whereas the others will depend on the value of $`m_A`$, which is a free parameter. In any case it will be imposible to distinguish these models from another type of split supersymmetry or from a general two Higgs model just at LHC. A careful measurement of Higgs coupling will be needed to disentangle between the different models that could be achieved at the ILC.
CONCLUSION Before closing, we would like to make a few comments concerning the number of parameters and fine-tuning issues. Dirac gaugino masses, and consequently scalar masses related through Eq. (1), are fixed by the one-loop unification condition according to Table 1. In the two $`n_H=1`$ models (with $`N=4`$ or $`N=2`$ gauge sector) the Higgsino mass $`\mu `$ is fixed by the DM constraint to $`1.1`$ TeV. Since the supersymmetry breaking scale $`m_0`$ is high, a fine-tuning is needed like in split supersymmetry to keep one Higgs scalar light. On the other hand, in the $`n_H=2`$ model with $`N=2`$, $`\mu `$ is determined to the right scale by the effective operator (3) and the required Majorana component of Bino mass by (2). Obviously in this case one needs more fine-tuning conditions in order to keep both eigenvalues of the Higgs mass-squared matrix light. One may wonder that for a two-by-two symmetric matrix, this implies three conditions on all its elements. However, it is easy to see that only two conditions are really needed. The reason is that the off-diagonal element is protected by two global low energy symmetries, namely Pecei-Quinn and $`R`$-invariance. As a result, if its tree-level value vanishes, quantum contributions will be proportional to $`\mu m_{1/2}^M𝒪(\mathrm{TeV})^2`$.
In summary, we presented three viable models of high scale supersymmetry that naturally emerge in simple string constructions with intersecting branes and are compatible with gauge coupling unification and the existence of a dark matter candidate. Gauginos come in multiplets of extended supersymmetry and get high scale Dirac masses by dimension-five effective operators without breaking $`R`$-symmetry, consistently with gauge coupling unification. The low energy sector contains besides the Standard Model particle content a Higgsino pair providing, in general through mixing with the Bino, a natural Dark Matter candidate.
This work was supported in part by the European Commission under the RTN contract MRTN-CT-2004-503369, in part by CICYT, Spain, under contracts FPA2001-1806 and FPA-2002-00748, in part by the INTAS contract 03-51-6346 and in part by IN2P3-CICYT under contract Pth 03-1. |
warning/0507/quant-ph0507093.html | ar5iv | text | # Quantum optics in different representations of the algebra of canonical commutation relations (I): Unexpected properties of Rabi oscillations — theory and experiment
## I Introduction
It is known that different representations of the same Lie algebra correspond to different types of physical systems. Parameters that characterize the representations have a meaning of quantum numbers (angular momentum for the rotation group, mass and spin for the Poincaré group, and so on). In principle, if we find a new representation of some physically meaningful group or algebra we should seriously consider the possibility that the representation corresponds to a physical system, perhaps yet unknown. This is how the anyons were predicted anyons .
In quantum optics the central role is played by the Lie algebra of canonical commutation relations (CCR)
$`[a(𝒌),a(𝒌^{})^{}]`$ $`=`$ $`\delta (𝒌,𝒌^{})I(𝒌).`$ (1)
The element $`I(𝒌)`$ commutes with all the elements of the algebra and, by the Schur lemma, is proportional to the identity operator if the representation is irreducible. Fields with different boundary conditions correspond to different $`I(𝒌)`$ and $`\delta (𝒌,𝒌^{})`$, and thus to different representations of CCR.
Still, the problem seems deeper than that. First of all, the theorem of von Neumann vN states that there exists an infinite number of inequivalent irreducible representations of CCR for systems with infinite numbers of degrees of freedom (here corresponding to the infinite number of different wave vectors). It is not clear which representation to choose, even if one keeps boundary conditions fixed and restricts the analysis to free fields. Secondly, as stressed by Dirac in his last published paper bib:dirac2 , one should not ignore the physical potential inherent in reducible representations. These are unavoidable in quantum physics since tensor products of irreducible representations are not themselves irreducible. Composite quantum systems are inherently reducible.
One of the representations that gained particular popularity in quantum optics employs the Hilbert space of infinitely many harmonic oscillators. To each frequency of the field there corresponds a separate oscillator. Physically the representation is rather pathological, just to mention the problem with the infinite energy of the ground state. Mathematically, the representation is also pathological in the sense of involving a non-separable Hilbert space typical of infinite tensor products vN2 , a fact rarely mentioned in the quantum optics literature.
It is relevant to mention in this context that recently, in a series of papers 1 ; 2 ; 3 ; 4 one of us investigated the possibility of electromagnetic field quantization in terms of certain reducible representations of CCR. The main feature of the proposed formalism was that there was no link between the number of frequencies allowed by field boundary conditions and the number $`N`$ of oscillators used in construction of the Hilbert space. In particular, it made perfect sense to speak of fields modelled by a finite number of oscillators even if the number of frequencies in the field was infinite. And vice versa, a monochromatic field could be modelled with any $`N`$.
It is easy to understand why this is possible if one thinks of each oscillator as a wave packet involving all the possible frequencies. The parameter $`N`$ occurring in the representation has a status of a quantum number, and the vacuum is a Bose-Einstein condensate of $`N`$ oscillators at zero temperature. The limit $`N\mathrm{}`$ plays a role of a correspondence principle mapping the new formalism into the standard one. What is interesting, when it comes to computing averages it turns out that the wave function of the ground state plays a role of a cut-off function regularizing integrals at ultraviolet and infrared regimes. The new formalism thus automatically introduces many elements that are imposed in an ad hoc manner in the standard one.
However, if in reality $`N`$ is large but finite then quantum field theory we know nowadays is the $`N=\mathrm{}`$ approximation of a more fundamental, and hopefully more consistent, $`N<\mathrm{}`$ theory. The question is: Why the $`N=\mathrm{}`$ approximation works so well, and where to look for experimental manifestations of a finite $`N`$?
In this paper we give partial answers to both parts of the question. We concentrate on the Jaynes-Cummings model where we solve Heisenberg equations of motion for the atomic inversion operator. A difference with respect to the usual treatment of the problem is that we begin with a representation-independent level. The first conclusion we find at that stage is that all the irreducible representations of (1) produce physically equivalent results.
As the next step we choose the reducible representation parametrized by $`N<\mathrm{}`$ and compute evolution of atomic inversion with different vacuum, thermal, and coherent-state initial conditions.
Here the situation changes. But first of all it has to be stressed that in the limit $`N\mathrm{}`$ we recover the standard formulas. This fact is of crucial importance for the whole approach since it plays a role of a correspondence principle and guarantees that the new theory may be regarded as a generalization of the one based on irreducible representations.
For finite $`N`$ there are differences which we compare with experimental data of Haroche . What is interesting, it seems that the data are more consistent with a finite $`N`$ than $`N=\mathrm{}`$ ($`N=280`$ for the nearly vacuum state at $`T=0.8`$ K, $`N=420`$ for the coherent state with $`0.4`$ and $`0.85`$ photons in average). The main reason why the experiment does not fully support the standard theory is that for finite $`N`$ one expects a faster relaxation of Rabi oscillation than what one might expect on the basis of $`N=\mathrm{}`$, and this is precisely what seems to happen in the experiment.
The relaxation occurring for finite $`N`$ is not a decay but a beat and therefore waiting sufficiently long one should see a revival of the Rabi oscillation. We show that cavities with lifetimes of a few hundred $`\mu `$s in principle allow for observation of the revival. Of particular interest are the maser and mazer experimental setups Walther1 ; Walther2 ; Walther3 ; Lamb , but we leave it for a future work.
If our intuitions are correct and the finite $`N`$ representations are more physical than the limiting $`N=\mathrm{}`$ case, then many additional questions have yet to be answered. In the paper that accompanies the present one we address the issues of Lorentz and gauge covariance and test the formalism on another exactly solvable model: Quantum fields produced by a classical current.
The present paper is organized as follows. In Sec. II we discuss the Jaynes-Cummings model at a representation independent level. In Sec. III we assume that the representation is irreducible and conclude that all such representations imply equivalent physics. In Sec. IV we switch to the reducible representation involving $`N<\mathrm{}`$ oscillators. Then, in Sec. V, we write the explicit form of the inversion operator $`R_3(t)`$ in this representation and, in Sec. VI, compute atomic inversion with different initial conditions. In particular, we discuss in detail the limit $`N\mathrm{}`$ and its links to the law of large numbers and renormalized parameters. Some technicalities and remarks about the reducible representation are moved into Appendices.
## II Representation independent formulation of the Jaynes-Cummings model
The Jaynes-Cummings model JC ; Allen represents a two-level atom interacting with a single mode of electromagnetic field in a cavity. The cavity boundary conditions imply that the set of free-field momenta is discrete and it is convenient to work from the outset with the discrete notation
$`[a_j,a_k^{}]=\delta _{jk}I_j.`$ (2)
Here $`\delta _{jk}`$ is the Kronecker delta, the asterisk denotes the adjoint in the sense of -algebras, and $`I_j^{}=I_j`$. We assume there exists a free-field Hamiltonian $`H_0`$ satisfying
$`[a_j,H_0]`$ $`=`$ $`\omega _ja_j,`$ (3)
$`[a_j^{},H_0]`$ $`=`$ $`\omega _ja_j^{}.`$ (4)
Note that $`H_0`$ cannot, in general, be given by $`_j\omega _ja_j^{}a_j`$; the latter works only if $`I_j`$ is an identity, which we do not assume at the present stage.
Let us now select a frequency $`\omega _{j_0}=\omega `$ and assume that only this frequency couples to the two-level system. The corresponding CCR operators will be indexed by $`\omega `$, that is: $`a_{j_0}=a_\omega `$, $`a_{j_0}^{}=a_\omega ^{}`$, $`I_{j_0}=I_\omega `$. We also split $`H_0`$ into two parts: $`H_0^{}`$ commuting with $`a_\omega `$ and $`a_\omega ^{}`$, and $`H_0^{}=\omega N_\omega `$, where
$`[a_\omega ,N_\omega ]`$ $`=`$ $`a_\omega ,`$ (5)
$`[a_\omega ^{},N_\omega ]`$ $`=`$ $`a_\omega ^{}.`$ (6)
Alternatively, we may assume that $`N_\omega `$ exists and define $`H_0^{}=H_0\omega N_\omega `$.
The model is given by the full Hamiltonian
$`H=\omega _0R_3+H_0+gR_+a_\omega +\overline{g}R_{}a_\omega ^{}.`$ (7)
We employ the usual notation Allen where $`R_l=\sigma _l/2`$, $`R_\pm =R_1\pm iR_2`$, $`\sigma _l`$ are the Pauli matrices, and $`g`$ is a complex coupling parameter. It is useful to split $`H`$ into three mutually commuting parts:
$`H`$ $`=`$ $`H_0^{}+\omega 𝒩+\mathrm{\Omega },`$ (8)
$`𝒩`$ $`=`$ $`R_3+N_\omega ,`$ (9)
$`\mathrm{\Omega }`$ $`=`$ $`\mathrm{\Delta }R_3+gR_+a_\omega +\overline{g}R_{}a_\omega ^{}.`$ (10)
$`\mathrm{\Delta }=\omega _0\omega `$ is the detuning. The evolution operator factorizes:
$`U_t`$ $`=`$ $`e^{iH_0^{}t}e^{i\omega 𝒩t}V_t,`$ (11)
$`V_t`$ $`=`$ $`e^{i\mathrm{\Omega }t}=\mathrm{cos}(\mathrm{\Omega }_Rt)i{\displaystyle \frac{\mathrm{sin}(\mathrm{\Omega }_Rt)}{\mathrm{\Omega }_R^2}}\mathrm{\Omega },`$ (12)
$`\mathrm{\Omega }_R`$ $`=`$ $`\sqrt{\mathrm{\Delta }^2/4+|g|^2X},`$ (13)
$`X`$ $`=`$ $`R_3I_\omega +a_\omega ^{}a_\omega +{\displaystyle \frac{1}{2}}I_\omega .`$ (14)
Let us note that the operators $`=R_3I_\omega +a_\omega ^{}a_\omega `$, occurring in $`X`$, and $`𝒩=R_3+N_\omega `$, occurring in $`H`$, should not be identified, although both are equivalent in irreducible representations of CCR, and both commute with $`H`$ independently of the choice of representation. Constructing $``$ and $`𝒩`$ in the reducible representation we shall explicitly show that they are essentially different.
Only (12) is relevant in the context of Heisenberg-picture dynamics of $`R_3(t)`$ since $`[R_3,H_0^{}]=[R_3,𝒩]=0`$. One can verify by a straightforward calculation that $`idU_t/dt=U_tH`$, and it is clear that the solution is valid for any representation of the Lie algebra (2). Employing (12) one finds the Heisenberg picture evolution of $`R_3`$:
$`R_3(t)=R_3\left(12|g|^2X{\displaystyle \frac{\mathrm{sin}^2(\mathrm{\Omega }_Rt)}{\mathrm{\Omega }_R^2}}\right)+\left({\displaystyle \frac{\mathrm{\Delta }}{2}}{\displaystyle \frac{\mathrm{sin}^2(\mathrm{\Omega }_Rt)}{\mathrm{\Omega }_R^2}}i{\displaystyle \frac{\mathrm{sin}(2\mathrm{\Omega }_Rt)}{2\mathrm{\Omega }_R}}\right)gR_+a_\omega +\left({\displaystyle \frac{\mathrm{\Delta }}{2}}{\displaystyle \frac{\mathrm{sin}^2(\mathrm{\Omega }_Rt)}{\mathrm{\Omega }_R^2}}+i{\displaystyle \frac{\mathrm{sin}(2\mathrm{\Omega }_Rt)}{2\mathrm{\Omega }_R}}\right)\overline{g}R_{}a_\omega ^{}.`$
(15)
The proof of (15) is outlined in the Appendix.
The next important notion that can be introduced at a representation independent level is the displacement operator. The operator is defined in the usual way as
$`D(z)`$ $`=`$ $`\mathrm{exp}{\displaystyle \underset{j}{}}\left(z_ja_j^{}\overline{z}_ja_j\right)`$ (16)
and satisfies
$`D(z)^{}a_jD(z)`$ $`=`$ $`a_j+z_jI_j,`$ (17)
$`D(z)^{}a_j^{}D(z)`$ $`=`$ $`a_j^{}+\overline{z}_jI_j,`$ (18)
$`D(z)^{}I_jD(z)`$ $`=`$ $`I_j.`$ (19)
Acting with $`D(z)`$ on a vacuum vector we obtain a coherent state. Its form depends on what is meant by vacuum in a given representation.
## III Evolution of atomic inversion operator in irreducible representations
Assume we work in an irreducible representation with some carrier Hilbert space containing a vector $`|0`$ annihilated by all $`a_\omega `$. The abstract $``$-conjugation in the algebra can be replaced by Hermitian conjugation of operators in the representation. The representation satisfies
$`[a_j,a_k^{}]=𝒵\delta _{jk}1,`$ (20)
where $`𝒵`$ is a real positive number (by Schur’s lemma an element that commutes with all elements of the algebra is a constant times identity if the representation is irreducible, i.e. $`I_j=𝒵1`$ for some $`𝒵`$, and for all $`j`$; had $`𝒵`$ been negative we would have called $`a_j`$ a creation operator, taken vacuum annihilated by $`a_j^{}`$, and appropriately adjusted the notation). The solution (15) involves the operator $`|g|^2X`$ whose representation reads
$`|g|^2X`$ $`=`$ $`|g|^2\left(𝒵R_3+a_\omega ^{}a_\omega +𝒵/2\right)`$ (21)
$`=`$ $`|\stackrel{~}{g}|^2\left(R_3+\stackrel{~}{a}_\omega ^{}\stackrel{~}{a}_\omega +1/2\right)=|\stackrel{~}{g}|^2\stackrel{~}{X},`$ (22)
where $`\stackrel{~}{a}_\omega =a_\omega /\sqrt{𝒵}`$, $`[\stackrel{~}{a}_j,\stackrel{~}{a}_k^{}]=\delta _{jk}1`$, and $`\stackrel{~}{g}=\sqrt{𝒵}g`$. Now,
$`R_3(t)=R_3\left(12|\stackrel{~}{g}|^2\stackrel{~}{X}{\displaystyle \frac{\mathrm{sin}^2(\stackrel{~}{\mathrm{\Omega }}_Rt)}{\stackrel{~}{\mathrm{\Omega }}_R^2}}\right)+\left({\displaystyle \frac{\mathrm{\Delta }}{2}}{\displaystyle \frac{\mathrm{sin}^2(\stackrel{~}{\mathrm{\Omega }}_Rt)}{\stackrel{~}{\mathrm{\Omega }}_R^2}}i{\displaystyle \frac{\mathrm{sin}(2\stackrel{~}{\mathrm{\Omega }}_Rt)}{2\stackrel{~}{\mathrm{\Omega }}_R}}\right)\stackrel{~}{g}R_+\stackrel{~}{a}_\omega +\left({\displaystyle \frac{\mathrm{\Delta }}{2}}{\displaystyle \frac{\mathrm{sin}^2(\stackrel{~}{\mathrm{\Omega }}_Rt)}{\stackrel{~}{\mathrm{\Omega }}_R^2}}+i{\displaystyle \frac{\mathrm{sin}(2\stackrel{~}{\mathrm{\Omega }}_Rt)}{2\stackrel{~}{\mathrm{\Omega }}_R}}\right)\overline{\stackrel{~}{g}}R_{}\stackrel{~}{a}_\omega ^{},`$
(23)
where $`\stackrel{~}{\mathrm{\Omega }}=\sqrt{\mathrm{\Delta }^2/4+|\stackrel{~}{g}|^2\stackrel{~}{X}}`$. Clearly, the only difference between different representations is in the values of the coupling constant $`\stackrel{~}{g}=g\sqrt{𝒵}`$. Since $`g`$ is proportional to the electron charge, different representations effectively differ by the value of the electron charge $`\stackrel{~}{e}_0=e_0\sqrt{𝒵}`$. This type of rescaling is exactly what occurs in transition from the bare charge $`e_0`$ to the physical, renormalized charge $`e_{\mathrm{ph}}=e_0\sqrt{Z_3}`$ typical of renormalized electromagnetic fields. Let us note finally that
$`\stackrel{~}{a}_jD(z)|0=z_j\sqrt{𝒵}D(z)|0`$ (24)
and thus the relation between the parameter $`z`$ and the ‘physical’ amplitude $`\stackrel{~}{z}`$ is also renormalized: $`\stackrel{~}{z}=z\sqrt{𝒵}`$. After these rescalings of the ‘bare’ parameters we will obtain identical formulas for coherent state averages of $`R_3(t)`$, independently of our choice of the irreducible representation. We conclude that the theorem of von Neumann does not bring to our problem anything physically important.
So let us switch to reducible representations.
## IV $`N<\mathrm{}`$ reducible representation
The representation is constructed as follows. For simplicity we ignore the polarization degree of freedom. Take an operator $`a`$ satisfying $`[a,a^{}]=1`$ and the kets $`|𝒌`$ corresponding to standing waves in some cavity. We define
$`a(𝒌)=|𝒌𝒌|a,I(𝒌)=|𝒌𝒌|1.`$ (25)
The operators (25) satisfy (1), where $`\delta (𝒌,𝒌^{})`$ is the 3D Kronecker delta. The fact that $`I(𝒌)`$ is not proportional to the identity means that the representation is reducible. In our terminology this is the ‘$`N=1`$ representation’. Its Hilbert space $``$ is spanned by the kets $`|𝒌,n=|𝒌|n`$, where $`a^{}a|n=n|n`$. Such a Hilbert space represents essentially a single harmonic oscillator of indefinite frequency 1 . An important property of the representation is that $`_kI(𝒌)=I`$ is the identity operator in $``$. A vacuum of this representation is given by any state annihilated by all $`a(𝒌)`$. The vacuum state is not unique and belongs to the subspace spanned by $`|𝒌,0`$. In our notation a $`N=1`$ vacuum state reads $`|O=_kO(𝒌)|𝒌,0`$ and is normalized by $`_k|O(𝒌)|^2=_kZ(𝒌)=1`$. Such a vacuum is exactly analogous to a single-oscillator ground-state center-of-mass wavepacket. As shown in 2 ; 3 (for a more complete discussion see the companion paper 4 ) in a fully relativistic formulation the maximal probability $`Z=\mathrm{max}_k\{Z(𝒌)\}`$ is a Poincaré invariant, an aspect worth keeping in mind. For $`N1`$ the representation space is given by the tensor power $`\underset{¯}{}=^N`$, i.e. we take the Hilbert space of $`N`$ (bosonic) harmonic oscillators. Let $`A:`$ be any operator for $`N=1`$. We denote $`A^{(n)}=I^{(n1)}AI^{(Nn)}`$, $`A^{(n)}:\underset{¯}{}\underset{¯}{}`$, for $`1nN`$. For arbitrary $`N`$ the representation is defined by
$`\underset{¯}{a}(𝒌)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N}}}{\displaystyle \underset{n=1}{\overset{N}{}}}a(𝒌)^{(n)},`$ (26)
$`\underset{¯}{I}(𝒌)`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{n=1}{\overset{N}{}}}I(𝒌)^{(n)},`$ (27)
$`[\underset{¯}{a}(𝒌),\underset{¯}{a}(𝒌^{})^{}]`$ $`=`$ $`\delta (𝒌,𝒌^{})\underset{¯}{I}(𝒌),`$ (28)
$`{\displaystyle \underset{k}{}}\underset{¯}{I}(𝒌)`$ $`=`$ $`\underset{¯}{I}=I^N`$ (29)
and the $`N`$-oscillator vacuum is the $`N`$-fold tensor power of the $`N=1`$ case, a kind of Bose-Einstein condensate consisting of $`N`$ wavepackets:
$`|\underset{¯}{O}=|O\mathrm{}|O=|O^N.`$ (30)
The free-field Hamiltonian is, for $`N=1`$ and $`\omega _k=|𝒌|`$,
$`H_0=_k\omega _ka(𝒌)^{}a(𝒌)=_k\omega _k|𝒌𝒌|a^{}a.`$ (31)
In each eigensubspace with fixed $`|𝒌`$ the operator $`H_0`$ is just an ordinary Hamiltonian of the oscillator with frequency $`\omega _k`$. For arbitrary $`N`$ the generator of free field evolution is the Hamiltonian of $`N`$ noninteracting oscillators, i.e. $`\underset{¯}{H}_0=_{n=1}^NH_0^{(n)}`$. Let us stress that $`\underset{¯}{H}_0`$ should not be confused with $`_k\omega _k\underset{¯}{a}(𝒌)^{}\underset{¯}{a}(𝒌)`$ which does not have a clear interpretation in this context and does not describe noninteracting oscillators. We have seen already that the operator $`\underset{¯}{a}(𝒌)^{}\underset{¯}{a}(𝒌)`$ will nevertheless play an important role in the Jaynes-Cummings problem. Our definition of $`\underset{¯}{H}_0`$ implies that $`[\underset{¯}{a}(𝒌),\underset{¯}{H}_0]=\omega _k\underset{¯}{a}(𝒌)`$ which is the formula we required at the representation independent level. Let us now select some $`𝒑`$ with frequency $`\omega =|𝒑|`$ to be the mode that interacts with the two-level atom. To simplify further calculations we denote $`\underset{¯}{a}(𝒑)=\underset{¯}{a}_\omega `$, $`\underset{¯}{I}(𝒑)=\underset{¯}{I}_\omega `$, $`|𝒑𝒑|1=P_1`$, $`P_0=IP_1`$.
The projector $`P_1`$ is related to $`\underset{¯}{I}_\omega `$ by
$`\underset{¯}{I}_\omega `$ $`=`$ $`{\displaystyle \frac{1}{N}}\left(P_1I\mathrm{}I+\mathrm{}+I\mathrm{}IP_1\right).`$ (32)
One recognizes in (32) the frequency-of-success operator discussed in the context of quantum laws of large numbers Hartle ; Farhi ; Gutman ; Casinello ; Aharonov ; Finkelstein ; Caves . Eigenvalues of $`\underset{¯}{I}_\omega `$ coincide with all the possible frequences of ‘heads’ in $`N`$ trials of coin tossing, i.e. $`\underset{¯}{I}_\omega =_{s=0}^N(s/N)P(s/N)`$ with spectral projectors
$`P(s/N)`$ $`=`$ $`{\displaystyle \underset{s_1+\mathrm{}+s_N=s;s_j=0,1}{}}P_{s_1}\mathrm{}P_{s_N}.`$ (33)
The explicit form (33) shows that $`P(s/N)`$ commute with all $`\underset{¯}{a}(𝒌)`$, $`\underset{¯}{a}(𝒌)^{}`$, and $`\underset{¯}{I}(𝒌)`$. Of particular importance is the splitting of the Jaynes-Cummings problem into subspaces corresponding to a given $`s/N`$. We define
$`a(s)`$ $`=`$ $`\underset{¯}{a}_\omega P(s/N),`$ (34)
$`a(s)^{}`$ $`=`$ $`\underset{¯}{a}_\omega ^{}P(s/N),`$ (35)
and obtain the representation
$`[a(s),a(s)^{}]`$ $`=`$ $`(s/N)P(s/N)`$ (36)
of CCR in the Hilbert space $`(s)=P(s/N)\underset{¯}{}`$.
Each $`(s)`$ is an invariant subspace for the Jaynes-Cummings dynamics, and in each such a subspace we effectively deal with a representation given by (20) whose $`𝒵=s/N`$, as we shall see in the next Section.
## V Evolution of atomic inversion operator in $`N<\mathrm{}`$ representation
As the general formula (15) is valid also in this representation and all the operators that occur there commute with $`P(s/N)`$, we begin with splitting $`R_3(t)`$ into parts acting in the invariant subspaces $`(s)`$. We employ the usual notation $`R_3=(|++|||)/2`$ and treat $`|++|+||`$ as the identity operator (to be more exact we should tensor the atomic-space identity with the field-space identity, but we prefer this simplified convention).
Denoting $`X(s)=XP(s/N)`$ we get $`X=_{s=0}^NX(s)`$ and
$`X(s)`$ $`=`$ $`\left({\displaystyle \frac{s}{N}}R_3+a(s)^{}a(s)+{\displaystyle \frac{s}{2N}}\right)P(s/N)`$ (37)
$`=`$ $`{\displaystyle \frac{s}{N}}|++|P(s/N)+a(s)^{}a(s).`$ (38)
Commutation relation (36) implies that eigenvalues of $`a(s)^{}a(s)`$ are $`sn/N`$, $`n=0,1,2,\mathrm{}`$. In spectral representation
$`a(s)^{}a(s)={\displaystyle \frac{s}{N}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}n\mathrm{\Pi }(n,s).`$ (39)
The spectral projectors satisfy $`\mathrm{\Pi }(n,s)P(s/N)=\mathrm{\Pi }(n,s)`$, $`_{n=0}^{\mathrm{}}\mathrm{\Pi }(n,s)=P(s/N)`$. Spectral representation of $`X(s)`$ therefore reads
$`X(s)`$ $`=`$ $`{\displaystyle \frac{s}{N}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}(|++|+n)\mathrm{\Pi }(n,s)`$ (40)
$`=`$ $`{\displaystyle \frac{s}{N}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}n\widehat{\mathrm{\Pi }}(n,s).`$ (41)
Spectral projectors $`\widehat{\mathrm{\Pi }}(n,s)`$ of $`X(s)`$ are related to $`\mathrm{\Pi }(n,s)`$ of $`a(s)^{}a(s)`$ by
$`\widehat{\mathrm{\Pi }}(n,s)`$ $`=`$ $`|++|\mathrm{\Pi }(n1,s)+||\mathrm{\Pi }(n,s),`$ (42)
for $`n>0`$, and $`\widehat{\mathrm{\Pi }}(0,s)=||\mathrm{\Pi }(0,s)`$. In the Appendix we show that $`[R_+a(s),\widehat{\mathrm{\Pi }}(n,s)]=[R_{}a(s)^{},\widehat{\mathrm{\Pi }}(n,s)]`$. Evolution of the atomic inversion operator is given in our reducible representation by
$`R_3(t)`$ $`=`$ $`R_3{\displaystyle \underset{s=0}{\overset{N}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}2R_3|g|^2{\displaystyle \frac{ns}{N}}{\displaystyle \frac{\mathrm{sin}^2\left(t\sqrt{\mathrm{\Delta }^2/4+|g|^2ns/N}\right)}{\mathrm{\Delta }^2/4+|g|^2ns/N}}\widehat{\mathrm{\Pi }}(n,s)`$ (43)
$`+{\displaystyle \underset{s=0}{\overset{N}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{\mathrm{\Delta }}{2}}{\displaystyle \frac{\mathrm{sin}^2\left(t\sqrt{\mathrm{\Delta }^2/4+|g|^2ns/N}\right)}{\mathrm{\Delta }^2/4+|g|^2ns/N}}i{\displaystyle \frac{\mathrm{sin}\left(2t\sqrt{\mathrm{\Delta }^2/4+|g|^2ns/N}\right)}{2\sqrt{\mathrm{\Delta }^2/4+|g|^2ns/N}}}\right)\widehat{\mathrm{\Pi }}(n,s)gR_+a(s)`$
$`+{\displaystyle \underset{s=0}{\overset{N}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{\mathrm{\Delta }}{2}}{\displaystyle \frac{\mathrm{sin}^2\left(t\sqrt{\mathrm{\Delta }^2/4+|g|^2ns/N}\right)}{\mathrm{\Delta }^2/4+|g|^2ns/N}}+i{\displaystyle \frac{\mathrm{sin}\left(2t\sqrt{\mathrm{\Delta }^2/4+|g|^2ns/N}\right)}{2\sqrt{\mathrm{\Delta }^2/4+|g|^2ns/N}}}\right)\overline{g}R_{}a(s)^{}\widehat{\mathrm{\Pi }}(n,s).`$
Formula (43) can be directly compared to (23) by means of the following simple rule: Skip the sum over $`s`$ and set $`s/N=1`$. An alternative recipe is to skip the sum over $`s`$ and set $`s/N=𝒵`$, which corresponds to CCR with $`𝒵1`$ at the right-hand-side. Yet another intuitive rule can be found by means of the law of large numbers and works for the weak limit $`N\mathrm{}`$. To make it precise one has to switch to the level of averages. We will see that the weak law of large numbers plays a role of a correspondence principle between our $`N<\mathrm{}`$ formalism and the standard regularized one.
## VI Evolution of atomic inversion in $`N<\mathrm{}`$ representation
Acting on the vacuum state (30) with the displacement operator (16) one obtains a coherent state. Here we are interested in a monochromatic coherent state with frequency $`\omega `$
$`|\underset{¯}{z}`$ $`=`$ $`\mathrm{exp}\left(z\underset{¯}{a}_\omega ^{}\overline{z}\underset{¯}{a}_\omega \right)|\underset{¯}{O}.`$ (44)
The coherent state is not an eigenstate of the annihilation operator $`\underset{¯}{a}_\omega `$ but a direct sum of its eigenstates. Indeed, the decomposition
$`|\underset{¯}{z}`$ $`=`$ $`{\displaystyle \underset{s=0}{\overset{N}{}}}P(s/N)|\underset{¯}{z}={\displaystyle \underset{s=0}{\overset{N}{}}}|z(s)`$ (45)
accompanying $`\underset{¯}{a}_\omega =_{s=0}^N\underset{¯}{a}_\omega P(s/N)=_{s=0}^Na(s)`$ implies that
$`a(s)|z(s)=(s/N)z|z(s).`$ (46)
The analogy of the latter eigenvalue problem to (24) is evident. Alternatively, one can say that the coherent state is a generalized eigenvector of $`\underset{¯}{a}_\omega `$, i.e.
$`\underset{¯}{a}_\omega |\underset{¯}{z}`$ $`=`$ $`z\underset{¯}{I}_\omega |\underset{¯}{z}.`$ (47)
Another state of interest, particularly in the context of experiments, is the mixture
$`\rho `$ $`=`$ $`(p_+|++|+p_{}||)\rho _{\mathrm{field}}.`$ (48)
### VI.1 Reduced inversion operator: Coherent state
Evaluating an average of (43) in an arbitrary coherent state $`|\underset{¯}{z}`$ we obtain the reduced operator
$`R_z(t)=\underset{¯}{z}|R_3(t)|\underset{¯}{z}`$ (49)
involving only the atomic degrees of freedom:
$`R_z(t)`$ $`=`$ $`R_3+{\displaystyle \underset{s=0}{\overset{N}{}}}{\displaystyle \frac{s}{N}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\underset{¯}{z}|\mathrm{\Pi }(n,s)|\underset{¯}{z}`$ (50)
$`\times (|g|^2(n+1){\displaystyle \frac{\mathrm{sin}^2\left(t\sqrt{\mathrm{\Delta }^2/4+|g|^2(n+1)s/N}\right)}{\mathrm{\Delta }^2/4+|g|^2(n+1)s/N}}|++|+|g|^2n{\displaystyle \frac{\mathrm{sin}^2\left(t\sqrt{\mathrm{\Delta }^2/4+|g|^2ns/N}\right)}{\mathrm{\Delta }^2/4+|g|^2ns/N}}||`$
$`+zg\left({\displaystyle \frac{\mathrm{\Delta }}{2}}{\displaystyle \frac{\mathrm{sin}^2\left(t\sqrt{\mathrm{\Delta }^2/4+|g|^2(n+1)s/N}\right)}{\mathrm{\Delta }^2/4+|g|^2(n+1)s/N}}i{\displaystyle \frac{\mathrm{sin}\left(2t\sqrt{\mathrm{\Delta }^2/4+|g|^2(n+1)s/N}\right)}{2\sqrt{\mathrm{\Delta }^2/4+|g|^2(n+1)s/N}}}\right)R_+`$
$`+\overline{z}\overline{g}({\displaystyle \frac{\mathrm{\Delta }}{2}}{\displaystyle \frac{\mathrm{sin}^2\left(t\sqrt{\mathrm{\Delta }^2/4+|g|^2(n+1)s/N}\right)}{\mathrm{\Delta }^2/4+|g|^2(n+1)s/N}}+i{\displaystyle \frac{\mathrm{sin}\left(2t\sqrt{\mathrm{\Delta }^2/4+|g|^2(n+1)s/N}\right)}{2\sqrt{\mathrm{\Delta }^2/4+|g|^2(n+1)s/N}}})R_{}),`$
where (see Appendix)
$`\underset{¯}{z}|\mathrm{\Pi }(n,s)|\underset{¯}{z}`$ (51)
$`={\displaystyle \frac{|z\sqrt{s/N}|^{2n}}{n!}}e^{|z\sqrt{s/N}|^2}\left(\begin{array}{c}N\\ s\end{array}\right)Z_\omega ^s(1Z_\omega )^{Ns}.`$ (54)
Here
$`Z_\omega =\underset{¯}{O}|\underset{¯}{I}_\omega |\underset{¯}{O}=O|I_\omega |O=|O(𝒑)|^2`$ (55)
is the probability of finding the momentum $`𝒑`$ corresponding to the resonant mode $`\omega =|𝒑|`$, if the vacuum state of the field is $`|\underset{¯}{O}`$. In the last term of (51) one recognizes the binomial distribution for $`N`$ Bernoulli trials with probability of success equal to $`Z_\omega `$.
### VI.2 Vacuum-state initial condition: $`|\underset{¯}{z}=|\underset{¯}{O}`$
Assume initially there are no photons and the atom is in either ground or excited state, i.e.
$`|\mathrm{\Psi }=|\pm |\underset{¯}{O}.`$ (56)
In this case only the $`n=0`$ term counts in (50). The case of atomic ground state, $`|\mathrm{\Psi }=||\underset{¯}{O}`$, is trivial since
$`w(t)=\mathrm{\Psi }|R_3(t)|\mathrm{\Psi }=|R_z(t)|=1/2.`$ (57)
However, starting with the excited state $`|\mathrm{\Psi }=|+|\underset{¯}{O}`$ we find
$`w(t)=\mathrm{\Psi }|R_3(t)|\mathrm{\Psi }=+|R_z(t)|+={\displaystyle \frac{1}{2}}{\displaystyle \underset{s=0}{\overset{N}{}}}|g|^2{\displaystyle \frac{s}{N}}{\displaystyle \frac{\mathrm{sin}^2\sqrt{\mathrm{\Delta }^2/4+|g|^2s/N}t}{\mathrm{\Delta }^2/4+|g|^2s/N}}\left(\begin{array}{c}N\\ s\end{array}\right)Z_\omega ^s(1Z_\omega )^{Ns}.`$ (60)
The law of large numbers for the binomial distribution implies that
$`\underset{N\mathrm{}}{lim}w(t)={\displaystyle \frac{1}{2}}|g|^2Z_\omega {\displaystyle \frac{\mathrm{sin}^2\sqrt{\mathrm{\Delta }^2/4+|g|^2Z_\omega }t}{\mathrm{\Delta }^2/4+|g|^2Z_\omega }},`$ (61)
i.e. the frequency $`s/N`$ approaches the probability of success in a single trial, that is $`s/NZ_\omega `$.
It is instructive to compare this result with the one we would have found had we started with the general irreducible representation whose right-hand-side is $`𝒵1`$. The corresponding result reads
$`w(t)={\displaystyle \frac{1}{2}}|g|^2𝒵{\displaystyle \frac{\mathrm{sin}^2\sqrt{\mathrm{\Delta }^2/4+|g|^2𝒵}t}{\mathrm{\Delta }^2/4+|g|^2𝒵}}.`$ (62)
These two formulas appear so similar that one can easily overlook an important difference: The parameter $`𝒵`$ in (62) is a number independent of $`\omega `$, whereas $`Z_\omega `$ in (61) is a value of a function vanishing for $`\omega \mathrm{}`$. Charge renormalization cannot be performed in the same way in both models, because charge must remain a relativistic invariant.
In the irreducible case the recipe is simple: $`e_{\mathrm{ph}}=e_0\sqrt{𝒵}`$. In the reducible case one first extracts from $`Z_\omega `$ the relativistic invariant $`Z=\mathrm{max}_k\{Z(𝒌)\}`$ 2 ; 3 ; 4 . Then one writes $`Z_\omega =Z\chi _\omega `$, redefines charge $`e_{\mathrm{ph}}=e_0\sqrt{Z}`$, and finally
$`\underset{N\mathrm{}}{lim}w(t)={\displaystyle \frac{1}{2}}|g_{\mathrm{ph}}|^2\chi _\omega {\displaystyle \frac{\mathrm{sin}^2\sqrt{\mathrm{\Delta }^2/4+|g_{\mathrm{ph}}|^2\chi _\omega }t}{\mathrm{\Delta }^2/4+|g_{\mathrm{ph}}|^2\chi _\omega }}.`$ (63)
The solutions (62) and (63) are identical up to the presence of the cut-off function $`\chi _\omega `$ in (63). Needless to say the function would necessarily regularize the interaction if we decided to work with extremely high frequencies $`\omega `$.
Let us note that we did not introduce any cut-off in the Hamiltonian. The cut-off has appeared automatically through the structure of the vacuum state typical of the reducible representation. The theory gets regularized even though we do not really need it in such a simple example. The regularization neither solves here any problem nor spoils anything. If we assume that for optical frequencies $`\chi _\omega 1`$ we obtain exact agreement between the irreducible case and the $`N\mathrm{}`$ limit of the reducible one. This is an example of the correspondence principle we have mentioned in the introduction.
The next step is to understand if one really needs $`N=\mathrm{}`$, and if some finite $`N`$ cannot, in fact, be consistent with experimental data. We first rewrite (60) by means of the renormalized coupling
$`w(t)={\displaystyle \frac{1}{2}}{\displaystyle \underset{s=0}{\overset{N}{}}}{\displaystyle \frac{|g_{\mathrm{ph}}|^2}{Z}}{\displaystyle \frac{s}{N}}{\displaystyle \frac{\mathrm{sin}^2\sqrt{\mathrm{\Delta }^2/4+|g_{\mathrm{ph}}|^2s/(ZN)}t}{\mathrm{\Delta }^2/4+|g_{\mathrm{ph}}|^2s/(ZN)}}\left(\begin{array}{c}N\\ s\end{array}\right)Z_\omega ^s(1Z_\omega )^{Ns}.`$ (66)
For $`N`$ large enough the binomial distribution can be approximated by a Gaussian
$`\left(\begin{array}{c}N\\ s\end{array}\right)Z_\omega ^s(1Z_\omega )^{Ns}{\displaystyle \frac{e^{\frac{(sNZ)^2}{2NZ(1Z)}}}{\sqrt{2\pi NZ(1Z)}}}{\displaystyle \frac{e^{\frac{(sNZ)^2}{2NZ}}}{\sqrt{2\pi NZ}}}`$ (69)
whose shape is controlled mainly by the product $`NZ`$. The smaller $`Z`$ the less important its exact value (the second approximate equality holds for small $`Z`$). So $`NZ`$ is the parameter we may have a chance of seeing in experiments. Let us stress that both $`N`$ and $`Z`$ are relativistically invariant.
The atomic inversion $`w(t)`$ is here a sum of $`N`$ oscillations, each at a different Rabi frequency. It is clear that for finite $`N`$ the evolution of $`w(t)`$ will reveal collapses and revivals, and not a simple Rabi oscillation as would be expected on the basis of irreducible representations (cf. the analysis of the experiment in the next subsection, and in particular Fig. 6). The solution is similar to those known from the standard analysis of collapses and revivals of coherent-state evolutions Narozhny .
### VI.3 Experiment: Thermal mixture as the initial condition
Replacing the Poisson statistics in (51) by thermal probability we obtain a formula applicable to situations where the coherent light is replaced by a thermal mixture. We thus assume (48) with
$`\mathrm{Tr}\rho _{\mathrm{field}}\mathrm{\Pi }(n,s)`$ $`=`$ $`𝙿(n)\left(\begin{array}{c}N\\ s\end{array}\right)Z_\omega ^s(1Z_\omega )^{Ns},`$ (72)
$`𝙿(n)`$ $`=`$ $`{\displaystyle \frac{\overline{n}^n}{(1+\overline{n})^{(n+1)}}},`$ (73)
and find the formula directly applicable to experiments:
$`w(t)`$ $`=`$ $`{\displaystyle \frac{p_+p_{}}{2}}+{\displaystyle \underset{s=0}{\overset{N}{}}}\left(\begin{array}{c}N\\ s\end{array}\right)Z_\omega ^s(1Z_\omega )^{Ns}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}𝙿(n)\left({\displaystyle \frac{p_{}\overline{n}}{1+\overline{n}}}p_+\right)|g|^2{\displaystyle \frac{(n+1)s}{N}}{\displaystyle \frac{\mathrm{sin}^2\left(t\sqrt{\frac{\mathrm{\Delta }^2}{4}+|g|^2\frac{(n+1)s}{N}}\right)}{\mathrm{\Delta }^2/4+|g|^2(n+1)s/N}}.`$ (76)
Taking the limit $`N\mathrm{}`$ we obtain
$`w(t)`$ $`=`$ $`{\displaystyle \frac{p_+p_{}}{2}}+{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}𝙿(n)\left({\displaystyle \frac{p_{}\overline{n}}{1+\overline{n}}}p_+\right)|g_{\mathrm{ph}}|^2\chi _\omega {\displaystyle \frac{\mathrm{sin}^2\left(t\sqrt{\mathrm{\Delta }^2/4+|g_{\mathrm{ph}}|^2(n+1)\chi _\omega }\right)}{\mathrm{\Delta }^2/4+|g_{\mathrm{ph}}|^2(n+1)\chi _\omega }},`$ (77)
which, up to $`\chi _\omega `$, is known from irreducible representations.
In what follows we compare theoretical results with the precise data on optical Rabi oscillations reported by the Kastler-Brossel Laboratory group from Paris Haroche .
The relevant plot is Fig. 2A in Haroche . In Fig. 1 we show the data vs. standard theoretical predictions based on irreducible representations. We assume exact resonance condition $`\mathrm{\Delta }=0`$, $`g_{\mathrm{ph}}=47`$ kHz, and $`p_+=0.99`$. As the initial field state we take the thermal mixture $`𝙿(n)`$ for $`\overline{n}=0.05`$ (Fig. 2 shows that $`𝙿(n)`$ agrees with Fig. 2$`\alpha `$ from Haroche ). The blue curve represents the Rabi oscillation additionally damped by the factor $`\mathrm{exp}(t/T_{\mathrm{cav}})`$, $`T_{\mathrm{cav}}=220`$ $`\mu `$s, which could be expected on the basis of cavity relaxation parameters. The red curve is the same Rabi oscillation but with stronger damping, $`T_{\mathrm{cav}}=45`$ $`\mu `$s. As we can see the red curve almost agrees with the data; the first two minima are somewhat lower than the data but the analysis does not include center-of-mass motion, finite duration of experiment, and the resulting frequency spread and detunings.
So the red curve yields a reasonable agreement with experiment. The problem is that the cavity lifetime was 220 $`\mu `$s and not 45 $`\mu `$s, and the data should coincide with the blue curve. This observation agrees with the remark of Brune et al. that ‘cavity relaxation plays a marginal role in the decrease of oscillation’. The authors further write: ‘Dark counts (…) are one of the main causes of oscillation damping(…). Decoherence by collisions with background gas may also contribute to the oscillation relaxation’ (Haroche , p. 1801). The relaxation observed experimentally thus appears to be stronger than expected, but the authors do not discuss the subtlety in much detail. In order to fit the data they used ‘damped sinusoids’, but did not explain what kind of a damping factor (or factors?) were employed. Also the Fourier analysis was performed on time-symmetrized signals, and not just on damped ones. This is important since symmetrized signals are closer to revivals-collapses typical of beats than to damped oscillations. In the light of the results we discuss in the present paper it is clear that a more detailed analysis of experimental damping factors is required.
In Fig. 3 we show the result with the same mixed initial condition but now computed by means of the reducible representation. The parameters are $`\mathrm{\Delta }=0`$, $`g_{\mathrm{ph}}=47`$ kHz, $`Z=0.1`$, and $`N=280`$. The blue curve involves no damping (ideal cavity); the red curve is additionally damped by $`\mathrm{exp}(t/T_{\mathrm{cav}})`$, $`T_{\mathrm{cav}}=220`$ $`\mu `$s. The first two minima are again lower than the data, but the general agreement with experiment is acceptable. As one can see the damping is caused here mainly by the beats. Unfortunately we did not possess all the data needed for a realistic comparison with the experiment.
Now, can we distinguish between true damping and beats? In principle yes, but we must wait longer and have better cavities. Fig. 4 shows the dynamics of $`w(t)`$ in the reducible representation with exact vacuum initial condition, ideal cavity, $`Z=0.1`$, $`N=280`$, $`\mathrm{\Delta }=0`$, and $`g_{\mathrm{ph}}=47`$ kHz for $`0<t<0.0015`$ s. Fig. 5 shows the same dynamics but now additionally damped by $`\mathrm{exp}(t/T_{\mathrm{cav}})`$, $`T_{\mathrm{cav}}=500`$ $`\mu `$s. In principle, in a roughly twice better cavity we might see the revival. It must be stressed that this revival has no counterpart if field is quantized in irreducible representation, and vanishes if $`N\mathrm{}`$ (the limit taken with $`Z`$ fixed). Fig. 6 shows the dynamics of $`w(t)`$ in ideal cavity, exact vacuum, and for $`N=10000`$, the remaining parameters being kept as before. This case is interesting since up to 100 $`\mu `$s, that is the time available in the discussed experiment, the plot is practically indistinguishable from the result based on irreducible representations, whereas for longer times we observe collapse and revival. The greater $`NZ`$, the later comes the revival.
The examples show that the origin of relaxation should be carefully reexamined.
### VI.4 General coherent-state initial condition
Let us now replace vacuum by a general $`|\underset{¯}{z}`$ but for simplicity keep the initial atomic state to be $`|+`$. Then
$`w(t)`$ $`=`$ $`+|R_z(t)|+`$ (78)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{s=0}{\overset{N}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{|g_{\mathrm{ph}}|^2(n+1)}{Z}}{\displaystyle \frac{s}{N}}{\displaystyle \frac{\mathrm{sin}^2\left(t\sqrt{\mathrm{\Delta }^2/4+|g_{\mathrm{ph}}|^2(n+1)s/(ZN)}\right)}{\mathrm{\Delta }^2/4+|g_{\mathrm{ph}}|^2(n+1)s/(ZN)}}{\displaystyle \frac{|z\sqrt{\frac{s}{N}}|^{2n}}{n!}}e^{|z\sqrt{\frac{s}{N}}|^2}\left(\begin{array}{c}N\\ s\end{array}\right)Z_\omega ^s(1Z_\omega )^{Ns}.`$ (81)
The limiting form, for $`N\mathrm{}`$, is again familiar
$`\underset{N\mathrm{}}{lim}w(t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}|g_{\mathrm{ph}}|^2(n+1)\chi _\omega {\displaystyle \frac{\mathrm{sin}^2\left(t\sqrt{\mathrm{\Delta }^2/4+|g_{\mathrm{ph}}|^2(n+1)\chi _\omega }\right)}{\mathrm{\Delta }^2/4+|g_{\mathrm{ph}}|^2(n+1)\chi _\omega }}{\displaystyle \frac{|z_{\mathrm{ph}}\chi _\omega ^{}|^{2n}}{n!}}e^{|z_{\mathrm{ph}}\chi _\omega ^{}|^2}.`$ (82)
For the same reasons as those discussed in the preceding subsections we obtain the standard formula but with the cut-offs $`\chi _\omega =Z_\omega /Z`$, $`\chi _\omega ^{}=\sqrt{\chi _\omega }`$, and renormalized $`e_{\mathrm{ph}}=e_0\sqrt{Z}`$, $`z_{\mathrm{ph}}=z\sqrt{Z}`$. The rescaling $`zz_{\mathrm{ph}}`$ is exactly analogous to (24), up to the presence of the cut-off, a negligible modification if $`\chi _\omega 1`$.
We again compare theory with the experiment of the Paris group. Fig. 7 shows the predictions based on the standard (irreducible) formalism and corresponding to Fig. 2B from Haroche . We have the same problem as before: In order to fit the data we have to assume the cavity lifetime $`T_{\mathrm{cav}}=50`$ $`\mu `$s, which is much worse than in the experiment. On the contrary, the reducible formalism leads to a reasonably looking curve even for $`T_{\mathrm{cav}}=220`$ $`\mu `$s if we choose $`NZ=42`$ (Fig. 8). In these data we did not have access to the error bars. Fig. 9 and Fig. 10 show the plots corresponding to Fig. 2C in Haroche . Finally, Fig. 11 illustrates the correspondence principle: Given any data collected in a finite time interval the reducible representation is capable of reconstructing the experimental points with arbitrary precision.
## VII Conclusions
The main conclusions are the following. First of all, the evolution of atomic inversion during a finite time interval can always be reconstructed by means of some $`N<\mathrm{}`$ representation. So, finite $`N`$ representation can be regarded as a generalization of standard quantum optics. This is a consequence of the fact that there exists the correspondence principle $`N\mathrm{}`$, analogous to $`\mathrm{}0`$ or $`c\mathrm{}`$. To our surprise at least some part of the available data seems to be more consistently explained with finite $`N`$ than with the standard formalism. In order to find the proof that $`N<\mathrm{}`$ is physical, we need to find a revival of a decaying vacuum Rabi oscillation.
###### Acknowledgements.
This work is a part of the Polish Ministry of Scientific Research and Information Technology (solicited) project PZB-MIN 008/P03/2003.
## Appendix A: ‘Reducible field quantization’ in questions and answers
Here we have collected the points we think are crucial for a correct understanding of the formalism based on $`N<\mathrm{}`$ representations.
### VII.1 What if $`Z_\omega =1`$?
In principle this case is not excluded. Then the vacuum is monochromatic, i.e.
$`|O`$ $`=`$ $`|𝒑,0,`$ (83)
$`|\underset{¯}{O}`$ $`=`$ $`|𝒑,0\mathrm{}|𝒑,0,`$ (84)
$`Z=1`$, and $`\chi _\omega =1`$ (exactly). In such a vacuum a photon with momentum different from $`𝒑`$ cannot occur. In particular no resonance fluorescence is then possible. The law of large numbers is then trivial since we are dealing with the Bernoulli process with probability of success equal to 1. The parameter $`N`$ then cancels in all the formulas for $`w(t)`$. However, this state is not very physical and there is no reason to believe such a vacuum can be encountered in experiment.
### VII.2 What is the meaning of $`\chi _\omega 1`$?
What it means is that the probability $`Z(𝒌)`$, treated as a function $`𝒌Z(𝒌)`$, is very flat in some part of its domain (termed in 2 ; 3 ; 4 the quantum optics regime). Our resonant frequency is assumed to belong to the flat region of the distribution, i.e. $`\chi (𝒑)`$ is close (or even equal) to 1. The function $`\chi (𝒌)=Z(𝒌)/Z`$ has then all the features typical of the cut-off functions employed in quantum optics. However, $`Z(𝒌)`$ cannot be a constant function since then the wave function normalization condition $`_kZ(𝒌)=_k|O(𝒌)|^2=1`$ would not be fulfilled. The cut-off is a consequence of square-integrability of the wave function. Let us add that the wave function $`O(𝒌)`$ has a status analogous to the wave function of the Universe.
### VII.3 How small is $`Z`$?
We do not know. However, if differences between $`\chi (𝒌)`$ and 1 are negligible in the quantum optics regime then many different momenta are equally probable and $`Z`$ (which is the maximum of the probability) must be a very small but nonzero number. Then, on the other hand, $`N`$ must be very large, and practically the parameter that controls finite $`N`$ representations is the product $`NZ`$ which, as we have seen, may be of the order of hundreds or thousands. In a sense, the smaller $`Z`$ the better since then instead of two parameters, $`N`$ and $`Z`$, we have a single one: $`NZ`$.
### VII.4 What is the link between $`N<\mathrm{}`$ representations and oscillator wave packets?
Take a nonrelativistic oscillator with the Hamiltonian
$`H=P^2/2m+m\mathrm{\Omega }^2Q^2/2`$ (85)
and assume that $`\mathrm{\Omega }`$ is not a parameter but an operator which nevertheless commutes with $`P`$ and $`Q`$. This modification seems trivial, but is not if $`\mathrm{\Omega }`$ is nontrivial. If $`\mathrm{\Omega }=_\omega \omega I_\omega `$ is its spectral representation with spectral projectors $`I_\omega `$, then $`[I_\omega ,P]=[I_\omega ,Q]=[I_\omega ,H]=0`$. Define $`H_\omega =I_\omega H`$. Then
$`H`$ $`=`$ $`{\displaystyle \underset{\omega }{}}H_\omega ,`$ (86)
$`H_\omega `$ $`=`$ $`{\displaystyle \frac{\mathrm{}\omega }{2}}\left(a_\omega ^{}a_\omega +a_\omega a_\omega ^{}\right)`$ (87)
$`=`$ $`\mathrm{}\omega a_\omega ^{}a_\omega +{\displaystyle \frac{\mathrm{}\omega }{2}}I_\omega ,`$
$`[a_\omega ,a_\omega ^{}^{}]`$ $`=`$ $`\delta _{\omega \omega ^{}}I_\omega ,`$ (88)
$`{\displaystyle \underset{\omega }{}}I_\omega `$ $`=`$ $`I.`$ (89)
This is precisely an example of $`N=1`$ representation. Taking $`N`$ such (noninteracting!) oscillators we arrive at a $`N<\mathrm{}`$ representation. The oscillator can exist in superposition of different eigenstates of $`\mathrm{\Omega }`$. A gas consisting of $`N`$ such oscillators, all in identical ground-state wavepackets, is a Bose-Einstein condensate. But this is simultaneously our vacuum state.
### VII.5 What about Lorentz invariance?
No problem. A relativistic formalism is developed in 2 ; 3 ; 4 (although 2 is yet preliminary — one should work with 4 ). We plan to redo the Jaynes-Cummings calculations in the representation of 4 .
### VII.6 Are there modifications of the blackbody radiation?
Not really. The analysis of the problem given in 1 was premature since the role of the large-$`N`$ limit and the law of large numbers was not yet understood in this context. A single oscillator, i.e. the $`N=1`$ representation, has spectrum identical to the standard one (i.e. $`E_n=\mathrm{}\omega (n+1/2)`$; the difference is that here $`\omega `$ is also an eigenvalue and not a parameter). The field is a gas of $`N`$ such oscillators. The only subtlety is that the gas is finite and, as shown in 4 , the statistics may be based on Rényi $`\alpha `$-entropies with $`\alpha =11/N`$. The limit $`N\mathrm{}`$ is then equivalent to $`\alpha 1`$, but this is the Shannon limit of $`\alpha `$-entropies.
### VII.7 What about the divergences?
This is not yet completely clear, but all the examples discussed so far show that the cut-offs occur in the correct places. The work on full QED, loop integrals included, is in progress. There is one element that is not completely controlled yet. Namely, if we work in full space and not in a cavity (such as the classical current example discussed in 4 ) then it is natural to start with $`N=1`$ representations that involve spectral projectors $`I_\omega `$ corresponding to plane waves (i.e. these are not really projectors because of Dirac delta normalizations). Such a procedure is useful in some cases, but the final formalism should always use actual projectors in order to avoid introducing artificial infinities coming from the terms $`I_\omega ^2`$. The ‘modes’ should be associated with basis vectors in a Hilbert space, and not with plane waves. A trivial way out is to work always with finite volumes, but the formalism then lacks elegance. Some polishing of the formalism is here yet needed.
### VII.8 What about vacuum energy?
For a finite $`N`$ the vacuum part of the Hamiltonian is a well defined operator from the center of the algebra (i.e. commutes with everything). We can remove it by a unitary transformation which is well defined (this is what we implicitly do in the present paper, and sometimes refer to the procedure as a vacuum picture 2 ). The average energy of a single oscillator is finite by assumption (this is a condition on the domain of the $`N=1`$ Hamiltonian and means that $`_\omega \omega Z_\omega <\mathrm{}`$). For an arbitrary $`N`$ the vacuum energy is $`N`$ times the average energy of the $`N=1`$ case and, of course, diverges with $`N\mathrm{}`$. However, in this sense the mass of a glass of water diverges if one treats the thermodynamic limit too literally. By the way, the vacuum energy of Dirac electrons, as discussed in 3 , is negative. The discussion of vacuum energy in QED must involve both fermions and bosons, and then we have a difference of two finite expressions which may be well defined even in the limit of large $`N`$.
### VII.9 Isn’t what we do a cut-off regularization in disguise?
No, because there are no cut-offs in operators, and in the Hamiltonians in particular. It is true that in effect the end result is similar, especially in the limit $`N\mathrm{}`$. But since there is no cut-off in the Hamiltonian, there cannot be any cut-off dependence in its spectrum! This point is very important since, in principle, it can lead to yet another direct test of the $`N<\mathrm{}`$ representation. In a forthcoming paper we shall discuss the spectral properties of the Jaynes-Cummings model. The experimental aspects have been worked out by the Caltech group Kimble1 ; Kimble2 .
## Appendix B: Technicalities
### VII.10 Spectral projectors of $`a(s)^{}a(s)`$ vs. $`a(s)`$ and $`a(s)^{}`$
For $`s=0`$ we find $`a(s)=0`$. So consider $`s>0`$, $`m>0`$, and
$`\mathrm{\Pi }(m,s)a(s)\mathrm{\Pi }(n,s)`$ $`=`$ $`{\displaystyle \frac{N}{ms}}\mathrm{\Pi }(m,s)a(s)^{}a(s)a(s)\mathrm{\Pi }(n,s)`$
$`=`$ $`{\displaystyle \frac{n1}{m}}\mathrm{\Pi }(m,s)a(s)\mathrm{\Pi }(n,s).`$
Hence $`\frac{n1m}{m}\mathrm{\Pi }(m,s)a(s)\mathrm{\Pi }(n,s)=0`$ and for $`0<mn1`$ one finds
$`\mathrm{\Pi }(m,s)a(s)\mathrm{\Pi }(n,s)`$ $`=`$ $`0`$
and, by Hermitian conjugation,
$`\mathrm{\Pi }(n,s)a(s)^{}\mathrm{\Pi }(m,s)`$ $`=`$ $`0.`$
Analogously, for $`m=0`$, $`n>0`$,
$`\mathrm{\Pi }(0,s)a(s)\mathrm{\Pi }(n,s)`$ $`=`$ $`{\displaystyle \frac{1}{n}}\mathrm{\Pi }(0,s)a(s)\mathrm{\Pi }(n,s)`$
and
$`{\displaystyle \frac{n1}{n}}\mathrm{\Pi }(0,s)a(s)\mathrm{\Pi }(n,s)=0.`$
Now assume $`a(s)\mathrm{\Pi }(0,s)0`$. Then there exists a vector $`|\psi `$ such that $`a(s)\mathrm{\Pi }(0,s)|\psi =:|\mathrm{\Psi }0`$. However,
$`\mathrm{\Psi }|\mathrm{\Psi }`$ $`=`$ $`\psi |\mathrm{\Pi }(0,s)a(s)^{}a(s)\mathrm{\Pi }(0,s)|\psi =0.`$
Contradiction. Therefore $`a(s)\mathrm{\Pi }(0,s)=0`$.
### VII.11 Spectral projectors of $`X(s)`$ vs. $`\mathrm{\Omega }(s)`$
For $`n>0`$
$`[R_+a(s),\widehat{\mathrm{\Pi }}(n,s)]`$ $`=`$ $`R_+\left(a(s)\mathrm{\Pi }(n,s)\mathrm{\Pi }(n1,s)a(s)\right).`$
Employing the formulas from the previous subsection we find
$`a(s)\mathrm{\Pi }(n,s)`$ $`=`$ $`\mathrm{\Pi }(n1,s)a(s)\mathrm{\Pi }(n,s),`$
$`\mathrm{\Pi }(n1,s)a(s)`$ $`=`$ $`\mathrm{\Pi }(n1,s)a(s)\mathrm{\Pi }(n,s),`$
and $`[R_+a(s),\widehat{\mathrm{\Pi }}(n,s)]=0`$. By Hermitian conjugation $`[R_{}a(s)^{},\widehat{\mathrm{\Pi }}(n,s)]=0`$. For $`n=0`$
$`[R_+a(s),\widehat{\mathrm{\Pi }}(0,s)]`$ $`=`$ $`R_+a(s)\mathrm{\Pi }(0,s)=0.`$
By conjugation $`[R_{}a(s)^{},\widehat{\mathrm{\Pi }}(0,s)]=0`$ and, for all $`n0`$, $`[\mathrm{\Omega }(s),\widehat{\mathrm{\Pi }}(n,s)]=0`$.
### VII.12 Average $`\underset{¯}{z}|\mathrm{\Pi }(n,s)|\underset{¯}{z}`$
By construction $`a(s)|\underset{¯}{O}=0`$. Introducing the rescaled operators $`\stackrel{~}{a}(s)=a(s)/\sqrt{s/N}`$, satisfying $`[\stackrel{~}{a}(s),\stackrel{~}{a}(s)^{}]=P(s/N)`$, we check that
$`a(s)^{}a(s)\left(a(s)^{}\right)^n|\underset{¯}{O}=\left({\displaystyle \frac{s}{N}}\right)^{1+\frac{n}{2}}\stackrel{~}{a}(s)^{}\stackrel{~}{a}(s)\left(\stackrel{~}{a}(s)^{}\right)^n|\underset{¯}{O}`$
$`=\left({\displaystyle \frac{s}{N}}\right)^{1+\frac{n}{2}}n\left(\stackrel{~}{a}(s)^{}\right)^n|\underset{¯}{O}={\displaystyle \frac{s}{N}}n\left(a(s)^{}\right)^n|\underset{¯}{O}.`$
Therefore
$`\mathrm{\Pi }(n^{},s)\left(a(s)^{}\right)^n|\underset{¯}{O}`$ $`=`$ $`\delta _{n,n^{}}\left(a(s)^{}\right)^n|\underset{¯}{O},`$
$`\mathrm{\Pi }(n,s)|\underset{¯}{z}`$ $`=`$ $`{\displaystyle \frac{z^n}{n!}}e^{\frac{1}{2}|z|^2(s/N)}\left(a(s)^{}\right)^n|\underset{¯}{O},`$
and
$`\underset{¯}{z}|\mathrm{\Pi }(n,s)|\underset{¯}{z}=(s/N)^n{\displaystyle \frac{|z|^{2n}}{n!}}e^{|z|^2(s/N)}\underset{¯}{O}|P(s/N)|\underset{¯}{O}`$
$`={\displaystyle \frac{|z\sqrt{s/N}|^{2n}}{n!}}e^{|z\sqrt{s/N}|^2}\left(\begin{array}{c}N\\ s\end{array}\right)Z_\omega ^s(1Z_\omega )^{Ns}.`$ (92)
### VII.13 Explicit form of $`R_3(t)`$
We have to simplify
$`R_3(t)`$ (93)
$`=\left[\mathrm{cos}\left(\sqrt{\mathrm{\Delta }^2/4+|g|^2X}t\right)+i{\displaystyle \frac{\mathrm{sin}\left(\sqrt{\mathrm{\Delta }^2/4+|g|^2X}t\right)}{\sqrt{\mathrm{\Delta }^2/4+|g|^2X}}}\mathrm{\Omega }\right]R_3\left[\mathrm{cos}\left(\sqrt{\mathrm{\Delta }^2/4+|g|^2X}t\right)i{\displaystyle \frac{\mathrm{sin}\left(\sqrt{\mathrm{\Delta }^2/4+|g|^2X}t\right)}{\sqrt{\mathrm{\Delta }^2/4+|g|^2X}}}\mathrm{\Omega }\right].`$
One checks that $`[X,R_3]=[X,\mathrm{\Omega }]=[X,R_+a_\omega ]=[X,R_{}a_\omega ^{}]=[X,a_\omega ^{}a_\omega ]=0`$, $`[\mathrm{\Omega },R_3]=\overline{g}R_{}a_\omega ^{}gR_+a_\omega `$,
$$\mathrm{\Omega }R_3\mathrm{\Omega }=\left(\mathrm{\Delta }^2/4|g|^2X\right)R_3+\frac{\mathrm{\Delta }}{2}\left(gR_+a_\omega +\overline{g}R_{}a_\omega ^{}\right).$$
Employing these formulas we arrive at (15). |
warning/0507/hep-ph0507232.html | ar5iv | text | # 2005 International Linear Collider Workshop - Stanford, U.S.A. CERN-PH-TH/2005-120 Resolving ambiguities in mass determinations at future colliders
## I INTRODUCTION
If R-parity conserving supersymmetry exists below the TeV-scale, new particles will be produced and decay in cascades at the LHC. The lightest supersymmetric particle will escape the detectors, thereby complicating the full reconstruction of the decay chains. However, the masses of the particles in a cascade like
$$\stackrel{~}{q}_L\stackrel{~}{\chi }_2^0q\stackrel{~}{l}_{R}^{}{}_{}{}^{}l_\mathrm{n}^\pm q\stackrel{~}{\chi }_1^0l_\mathrm{f}^{}l_\mathrm{n}^\pm q$$
(1)
can be determined from endpoints of kinematical distributions Baer:1995va ; Allanach:2000kt ; Gjelsten:2004ki . Here, $`\stackrel{~}{q}_L`$ denotes a (left-handed) squark, $`\stackrel{~}{l}_R`$ a (right-handed) slepton, whereas $`\stackrel{~}{\chi }_2^0`$ and $`\stackrel{~}{\chi }_1^0`$ denote the lightest neutralinos, the latter being stable. The observed leptons and quark (jet) are denoted $`l_\mathrm{n}`$, $`l_\mathrm{f}`$ and $`q`$, where the subscripts “n” and “f” denote “near” and “far”.
Attention is focused on the mSUGRA benchmark point SPS 1a (line) Allanach:2002nj
$$m_0=A_0=0.4m_{1/2},\mathrm{tan}\beta =10,\mu >0,$$
(2)
and in particular two points on the line, SPS 1a ($`\alpha `$) with $`m_0=100\text{GeV}`$ and $`m_{1/2}=250\text{GeV}`$, and SPS 1a ($`\beta `$) with $`m_0=160\text{GeV}`$ and $`m_{1/2}=400\text{GeV}`$. The resulting low-energy masses entering in the cascade (1) are given in Table 1 Gjelsten:2004ki .
Invariant masses involving various subsets of the observed particles, namely the quark, $`q`$ and the two leptons, $`l_\mathrm{n}`$ and $`l_\mathrm{f}`$, can be reconstructed. Since the “near” and “far” leptons can not be distinguished, one must instead, on an event-by-event basis, define a “high” and “low” distribution, such that
$$m_{ql(\mathrm{low})}=\mathrm{min}(m_{ql_\mathrm{n}},m_{ql_\mathrm{f}}),m_{ql(\mathrm{high})}=\mathrm{max}(m_{ql_\mathrm{n}},m_{ql_\mathrm{f}}).$$
(3)
From a knowledge of kinematical endpoints, in particular those of the $`m_{qll}`$, $`m_{ql(\mathrm{low})}`$, $`m_{ql(\mathrm{high})}`$ and $`m_{ll}`$ distributions, the masses of the unstable particles can be reconstructed. Indeed, the endpoints of these distributions can be expressed explicitly in terms of $`m_{\stackrel{~}{q}_L}`$, $`m_{\stackrel{~}{\chi }_2^0}`$, $`m_{\stackrel{~}{l}_R}`$ and $`m_{\stackrel{~}{\chi }_1^0}`$ Allanach:2000kt ; Gjelsten:2004ki . Additional information can be obtained from threshold determinations Allanach:2000kt , and the method can be extended to the case of a gluino at the head of the chain Gjelsten:2005aw .
## II COMPOSITE FORMULAS
Many of the problems which arise in this method are related to the fact that the kinematical endpoints are composite functions of the unknown masses. Indeed, the functional form for $`m_{qll}^{\mathrm{max}}`$, $`m_{ql(\mathrm{low})}^{\mathrm{max}}`$ and $`m_{ql(\mathrm{high})}^{\mathrm{max}}`$ depend on the relative mass ratios. For $`m_{ql(\mathrm{low})}^{\mathrm{max}}`$ and $`m_{ql(\mathrm{high})}^{\mathrm{max}}`$, they are given by Allanach:2000kt ; Gjelsten:2004ki :
$$(m_{ql(\mathrm{low})}^{\mathrm{max}},m_{ql(\mathrm{high})}^{\mathrm{max}})=\left\{\begin{array}{cccc}(m_{ql_\mathrm{n}}^{\mathrm{max}},m_{ql_\mathrm{f}}^{\mathrm{max}})\hfill & \mathrm{for}\hfill & 2m_{\stackrel{~}{l}_R}^2>m_{\stackrel{~}{\chi }_1^0}^2+m_{\stackrel{~}{\chi }_2^0}^2>2m_{\stackrel{~}{\chi }_1^0}m_{\stackrel{~}{\chi }_2^0}& (\mathit{1})\\ (m_{ql(\mathrm{eq})}^{\mathrm{max}},m_{ql_\mathrm{f}}^{\mathrm{max}})\hfill & \mathrm{for}\hfill & m_{\stackrel{~}{\chi }_1^0}^2+m_{\stackrel{~}{\chi }_2^0}^2>2m_{\stackrel{~}{l}_R}^2>2m_{\stackrel{~}{\chi }_1^0}m_{\stackrel{~}{\chi }_2^0}& (\mathit{2})\\ (m_{ql(\mathrm{eq})}^{\mathrm{max}},m_{ql_\mathrm{n}}^{\mathrm{max}})\hfill & \mathrm{for}\hfill & m_{\stackrel{~}{\chi }_1^0}^2+m_{\stackrel{~}{\chi }_2^0}^2>2m_{\stackrel{~}{\chi }_1^0}m_{\stackrel{~}{\chi }_2^0}>2m_{\stackrel{~}{l}_R}^2& (\mathit{3})\end{array}\right\}$$
(4)
with
$`\left(m_{ql_\mathrm{n}}^{\mathrm{max}}\right)^2`$ $`=`$ $`\left(m_{\stackrel{~}{q}_L}^2m_{\stackrel{~}{\chi }_2^0}^2\right)\left(m_{\stackrel{~}{\chi }_2^0}^2m_{\stackrel{~}{l}_R}^2\right)/m_{\stackrel{~}{\chi }_2^0}^2`$ (5)
$`\left(m_{ql_\mathrm{f}}^{\mathrm{max}}\right)^2`$ $`=`$ $`\left(m_{\stackrel{~}{q}_L}^2m_{\stackrel{~}{\chi }_2^0}^2\right)\left(m_{\stackrel{~}{l}_R}^2m_{\stackrel{~}{\chi }_1^0}^2\right)/m_{\stackrel{~}{l}_R}^2`$ (6)
$`\left(m_{ql(\mathrm{eq})}^{\mathrm{max}}\right)^2`$ $`=`$ $`\left(m_{\stackrel{~}{q}_L}^2m_{\stackrel{~}{\chi }_2^0}^2\right)\left(m_{\stackrel{~}{l}_R}^2m_{\stackrel{~}{\chi }_1^0}^2\right)/\left(2m_{\stackrel{~}{l}_R}^2m_{\stackrel{~}{\chi }_1^0}^2\right)`$ (7)
In this report we will focus on $`m_{ql(\mathrm{low})}^{\mathrm{max}}`$ and $`m_{ql(\mathrm{high})}^{\mathrm{max}}`$ as functions of $`m_{\stackrel{~}{\chi }_1^0}`$, since the behaviour of these two endpoints is most important for the mSUGRA points studied. In the case of $`m_{ql(\mathrm{low})}^{\mathrm{max}}`$, the functional form changes when $`m_{\stackrel{~}{\chi }_1^0}`$ crosses $`\sqrt{2m_{\stackrel{~}{l}_R}^2m_{\stackrel{~}{\chi }_2^0}^2}`$, whereas for $`m_{ql(\mathrm{high})}^{\mathrm{max}}`$ it changes when $`m_{\stackrel{~}{\chi }_1^0}`$ crosses $`m_{\stackrel{~}{l}_R}^2/m_{\stackrel{~}{\chi }_2^0}`$. The three cases given in eq. (4) are in Gjelsten:2004ki referred to as regions (1,1), (1,2) and (1,3), where the first index (“1”) refers to the expression for $`m_{qll}^{\mathrm{max}}`$, which remains unchanged. The corresponding critical mass values are given in Table 2 for $`m_{\stackrel{~}{\chi }_1^0}`$, keeping the other masses at their nominal SPS 1a values.
Thus, the functional forms for $`m_{ql(\mathrm{low})}^{\mathrm{max}}`$ and $`m_{ql(\mathrm{high})}^{\mathrm{max}}`$ change very close to the nominal values of $`m_{\stackrel{~}{\chi }_1^0}`$ for SPS 1a ($`\alpha `$) and ($`\beta `$), respectively, as is illustrated in Fig. 1. These points are close to the transitions from region (1,1) to (1,2), and from region (1,2) to (1,3), respectively.
## III AMBIGUITIES
Ambiguities in the masses extracted from endpoint measurements are principally caused by the analytic form of the endpoint expressions, as seen above. In particular, they frequently lead to multiple solutions in mass-space for a particular set of endpoint values, even when experimental uncertainties are neglected. Furthermore, inclusion of extra overconstraining measurements moves solutions about in mass-space and can destroy or create new ones, adding to the complexity.
### III.1 Multiple solutions
In this subsection we restrict the discussion to the case in which the number of available (linearly independent) endpoints is equal to the number of masses involved. In particular we have in mind the standard case where $`m_{ll}^{\mathrm{max}}`$, $`m_{ql(\mathrm{low})}^{\mathrm{max}}`$, $`m_{ql(\mathrm{high})}^{\mathrm{max}}`$ and $`m_{qll}^{\mathrm{max}}`$ are measured. In this situation analytic expressions for the masses in terms of the endpoints can be obtained Gjelsten:2004ki , which in a numerical fit would correspond to solutions with $`\chi ^2=0`$. While the endpoints are obviously single-valued functions of the masses, due to the composite (and quadratic) form of the endpoint expressions the inverse is in general not true. A specific set of endpoint values can often be produced by several sets of masses. This ambiguity is illustrated in Fig. 2, where multiple solutions in the $`\chi ^2`$ fit become evident. In this fit, the four kinematical endpoints $`m_{ll}^{\mathrm{max}}`$, $`m_{ql(\mathrm{low})}^{\mathrm{max}}`$, $`m_{ql(\mathrm{high})}^{\mathrm{max}}`$ and $`m_{qll}^{\mathrm{max}}`$ are taken at their nominal values, and for each value of $`m_{\stackrel{~}{\chi }_1^0}`$ in the figure, the other three masses are allowed to vary. The left panel is devoted to SPS 1a ($`\alpha `$), whereas the right panel is devoted to SPS 1a ($`\beta `$). Three mass regions are realised; (1,1), (1,2) and (1,3), as well as the border between (1,2) and (1,3), denoted “B”. Vertical lines separate the regions in the $`m_{\mathrm{LSP}}`$ scans.
We clearly see that for four endpoint measurements, multiple mass sets provide the same endpoints and create an ambiguity at both the investigated mSUGRA points. For each of the two plots the non-smooth parts of the curve correspond to discontinuities in the mass-space trajectory defined by the (ordered) collection of best-fit masses found from the scan. At these points two identically good minima exist and a jump is made from one to the other. While the scanned mass $`m_{\stackrel{~}{\chi }_1^0}`$ obviously changes smoothly, the other masses have discontinuities at these points. For the two investigated mSUGRA scenarios and in particular SPS 1a $`(\beta )`$, $`m_{\stackrel{~}{l}_R}`$ and $`m_{\stackrel{~}{q}_L}`$, more so than $`m_{\stackrel{~}{\chi }_2^0}`$, change considerably at these discontinuities, the first by roughly 5%.
### III.2 Overconstrained systems
When an extra endpoint measurement, such as the $`qll`$ “threshold”, is introduced, the system becomes overconstrained, with five measurements determining only four masses. One might expect this to alleviate the ambiguity from multiple solutions seen in the previous section, since the extra endpoint should ‘pick out’ one of the solutions. However, this is not always the case: the uncertainty of endpoint measurements and the introduction of an overconstraining measurement will move the solutions around in mass-space, possibly creating new solutions, or destroying old ones. This may then cause a further ambiguity, additional to that seen in III.1, and actually does so for SPS 1a $`(\beta )`$. We will demonstrate the principles involved in the creation of new solutions by using a simplified case where only one mass is unknown and two endpoints are measured.
For our simplified example, let us keep all masses other than $`m_{\stackrel{~}{\chi }_1^0}`$ fixed at their nominal values, but let the two endpoint values be offset from their nominal values, as is easily imagined due to statistical fluctuations. We show in Fig. 3 a simplified $`\chi ^2`$ function:
$$\chi ^2=a(m_{ql(\mathrm{low})}^{\mathrm{max}}m_{ql(\mathrm{low})}^{\mathrm{max}}{}_{\mathrm{exp}}{}^{})^2+b(m_{ql(\mathrm{high})}^{\mathrm{max}}m_{ql(\mathrm{high})}^{\mathrm{max}}{}_{\mathrm{exp}}{}^{})^2$$
(8)
for SPS 1a ($`\beta `$), where two cases of “experimental” data are considered, nominal values (left panel) and off-set values (right panel). The coefficients $`a`$ and $`b`$ represent the experimental uncertainties, $`\sigma (m_{ql(\mathrm{low})}^{\mathrm{max}})=6.3`$ GeV, and $`\sigma (m_{ql(\mathrm{high})}^{\mathrm{max}})=5.5`$ GeV Gjelsten:2004ki . Individual $`\chi ^2`$-values from $`m_{ql(\mathrm{low})}^{\mathrm{max}}`$ and $`m_{ql(\mathrm{high})}^{\mathrm{max}}`$ are shown (labeled $`\chi ^2(\mathrm{low})`$ and $`\chi ^2(\mathrm{high})`$), together with the sum. In this example, there is only one solution if the endpoints take their nominal values, despite the system being overconstrained and in spite of the compositeness of $`\chi ^2`$ (see left panel).
However, if the endpoint measurements give values which differ from the nominal ones, as will in general be the case, the situation can be more complicated, as is shown in the right panel of Fig. 3. Here, for the purpose of illustration, we consider $`m_{ql(\mathrm{low})}^{\mathrm{max}}`$ and $`m_{ql(\mathrm{high})}^{\mathrm{max}}`$ offset from their nominal values by $`20`$ GeV and $`+15`$ GeV, respectively. A secondary minimum has now emerged, at a higher value of $`m_{\stackrel{~}{\chi }_1^0}`$.
If one allows also the other masses ($`m_{\stackrel{~}{q}_L}`$, $`m_{\stackrel{~}{\chi }_2^0}`$, $`m_{\stackrel{~}{l}_R}`$) to vary, the transition between regions (1,2) and (1,3) will move, and secondly, the minima will get deeper. In a more realistic analysis, however, the additional measurements of $`m_{qll}^{\mathrm{max}}`$ and, in particular $`m_{ll}^{\mathrm{max}}`$, will disfavor the second minimum. Indeed, with the additional information of the $`qll`$ “threshold”, this ambiguity was observed in the more detailed analysis reported in Gjelsten:2004ki . In Tables 34 we report on the probability of having false minima, and corresponding mass values, for SPS 1a $`(\alpha )`$ and ($`\beta `$).
## IV LINEAR COLLIDER INPUT
With Linear Collider input on the LSP mass, taken to be determined with a precision of 0.05 GeV Aguilar-Saavedra:2001rg , the ambiguities are mostly resolved. For SPS 1a $`(\alpha )`$ the false solution is removed altogether. The masses and their errors are given in Table 5. For SPS 1a $`(\beta )`$ the high-mass solution no longer appears, but there is still an ambiguity in the low-mass region, see Table 6. These masses are however very close. Furthermore, as a Linear Collider will measure the slepton mass with an accuracy similar to that of the LSP mass Aguilar-Saavedra:2001rg , one may expect the remaining ambiguity to vanish as well. This possibility was not considered in the analyses Gjelsten:2004ki ; Gjelsten:2005aw and hence is not reflected in Tables 56.
A second effect, not discussed above, is very clear by comparison of the errors on the masses obtained with and without the LC measurement: Fixing the LSP mass strongly reduces the errors on all the masses. This comes from the fact that the endpoint method determines mass differences much better than the actual masses themselves, a feature due to the way masses enter the endpoint expressions; in terms of differences of masses (squared). Actually, the errors on the masses obtained when the LHC and the LC are combined, Tables 56, are very close to the errors on the mass differences as obtained by the LHC alone.
The removal of ambiguities, and the higher precision are crucial for an extrapolation to the GUT scale Allanach:2004ud , a major goal of studying the spectroscopy of supersymmetric particles at future accelerators.
###### Acknowledgements.
This research has been supported in part by the Research Council of Norway. |
warning/0507/hep-th0507140.html | ar5iv | text | # Contents
## 1 Introduction
Domain walls in gauge theories with eight supercharges have rather special properties. These walls were first studied by Abraham and Townsend who showed that in two-dimensions, where domain walls are known as kinks, they exhibit dyonic behaviour reminiscent of magnetic monopoles. Further similarities between kinks and magnetic monopoles, at both the classical and quantum level, were uncovered in . The physical explanation for this relationship was presented in , where new BPS solutions were described corresponding to magnetic monopoles in a phase with fully broken gauge symmetry. The Meissner effect ensures that monopoles are confined: the magnetic flux is unable to propagate through the vacuum and leaves the monopole in two collimated tubes. From the perspective of the flux tube, the monopole appears as a kink. The idea of describing confined monopoles as kinks in $`Z_N`$ strings occurred previously in . The relationship between the confined magnetic monopoles and the kink was further explored in and related systems were studied in .
In this paper we use D-brane techniques to study the moduli space of multiple domain walls. This allows us to develop a description of the domain wall dynamics in terms of a linearized Nahm equation, providing a direct relationship to the dynamics of monopoles. Specifically, we show that the moduli space of domain walls, which we denote as $`𝒲_\stackrel{}{g}`$, is isomorphic to a middle dimensional submanifold of the moduli space of magnetic monopoles $`_\stackrel{}{g}`$. This submanifold describes magnetic monopoles lying along a line, and can be described as the fixed point of an $`𝐒^1`$ action $`\widehat{k}`$, rotating the monopoles in a plane,
$`𝒲_\stackrel{}{g}_\stackrel{}{g}|_{\widehat{k}=0}`$ (1.1)
The correspondence captures the topology and asymptotic metric of the domain wall moduli space $`𝒲_\stackrel{}{g}`$. It does not extend to the full metric on $`𝒲_\stackrel{}{g}`$. Nevertheless, as we shall explain, it does correctly capture the most important feature of domain walls: their ordering along the line.
The relationship (1.1) plays companion to the result of , where the moduli space of vortices was shown to be a middle dimensional submanifold of the moduli space of instantons. Indeed, upon dimensional reduction, the self-dual instanton equations become the monopole equations, while the vortex equations descend to the domain wall equations.
We start in the following section by describing the domain walls in question, together with a review of their moduli spaces. We pay particular attention to the crudest physical feature of domain walls, namely the rules governing their spatial ordering along the line. Section 3 contains a brief review of magnetic monopoles in higher rank gauge groups, primarily in order to fix notation, allowing us to elaborate on the relationship (1.1). We also describe the Nahm construction of the monopole moduli space as it arises from D-branes. The meat of the paper is in Section 4. We present a D-brane embedding of domain wall solitons which gives a description of their dynamics in terms of a linear Nahm equation. This equation is somewhat trivial, with the content hidden in various boundary conditions. We show how these boundary conditions capture the prescribed ordering of domain walls.
## 2 Domain Walls
In this paper we will study a class of BPS domain wall solutions occurring in maximally supersymmetric theories with multiple, isolated vacua. The Lagrangian for these models includes a $`U(k)`$ gauge field $`A_\mu `$, a real adjoint scalar $`\sigma `$ and $`N`$ fundamental scalars $`q_i`$, each with real mass $`m_i`$
$`=\mathrm{Tr}\left[{\displaystyle \frac{1}{4e^2}}F_{\mu \nu }F^{\mu \nu }+{\displaystyle \frac{1}{2e^2}}|𝒟_\mu \sigma |^2+{\displaystyle \frac{e^2}{2}}(q_iq_i^{}v^2)^2\right]+{\displaystyle \underset{i=1}{\overset{N}{}}}\left[|𝒟_\mu q_i|^2+q_i^{}(\sigma m_i)^2q_i\right]`$
where there is an implicit sum over the flavor index $`i`$ in the adjoint valued term $`q_iq_i^{}`$. This Lagrangian can be embedded in a theory with 8 supercharges in any spacetime dimension $`1d5`$ (e.g. $`𝒩=2`$ SQCD in $`d=3+1`$). Such theories include further scalar fields which can be shown to vanish on the domain wall solutions<sup>1</sup><sup>1</sup>1If we promote the scalar field $`\sigma `$ and the masses $`m_i`$ to complex variables, then the theories admit an interesting array of domain wall junctions and dyonic walls .. The fermions do contribute zero modes but will not be important here.
When the Higgs expectation value $`v^2`$ is non-vanishing, and the masses $`m_i`$ are distinct ($`m_im_j`$ for $`ij`$), the theory has a set of isolated vacua. Each vacuum is labelled by a set $`\mathrm{\Xi }`$ of $`k`$ distinct elements, chosen from a possible $`N`$,
$`\mathrm{\Xi }=\{\xi (a):\xi (a)\xi (b)\mathrm{for}ab\}`$ (2.1)
Here $`a=1,\mathrm{},k`$ runs over the color index, while $`\xi (a)\{1,\mathrm{},N\}`$. Up to a gauge transformation, the vacuum associated to this set is given by,
$`\sigma =\mathrm{diag}(m_{\xi (1)},\mathrm{},m_{\xi (k)}),q_i^a=v\delta _{i=\xi (a)}^a`$ (2.2)
For $`N<k`$ there are no supersymmetric vacua; for $`Nk`$, the number of vacua is
$`N_{\mathrm{vac}}=\left({\displaystyle \genfrac{}{}{0pt}{}{N}{k}}\right)={\displaystyle \frac{N!}{k!(Nk)!}}`$ (2.3)
Each of these vacua is isolated and exhibits a mass gap. There are $`k^2`$ non-BPS massive gauge bosons and quarks with masses $`m_\gamma ^2e^2v^2+|m_{\xi (a)}m_{\xi (b)}|^2`$ and $`k(Nk)`$ BPS massive quark fields with masses $`m_q|m_{\xi (a)}m_i|`$ (with $`i\mathrm{\Xi }`$).
For vanishing masses $`m_i=0`$ the theory enjoys an $`SU(N)`$ flavor symmetry, rotating the $`q_i`$. When distinct masses are turned on this is broken explicitly to the Cartan-sub-algebra $`U(1)^{N1}`$. Meanwhile, the $`U(k)`$ gauge group is broken spontaneously in the vacuum by the expectation values (2.2).
The existence of multiple, gapped, isolated vacua is sufficient to guarantee the existence of co-dimension one domain walls (otherwise known as kinks). These walls are BPS objects, satisfying first order Bogomoln’yi equations which can be derived in the usual manner by completing the square. We first choose a flat connection $`F_{\mu \nu }=0`$, and allow the fields to depend only on a single coordinate, say $`x^3`$. Then the Hamiltonian is given by
$``$ $`=`$ $`\mathrm{Tr}\left[{\displaystyle \frac{1}{2e^2}}|𝒟_3\sigma |^2+{\displaystyle \frac{e^2}{2}}(q_iq_i^{}v^2)^2\right]+{\displaystyle \underset{i=1}{\overset{N}{}}}\left[|𝒟_3q_i|^2+q_i^{}(\sigma m_i)^2q_i\right]`$
$`=`$ $`{\displaystyle \frac{1}{2e^2}}\mathrm{Tr}\left[𝒟_3\sigma e^2(q_iq_i^{}v^2)\right]^2+{\displaystyle \underset{i=1}{\overset{N}{}}}|𝒟_3q_i(\sigma m_i)q_i|^2`$
$`+\mathrm{Tr}\left[𝒟_3\sigma (q_iq_i^{}v^2)\right]+{\displaystyle \underset{i=1}{\overset{N}{}}}\left[q_i^{}(\sigma m_i)𝒟_3q_i+𝒟_3q_i^{}(\sigma m_i)q_i\right]`$
$``$ $`_3\left(v^2\mathrm{Tr}\sigma \right)`$
Our domain wall interpolates between a vacuum $`\mathrm{\Xi }_{}`$ at $`x^3=\mathrm{}`$, as determined by a set (2.1), and a distinct vacuum $`\mathrm{\Xi }_+`$ at $`x^3=+\mathrm{}`$. The minus signs above have been chosen under the assumption that $`\mathrm{\Delta }m>0`$, where $`\mathrm{\Delta }m=_{i\mathrm{\Xi }_{}}m_i_{i\mathrm{\Xi }_+}m_i`$, so that a BPS domain wall satisfies the Bogomoln’yi equations,
$`𝒟_3\sigma =e^2(q_iq_i^{}v^2),𝒟_3q_i=(\sigma m_i)q_i`$ (2.5)
and has tension given by $`T=v^2\mathrm{\Delta }m`$. Analytic solutions to these equations can be found in the $`e^2\mathrm{}`$ limit , which give smooth approximations to the solution at large, but finite $`e^2`$ .
### 2.1 Classification of Domain Walls
Domain walls in field theories are classified by the choice of vacuum $`\mathrm{\Xi }_{}`$ and $`\mathrm{\Xi }_+`$ at left and right infinity. However, our theory contains an exponentially large number of vacua (2.3) and one may hope that there is a coarser, less unwieldy, classification which captures certain relevant properties of a given domain wall. Such a classification was offered in .
Firstly define the $`N`$-vector $`\stackrel{}{m}=(m_1,\mathrm{},m_N)`$. The tension of the BPS domain wall can then be written as
$`T_\stackrel{}{g}=v^2\mathrm{\Delta }mv^2\stackrel{}{m}\stackrel{}{g}`$ (2.6)
where the $`N`$-vector $`\stackrel{}{g}\mathrm{\Lambda }_R(su(N))`$, the root lattice of $`su(N)`$. Note that there do not exist domain wall solutions for all $`\stackrel{}{g}\mathrm{\Lambda }_R(su(N))`$; the only admissible vectors are of the form $`\stackrel{}{g}=(p_1,\mathrm{},p_N)`$ with $`p_i=0`$ or $`\pm 1`$. Note also that a choice of $`\stackrel{}{g}`$ does not specify a unique choice of vacua $`\mathrm{\Xi }_{}`$ and $`\mathrm{\Xi }_+`$ at left and right infinity. Nor, in fact, does it specify a unique domain wall moduli space $`𝒲_\stackrel{}{g}`$. Nevertheless, domain walls in sectors with the same $`\stackrel{}{g}`$ share common traits.
The dimension of the moduli space of domain wall solutions was computed in using an index theorem, following earlier results in . To describe the dimension of the moduli space, it is useful to decompose $`\stackrel{}{g}`$ in terms of simple roots<sup>1</sup><sup>1</sup>1The basis of simple roots is fixed by the requirement that $`\stackrel{}{m}\stackrel{}{\alpha }_i>0`$ for each $`i`$. A unique basis is defined in this way if $`\stackrel{}{m}`$ lies in a positive Weyl chamber, which occurs whenever the masses are distinct so that $`SU(N)U(1)^{N1}`$. If we choose the ordering $`m_1>m_2>\mathrm{}>m_N`$ we have simple roots $`\stackrel{}{\alpha }_1=(1,1,0,\mathrm{},0)`$ and $`\stackrel{}{\alpha }_2=(0,1,1,0,\mathrm{},0)`$ through to $`\stackrel{}{\alpha }_{N1}=(0,\mathrm{},1,1)`$. $`\stackrel{}{\alpha }_i`$,
$`\stackrel{}{g}={\displaystyle \underset{i=1}{\overset{N1}{}}}n_i\stackrel{}{\alpha }_i,n_i𝐙`$ (2.7)
The index theorem of reveals that domain wall solutions to (2.5) exist only if $`n_i0`$ for all $`i`$. If this holds, the number of zero modes of a solution is given by
$`\mathrm{dim}\left(𝒲_\stackrel{}{g}\right)=2{\displaystyle \underset{i=1}{\overset{N1}{}}}n_i`$ (2.8)
where $`𝒲_\stackrel{}{g}`$ denotes the moduli space of any set of domain walls with charge $`\stackrel{}{g}`$. This result has a simple physical interpretation. There exist $`N1`$ types of “elementary” domain walls corresponding to a $`\stackrel{}{g}=\stackrel{}{\alpha }_i`$, the simple roots. Each of these has just two collective coordinates corresponding to a position in the $`x^3`$ direction and a phase. As first explained in , the phase coordinate is a Goldstone mode arising because the domain wall configuration breaks the $`U(1)^{N1}`$ flavor symmetry as we review below. In general, a domain wall labelled by $`\stackrel{}{g}`$ can be thought to be constructed from $`_in_i`$ elementary domain walls, each with its own position and phase collective coordinate. Let us now turn to some examples.
### 2.2 The Moduli Space of Domain Walls: Some Examples
#### Example 1: $`\stackrel{}{g}=\stackrel{}{\alpha }_1`$
The simplest system admitting a domain wall is the abelian $`k=1`$ theory with $`N=2`$ charged scalars $`q_1`$ and $`q_2`$. The $`N_{vac}=2`$ vacua of the theory are given by $`\sigma =m_i`$ and $`|q_j|^2=v^2\delta _{ij}`$ for $`i=1,2`$. There is a single domain wall in this theory with $`\stackrel{}{g}=\stackrel{}{\alpha }_1`$ interpolating between these two vacua. Under the $`U(1)_F`$ flavor symmetry, $`q_1`$ has charge $`+1`$ and $`q_2`$ has charge $`1`$. In each of the vacua, the $`U(1)_F`$ symmetry coincides with the $`U(1)`$ gauge action but, in the core of the domain wall, both $`q_1`$ and $`q_2`$ are non-vanishing, and $`U(1)_F`$ acts non-trivially. The resulting goldstone mode is the phase collective coordinate. The moduli space of the domain wall is simply
$`𝒲_{\stackrel{}{g}=\stackrel{}{\alpha }_1}𝐑\times 𝐒^1`$ (2.9)
where the $`𝐑`$ factor describes the center of mass of the domain wall and the $`𝐒^1`$ the phase. One can show that the $`𝐒^1`$ has radius $`2\pi v^2/T_\stackrel{}{g}=2\pi /(m_1m_2)`$.
#### Example 2: $`\stackrel{}{g}=\stackrel{}{\alpha }_1+\stackrel{}{\alpha }_2`$
The simplest system admitting multiple domain walls is the abelian $`k=1`$ theory with $`N=3`$ charged scalars. There are now three vacua, given by $`\sigma =m_i`$ and $`|q_j|^2=v^2\delta _{ij}`$. In each a $`U(1)_{F_1}\times U(1)_{F_2}`$ flavor symmetry is unbroken, under which the scalars have charge
$`U(1)_{F_1}\times U(1)_{F_2}:\{\begin{array}{cc}q_1:& (+1,0)\hfill \\ q_2:& (1,1)\hfill \\ q_3:& (0,1)\hfill \end{array}`$ (2.13)
The first elementary domain wall $`\stackrel{}{g}=\stackrel{}{\alpha }_1`$ interpolates between vacuum 1 and vacuum 2, breaking $`U(1)_{F_1}`$ along the way. The second elementary domain wall interpolates between vacuum 2 and vacuum 3, breaking $`U(1)_{F_2}`$ along the way. Of interest here is the domain wall $`\stackrel{}{g}=\stackrel{}{\alpha }_1+\stackrel{}{\alpha }_2`$ interpolating between vacuum 1 and vacuum 3. It can be thought of as a composite of two domain walls. The moduli space for these two domain walls was studied in and is of the form,
$`𝒲_{\stackrel{}{g}=\stackrel{}{\alpha }_1+\stackrel{}{\alpha }_2}𝐑\times {\displaystyle \frac{𝐑\times \stackrel{~}{𝒲}_{\stackrel{}{\alpha }_1+\stackrel{}{\alpha }_2}}{𝐙}}`$ (2.14)
The first factor of $`𝐑`$ corresponds to the center of mass of the two domain walls; the second factor corresponds to the combined phase associated to the two domain walls. When the ratio of tensions of the two elementary domain walls $`T_{\stackrel{}{\alpha }_1}/T_{\stackrel{}{\alpha }_2}`$ is rational, the ratio of the periods of the two phases are similarly rational and the second $`𝐑`$ factor collapses to $`𝐒^1`$, while the quotient $`𝐙`$ reduces to a finite group. The relative moduli space $`\stackrel{~}{𝒲}_{\stackrel{}{\alpha }_1+\stackrel{}{\alpha }_2}`$ corresponds to the separation and relative phases of the two domain walls. Importantly, and unlike other solitons of higher co-dimension, the domain walls must obey a strict ordering on the $`x^3`$ line: the $`\stackrel{}{g}=\stackrel{}{\alpha }_1`$ domain wall must always be to the left of $`\stackrel{}{g}=\stackrel{}{\alpha }_2`$ domain wall. The separation between the walls is therefore the halfline $`𝐑^+`$. It was shown in that the relative phase is fibered over $`𝐑^+`$ to give rise to a smooth cylinder, with the tip corresponding to coincident domain walls. The resulting moduli space is shown in Figure 1.
Note that the moduli space (2.14) is toric, inheriting two isometries from the $`U(1)_{F_1}\times U(1)_{F_2}`$ symmetry. In an abelian gauge theory with arbitrary number of flavors $`N`$, the domain wall charge is always of the form $`\stackrel{}{g}=_in_i\stackrel{}{\alpha }_i`$ with $`n_i=0,1`$, and the moduli space is always toric, meaning that half of the dimensions correspond to $`U(1)`$ isometries.
#### Example 3: $`\stackrel{}{g}=\stackrel{}{\alpha }_1+2\stackrel{}{\alpha }_2+\stackrel{}{\alpha }_3`$
In non-abelian theories, the domain wall moduli spaces are no longer toric. The simplest such theory has a $`U(2)`$ gauge group with $`N=4`$ fundamental scalars. The 6 vacua, and 15 different domain walls, of this theory were detailed in . Under the $`U(1)_F^3`$ flavor symmetry, the fundamental scalars transform as
$`U(1)_{F_1}\times U(1)_{F_2}\times U(1)_{F_3}:\{\begin{array}{cc}q_1:& (+1,0,0)\hfill \\ q_2:& (1,1,0)\hfill \\ q_3:& (0,1,1)\hfill \\ q_4:& (0,0,1)\hfill \end{array}`$ (2.19)
With this convention, the elementary domain wall $`\stackrel{}{g}=\stackrel{}{\alpha }_i`$ picks up its phase from the action of the $`U(1)_{F_i}`$ flavor symmetry.
Here we concentrate on the domain wall system with the maximal number of zero modes which arises from the choice of vacua $`\mathrm{\Xi }_{}=(1,2)`$ and $`\mathrm{\Xi }_+=(3,4)`$ so that $`\stackrel{}{g}=\stackrel{}{\alpha }_1+2\stackrel{}{\alpha }_2+\stackrel{}{\alpha }_3`$. This system can be separated into four constituent domain walls. As explained in , the ordering of domain walls is no longer strictly fixed in this example. The two outer elementary domain walls, on the far left and far right, are each of the type $`\stackrel{}{g}=\stackrel{}{\alpha }_2`$. However, the relative positions of the middle two domain walls, $`\stackrel{}{g}=\stackrel{}{\alpha }_1`$ and $`\stackrel{}{\alpha }_3`$ are not ordered and they may pass through each other.
Unlike the situation for abelian gauge theories, the 8 dimensional domain wall moduli space for this example is no longer toric; $`𝒲_\stackrel{}{g}`$ inherits only three $`U(1)`$ isometries from (2.19). Physically this means that the two phases associated to the $`\stackrel{}{\alpha }_2`$ domain walls are not both Goldstone modes and they may interact as the domain walls approach. This behaviour is familiar from the study of the Atiyah-Hitchin metric describing the dynamics of two monopoles in $`SU(2)`$ gauge theory; we shall make the analogy more precise in the following.
### 2.3 The Ordering of Domain Walls
As we stressed in the above examples, in contrast to other solitons domain walls must obey some ordering on the line. This will be an important ingredient when we come to extract domain wall data from the linearized Nahm’s equations in Section 4. Here we linger to review this ordering.
The ordering of the elementary domain walls in non-abelian theories was studied in detail in . One can derive the result by considering the possible sequences of vacua as we move over each domain wall. For example, we could consider the “maximal domain wall”, interpolating between $`\mathrm{\Xi }_{}=\{1,2,\mathrm{},k\}`$ and $`\mathrm{\Xi }_+=\{Nk+1,\mathrm{},N\}`$. From the left, the first elementary domain wall that we come across must be $`\stackrel{}{g}=\stackrel{}{\alpha }_k`$, corresponding to $`\mathrm{\Xi }_{}=(1,2,\mathrm{},k1,k)(1,2,\mathrm{},k1,k+1)`$. The next elementary domain wall may be either $`\stackrel{}{\alpha }_{k1}`$ or $`\stackrel{}{\alpha }_{k+1}`$. These two walls are free to pass through each other, but cannot move further to the left than the $`\stackrel{}{\alpha }_k`$ wall. And so on. Iterating this procedure, one finds that two neighbouring elementary domain walls $`\stackrel{}{\alpha }_i`$ and $`\stackrel{}{\alpha }_j`$ may pass through each other whenever $`\stackrel{}{\alpha }_i\stackrel{}{\alpha }_j=0`$, but otherwise have a fixed ordering on the line. The net result of this analysis is summarized in Figure 2. The $`x^3`$ positions of the domain walls are shown on the vertical axis; the position on the horizontal axis denotes the type of elementary domain wall, starting on the left with $`\stackrel{}{\alpha }_1`$ and ending on the right with $`\stackrel{}{\alpha }_{N1}`$.
In summary, we see that for $`n_i=n_{i+1}1`$, the $`\stackrel{}{\alpha }_i`$ domain walls are trapped between the $`\stackrel{}{\alpha }_{i+1}`$ domain walls. The reverse holds when $`n_i=n_{i+1}+1`$. Finally, when $`n_i=n_{i+1}`$ the positions of the $`\stackrel{}{\alpha }_i`$ domain walls are interlaced with those of the $`\stackrel{}{\alpha }_{i+1}`$ domain walls. The last $`\stackrel{}{\alpha }_{i+1}`$ domain wall is unconstrained by the $`\stackrel{}{\alpha }_i`$ walls in its travel in the positive $`x^3`$ direction, although it may be trapped in turn by a $`\stackrel{}{\alpha }_{i+2}`$ wall.
While we have discussed the maximal domain wall above, other sectors can be reached either by removing some of the outer domain walls to infinity, or by taking non-interacting products such subsets. It’s important to note that labelling a topological sector $`\stackrel{}{g}`$ does not necessarily determine the ordering of domain walls<sup>1</sup><sup>1</sup>1An example: the $`\stackrel{}{g}=\stackrel{}{\alpha }_1+\stackrel{}{\alpha }_2+\stackrel{}{\alpha }_3`$ domain wall. In the abelian theory with $`k=1`$ and $`N=4`$, the ordering is $`\stackrel{}{\alpha }_1<\stackrel{}{\alpha }_2<\stackrel{}{\alpha }_3`$. However, in the non-abelian theory with $`k=2`$ and $`N=4`$, the ordering is $`\stackrel{}{\alpha }_1,\stackrel{}{\alpha }_3<\stackrel{}{\alpha }_2`$.. We shall show that domain walls with the same $`\stackrel{}{g}`$, but different orderings, descend from different submanifolds of the same monopole moduli space.
## 3 Monopoles
The main goal of this paper is to show how the moduli space of domain walls introduced in the previous section is isomorphic to a submanifold of a related monopole moduli space. In this section we review several relevant aspects of these monopole moduli spaces.
It will turn out that the domain walls of Section 2 are related to monopoles in an $`SU(N)`$ gauge theory. Note that the flavor group from Section 2 has been promoted to a gauge group; we shall see the reason behind this in Section 5. The Bogomoln’yi monopole equations are
$`B_\mu =𝒟_\mu \varphi `$ (3.1)
where $`B_\mu `$, $`\mu =1,2,3`$ is the $`SU(N)`$ magnetic field and $`\varphi `$ is an adjoint valued real scalar field. The monopoles exist only if $`\varphi `$ takes a vacuum expectation value,
$`\varphi =\mathrm{diag}(m_1,\mathrm{},m_N)`$ (3.2)
where we take $`m_im_j`$ for $`ij`$, ensuring breaking to the maximal torus, $`SU(N)U(1)^{N1}`$. It is not coincidence that we’ve denoted the vacuum expectation values by $`m_i`$, the same notation used for the masses in Section 2; it is for this choice of vacuum that the correspondence holds. (Specifically, the masses of the kinks will coincide with the masses of monopoles, ensuring that the asymptotic metrics on $`𝒲_\stackrel{}{g}`$ and $`_\stackrel{}{g}`$ also coincide).
As described long ago by Goddard, Nuyts and Olive , the allowed magnetic charges under each unbroken $`U(1)^{N1}`$ are specified by a root vector<sup>1</sup><sup>1</sup>1We ignore the factor of 2 difference between roots and co-roots. For simply laced groups, such as $`SU(N)`$, it can be absorbed into convention. of $`su(N)`$, $`\stackrel{}{g}=(p_1,\mathrm{},p_N)`$. It is customary to decompose this in terms of simple roots $`\stackrel{}{\alpha }_i`$,
$`\stackrel{}{g}={\displaystyle \underset{i=1}{\overset{N1}{}}}n_i\stackrel{}{\alpha }_i,n_i𝐙`$ (3.3)
Once again, the notation is identical to that used for domain walls (2.7) for good reason. Solutions to the monopole equations (3.1) exist for all values of $`n_i0`$. This is in contrast to domain walls where, as we have seen, configurations only exist in a finite number of sectors defined by $`p_i=0`$ or $`p_i=\pm 1`$. The mass of the magnetic monopole is $`M_{\mathrm{mono}}=(2\pi /e^2)\stackrel{}{m}\stackrel{}{g}`$.
The monopole moduli space $`_\stackrel{}{g}`$ is the space of solutions to (3.2) in a fixed topological sector $`\stackrel{}{g}`$. The dimension of this space, equal to the number of zero modes of given solution, was computed by E. Weinberg in using Callias’ version of the index theorem. The result is:
$`\mathrm{dim}(_\stackrel{}{g})=4{\displaystyle \underset{i=1}{\overset{N1}{}}}n_i`$ (3.4)
which is to be compared with (2.8).
### 3.1 The Relationship between $`_\stackrel{}{g}`$ and $`𝒲_\stackrel{}{g}`$
We are now in a position to describe the relationship between the moduli space of domain walls $`𝒲_\stackrel{}{g}`$ and the moduli space of magnetic monopoles $`_\stackrel{}{g}`$. We will show that $`𝒲_\stackrel{}{g}`$ is a complex, middle dimensional, submanifold of $`_\stackrel{}{g}`$, defined by the fixed point set of the action rotating the monopoles in a plane, together with a suitable gauge action. To do this, we first need to describe the symmetries of $`_\stackrel{}{g}`$.
The monopole moduli space $`_\stackrel{}{g}`$ admits a natural, smooth, hyperKähler metric . For generic $`\stackrel{}{g}`$ this metric enjoys $`(N1)`$ tri-holomorphic isometries arising from the action of the $`U(1)^{N1}`$ abelian gauge group. Further the metric has an $`SU(2)_R`$ symmetry, arising from rotations of the monopoles in $`𝐑^3`$, which acts on the three complex structures of $`_\stackrel{}{g}`$. In other words, any $`U(1)_RSU(2)_R`$ is a holomorphic isometry, preserving a single complex structure while revolving the remaining two. Let us choose $`U(1)_R`$ to rotate the monopoles in the $`(x^2x^3)`$ plane. In what follows we will be interested in a specific holomorphic $`\widehat{U}(1)`$ action which acts simultaneously by a $`U(1)_R`$ rotation and a linear combination of the gauge rotations $`U(1)^{N1}`$ (to be specified presently). We denote the Killing vector on $`_\stackrel{}{g}`$ associated to $`\widehat{U}(1)`$ as $`\widehat{k}`$. We claim
$`𝒲_\stackrel{}{g}=_\stackrel{}{g}|_{\widehat{k}=0}`$ (3.5)
This result holds at the level of topology and asymptotic metric of the spaces. The manifold $`𝒲_\stackrel{}{g}`$ inherits a metric from $`_\stackrel{}{g}`$ by this reduction: it does not coincide with the domain wall metric in the interior on $`𝒲_\stackrel{}{g}`$. (For example, corrections to the asymptotic metric on $`𝒲_\stackrel{}{g}`$ are exponentially suppressed while those of $`_\stackrel{}{g}`$ have power law behaviour). It would be interesting to examine if $`𝒲_\stackrel{}{g}`$ inherits the correct Kähler class and/or complex structure from $`_\stackrel{}{g}`$.
We defer a derivation of (3.5) to the following section, but first present some simple examples.
#### Example 1: $`\stackrel{}{g}=\stackrel{}{\alpha }_1`$
Monopoles in $`SU(2)`$ gauge theories are labelled by a single topological charge $`\stackrel{}{g}=n_1\stackrel{}{\alpha }_1`$. For a single monopole ($`n_1=1`$) the moduli space is simply
$`_{\stackrel{}{g}=\stackrel{}{\alpha }_1}𝐑^3\times 𝐒^1`$ (3.6)
where the $`𝐑^3`$ factor denotes the position of the monopole, while the $`𝐒^1`$ arises from global gauge transformations under the surviving $`U(1)`$. The radius of the $`𝐒^1`$ is $`2\pi /(m_1m_2)`$. In this case the $`\widehat{U}(1)`$ action coincides with the rotation $`U(1)_R`$ in the $`(x^2x^3)`$ plane and we have trivially
$`𝒲_{\stackrel{}{\alpha }_1}𝐑\times 𝐒^1_{\stackrel{}{\alpha }_1}|_{\widehat{k}=0}`$ (3.7)
The similarity between the domain wall and monopole moduli spaces for a single soliton was noted by Abraham and Townsend . In both cases, motion in the $`𝐒^1`$ factor gives rise to dyonic solitons.
Note that monopole moduli spaces for charges $`\stackrel{}{g}=n_1\stackrel{}{\alpha }_1`$ exist for all $`n_1𝐙^+`$. For example, the $`n_1=2`$ monopole moduli space is home to the famous Atiyah-Hitchin metric . However, there is no domain wall moduli space with this charge in the class of theories we discuss in Section 2.
#### Example 2: $`\stackrel{}{g}=\stackrel{}{\alpha }_1+\stackrel{}{\alpha }_2`$
Our second example is the $`\stackrel{}{g}=\stackrel{}{\alpha }_1+\stackrel{}{\alpha }_2`$ monopole in $`SU(3)`$ gauge theories (sometimes referred to as the $`(1,1)`$ monopole). The moduli space was determined in to be of the form
$`_{\stackrel{}{g}=\stackrel{}{\alpha }_1+\stackrel{}{\alpha }_2}𝐑^3\times {\displaystyle \frac{𝐑\times \stackrel{~}{}_{\stackrel{}{\alpha }_1+\stackrel{}{\alpha }_2}}{𝐙}}`$ (3.8)
where the relative moduli space $`\stackrel{~}{}_{\stackrel{}{\alpha }_1+\stackrel{}{\alpha }_2}`$ is the Euclidean Taub-NUT space, endowed with the metric
$`ds^2=\left(1+{\displaystyle \frac{1}{r}}\right)(dr^2+r^2d\theta ^2+^2\mathrm{sin}^2\theta d\varphi ^2)+\left(1+{\displaystyle \frac{1}{r}}\right)^1(d\psi +\mathrm{cos}\theta d\varphi )^2`$ (3.9)
Here $`r`$,$`\theta `$ and $`\varphi `$ are the familiar spherical polar coordinates. The coordinate $`\psi `$ arises from $`U(1)`$ gauge transformations. The manifold has a $`SU(2)_R\times U(1)`$ isometry, of which only a $`U(1)_R\times U(1)`$ are manifest in the above coordinates. The holomorphic $`U(1)_R`$ isometry acts by rotating the two monopoles: $`\varphi \varphi +c`$. The tri-holomorpic $`U(1)`$ isometry changes the relative phase of the monopoles: $`\psi \psi +c`$. Both of these actions have a unique fixed point at $`r=0`$, the “nut” of Taub-NUT. However, the combined action with Killing vector $`_\psi +_\varphi `$ has a fixed point along the half-line $`\theta =\pi `$, with $`\psi `$ fibered over this line to produce the cigar shown in Figure 1. This is the relative moduli space $`\stackrel{~}{𝒲}_{\stackrel{}{\alpha }_1+\stackrel{}{\alpha }_2}`$.
Similar calculations hold for monopoles of charge $`\stackrel{}{g}=_{i=1}^{N1}\stackrel{}{\alpha }_i`$, whose dynamics is described by a class of toric hyperKähler metrics, known as the Lee-Weinberg-Yi metrics . Once again, a suitable $`𝐒^1`$ action on these spaces can be identified such that the fixed points localize on $`𝒲_\stackrel{}{g}`$, the moduli space of domain walls in $`U(1)`$ gauge theories with $`N`$ charged scalars.
#### Example 3: $`\stackrel{}{g}=\stackrel{}{\alpha }_1+2\stackrel{}{\alpha }_2+\stackrel{}{\alpha }_1`$
As described in the previous section, the simplest domain wall charge $`\stackrel{}{g}=_in_i\stackrel{}{\alpha }_i`$ with some $`n_i>1`$ occurs for $`\stackrel{}{g}=\stackrel{}{\alpha }_1+2\stackrel{}{\alpha }_2+\stackrel{}{\alpha }_3`$, and corresponds to a monopole in a $`SU(4)`$ gauge theory. No explicit expression for the metric on this monopole moduli space is known although, given the results of , such a computation may be feasible. Without an explicit expression for the metric in this, and more complicated examples, we need a more powerful method to describe the moduli space. This is provided by the Nahm construction, which we now review
### 3.2 D-Branes and Nahm’s Equations
The moduli space of magnetic monopoles is isomorphic to the moduli space of Nahm data. Here we review the Nahm construction of the monopole moduli space and, in particular, the embedding within the framework of D-branes due to Diaconescu . This will be useful to compare to the domain walls of the next section.
In the D-brane construction, Nahm’s equations arise as the low-energy description of D-strings suspended between D3-branes . The $`SU(N)`$ Yang-Mills theory lives on the worldvolume of $`N`$ D3-branes separated in, say, the $`x_6`$ direction, with the $`i^{\mathrm{th}}`$ D3-brane placed at position $`(x_6)_i=m_i`$ in accord with the adjoint expectation value (3.2). The monopole of charge $`\stackrel{}{g}=_in_i\stackrel{}{\alpha }_i`$ corresponds to suspending $`n_i`$ D-strings between the $`i^{\mathrm{th}}`$ and $`(i+1)^{\mathrm{th}}`$ D3-brane. This configuration is shown in Figure 3.
The motion of the D-strings in each segment $`m_ix_6m_{i+1}`$ is governed by four hermitian $`n_i\times n_i`$ matrices, $`X_1,X_2,X_3`$ and $`A_6`$ subject to the covariant version of Nahm’s equations,
$`{\displaystyle \frac{dX_\mu }{dx_6}}i[A_6,X_\mu ]{\displaystyle \frac{i}{2}}ϵ_{\mu \nu \rho }[X_\nu ,X_\rho ]=0m_ix_6m_{i+1}`$ (3.10)
modulo $`U(n_i)`$ gauge transformations acting on the interval $`m_ix_6m_{i+1}`$, and vanishing at the boundaries. The $`X_\mu `$ form the triplet representation of the $`SU(2)_R`$ symmetry which rotates monopoles in $`𝐑^3`$. The $`U(1)^{N1}`$ surviving gauge transformations acts on the Nahm data by constant shifts of the $`(N1)`$ “Wilson lines” $`A_6A_6+c1_{n_i}`$.
The interactions between neighbouring segments depends on the relative size of the matrices and is given by
$`n_i=n_{i+1}`$ In this case the $`U(n_i)`$ gauge symmetry is extended to the interval $`m_ix_6m_{i+2}`$ and an impurity is added to the right-hand-side of Nahm’s equations, which now read
$`{\displaystyle \frac{dX_\mu }{dx_6}}i[A_6,X_\mu ]{\displaystyle \frac{i}{2}}ϵ_{\mu \nu \rho }[X_\nu ,X_\rho ]=\omega _\alpha \sigma _\mu ^{\alpha \beta }\omega _\beta ^{}\delta (x_6m_{i+1})`$ (3.11)
Here $`\sigma _\mu `$ are the Pauli matrices. The impurity degrees of freedom lie in the complex 2-vector, $`\omega _\alpha =(\psi ,\stackrel{~}{\psi }^{})`$ which is a doublet under the $`SU(2)_R`$ symmetry. Both $`\psi `$ and $`\stackrel{~}{\psi }^{}`$ are themselves complex $`n_i`$ vectors, transforming in the fundamental representation of the $`U(n_i)`$ gauge group. The combination $`\omega _\alpha \sigma _\mu ^{\alpha \beta }\omega _\beta ^{}`$ is thus an $`n_i\times n_i`$ matrix, transforming in the adjoint representation of the gauge group. The $`\omega _\alpha `$ fields can be thought of as a hypermultiplet arising from $`D1D3`$ strings
$`n_i=n_{i+1}1`$ In this case $`X_\mu (X_\mu )_{}`$, a set of three $`n_i\times n_i`$ matrices, as $`x_6(m_i)_{}`$ from the left. To the right of $`m_i`$, the $`X_\mu `$ are $`(n_i+1)\times (n_i+1)`$ matrices which must obey
$`X_\mu \left(\begin{array}{cc}y_\mu & a_\mu ^{}\\ a_\mu & (X_\mu )_{}\end{array}\right)\mathrm{as}x_6(m_i)_+`$ (3.14)
where $`y_\mu 𝐑`$ and each $`a_\mu `$ is a complex $`n_i`$-vector.
$`n_in_{i+1}2`$ Once again we take $`X_\mu (X_\mu )_{}`$ as $`x_6(m_i)_{}`$ but, from the other side, the matrices $`X_\mu `$ now have a simple pole at the boundary,
$`X_\mu \left(\begin{array}{cc}J_\mu /(x_6m_i)+Y_\mu & 0\\ 0& (X_\mu )_{}\end{array}\right)\mathrm{as}x_6(m_i)_+`$ (3.17)
where $`J_\mu `$ is the irreducible $`(n_{i+1}n_i)\times (n_{i+1}n_i)`$ representation of $`su(2)`$, and $`Y_\mu `$ are now constant $`(n_{i+1}n_i)\times (n_{i+1}n_i)`$ matrices.
Case 2 above is usually described as a subset of Case 3 (with the one-dimensional irreducible $`su(2)`$ representation given by $`J_\mu =0`$). Here we have listed Case 2 separately since when we come to describe a similar construction for domain walls, only Case 1 and 2 above will appear. The conditions for $`n_i<n_{i+1}`$ were derived in by starting with the impurity data (3.11) and taking several monopoles to infinity. Obviously, for $`n_i>n_{i+1}`$, one imposes the same boundary conditions described above, only flipped in the $`x_6`$ direction.
The space of solutions to Nahm’s equations, subject to the boundary conditions detailed above, is isomorphic to the monopole moduli space $`_\stackrel{}{g}`$. Moreover, there exists a natural hyperKähler metric on the solutions to Nahm’s equations which can be shown to coincide with the Manton metric on the monopole moduli space. For the $`\stackrel{}{g}=\stackrel{}{\alpha }_1+\stackrel{}{\alpha }_2`$ monopole, the metric on the associated Nahm data was computed in and shown to give rise to the Euclidean Taub-NUT metric (3.9). For the $`\stackrel{}{g}=_i\stackrel{}{\alpha }_i`$ monopoles, the corresponding computation was performed in , resulting in the Lee-Weinberg-Yi metrics .
## 4 D-Branes and Domain Walls
In this section we would like to realize the domain walls that we described in Section 2 on the worldvolume of D-branes, mimicking Diaconescu’s construction for monopoles. From the resulting D-brane set-up we shall read off the world-volume dynamics of the domain walls to find that they are described by a truncated version of Nahm’s equations (3.10). Nahm’s equations have also arisen as a description of domain walls in $`𝒩=1^{}`$ theories , although the relationship, if any, with the current work is unclear. Domain walls of the type described in Section 2 were previously embedded in D-branes in and several properties of the solitons were extracted (see in particular the latter reference). However, the worldvolume dynamics of the walls is difficult to determine in these set-ups and the relationship to magnetic monopoles obscured.
We start by constructing the theory with eight supercharges on the worldvolume of D-branes . For definiteness we choose to build the $`𝒩=2`$, $`d=3+1`$ theory in IIA string theory although, by T-duality, we could equivalently work with any spacetime dimension<sup>1</sup><sup>1</sup>1In fact, as explained in , the overall $`U(1)U(k)`$ is decoupled in the IIA brane set-up after lifting to M-theory. This effect will not concern us here.. The construction is well known and is drawn in Figure 4. We suspend $`k`$ D4-branes between two NS5-branes, and insert a further $`N`$ D6-branes to play the role of the fundamental hypermultiplets. The worldvolume dimensions of the branes are
$`NS5:`$ $`012345`$
$`D4:`$ $`01236`$
$`D6:`$ $`0123789`$
The gauge coupling $`e^2`$ and the Higgs vev $`v^2`$ are encoded in the separation of the NS5-branes in the $`x_6`$ and $`x_9`$ directions respectively, while the masses $`m_i`$ are determined by the positions of the D6-branes in the $`x_4`$ direction (we choose the D6-branes to be coincident in the $`x_5`$ direction, corresponding to choosing all masses to be real).
$`{\displaystyle \frac{1}{e^2}}{\displaystyle \frac{\mathrm{\Delta }x_6}{l_sg_s}}|_{NS5},v^2{\displaystyle \frac{\mathrm{\Delta }x_9}{l_s^3g_s}}|_{NS5},m_i{\displaystyle \frac{x_4}{l_s^2}}|_{D6_i}`$ (4.1)
After turning on the Higgs vev $`v^2`$, the D4-branes must split on the D6-branes in order to preserve supersymmetry. The S-rule ensures that each D6-brane may play host to only a single D4-brane. In this manner a vacuum of the theory is chosen by picking $`k`$ out of the $`N`$ D6-branes on which the D4-branes end, in agreement with equation (2.1).
The domain walls correspond to a configuration of D4-branes which start life at $`x^3=\mathrm{}`$ in a vacuum configuration $`\mathrm{\Xi }_{}`$, and end up at $`x^3=+\mathrm{}`$ in a distinct vacuum $`\mathrm{\Xi }_+`$. As is clear from Figure 4, as D4-branes walls interpolate in $`x^1`$, they must also move in both the $`x^4`$ direction and the $`x^9`$ direction . The NS-branes and D6-branes are linked, meaning that a D4-brane is either created or destroyed as they pass the NS5-branes in the $`x^6`$ direction . In the domain wall background, which of these possibilities occurs differs if we move the D6-branes to the left or right since the D4-brane charge varies from one end of the domain wall to the other.
As it stands, it is difficult to read off the dynamics of the D4-branes in Figure 4. However, we can make progress by taking the $`e^2\mathrm{}`$ limit, in which the two NS5-branes become coincident in the $`x^6`$ direction. After rotating our viewpoint, the system of branes now looks like the ladder configuration shown in Figure 5 (note that we have also rotated the branes relative to Figure 4, so the horizontal is the $`x^4`$ direction). We are left with a series of D4-branes, now with worldvolume $`02349`$, stretched between $`N`$ D6-branes, while simultaneously sandwiched between two NS5-branes. Following these manoeuvres, one finds that the domain wall $`\stackrel{}{g}=_in_i\stackrel{}{\alpha }_i`$ results in $`n_i`$ D4-branes stretched between the $`i^{\mathrm{th}}`$ and $`(i+1)^{\mathrm{th}}`$ D6-branes (counting from the top, since we have chosen the ordering $`m_i>m_{i+1}`$.
It may be worth describing how the domain wall charges arise directly in the set-up of Figure 5. We start in a chosen vacuum $`\mathrm{\Xi }_{}`$, denoted by placing $`k`$ pairs of white dots on $`N`$ distinct D6-branes, as shown in the figure. A domain wall arises every time a pair of dots is proceeds to another D6-brane, dragging a D4-brane behind it like clingwrap. The S-rule translates to the fact that two pairs of dots may not simultaneously lie on the same D6-brane. The final vacuum $`\mathrm{\Xi }_+`$ is denoted by the black dots in the figure and the domain wall charges $`n_i`$ are given by the number of times a D4-brane has been pulled between the $`i^{\mathrm{th}}`$ and $`(i+1)^{\mathrm{th}}`$ D6-branes.
### 4.1 Domain Wall Dynamics
We are now in a position to read off the dynamics of the domain walls. In the absence of the NS5-branes, the D4-branes would stretch to infinity in the $`x_9`$ direction, and the resulting D-brane set-up in Figure 5 is T-dual to the monopoles in Figure 3. The presence of the NS5-branes projects out half the degrees of freedom of the monopoles, leaving a simple linear set of equations. In each segment $`m_ix_4m_{i+1}`$ the domain walls are described by two $`n_i\times n_i`$ matrices $`X_3`$ and $`A_4`$ satisfying
$`{\displaystyle \frac{dX_3}{dx_4}}i[A_4,X_3]=0`$ (4.2)
modulo $`U(n_i)`$ gauge transformations acting on the interval $`m_ix_4m_{i+1}`$, and vanishing at the boundaries. As in the case of monopoles, the interactions between neighbouring segments depends on the relative size of the matrices:
$`n_i=n_{i+1}`$ Again, the $`U(n_i)`$ gauge symmetry is extended to the interval $`m_ix_4m_{i+2}`$ and an impurity is added to the right-hand-side of Nahm’s equations, which now read
$`{\displaystyle \frac{dX_3}{dx_4}}i[A_4,X_3]=\pm \psi \psi ^{}\delta (x_4m_{i+1})`$ (4.3)
where the impurity degree of freedom $`\psi `$ transforms in the fundamental representation of the $`U(n_i)`$ gauge group, ensuring the combination $`\psi \psi ^{}`$ is a $`n_i\times n_i`$ matrix transforming, like $`X_3`$, in the adjoint representation. These $`\psi `$ degrees of freedom are chiral multiplets which survive the NS5-brane projection. We shall see shortly that the choice of $`\pm `$ sign will dictate the relative ordering of the domain walls along the $`x_3`$ direction.
$`n_i=n_{i+1}1`$ In this case $`X_3(X_3)_{}`$, an $`n_i\times n_i`$ matrix, as $`x_4(m_i)_{}`$ from the left. To the right of $`m_i`$, $`X_3`$ is a $`(n_i+1)\times (n_i+1)`$ matrix obeying
$`X_3\left(\begin{array}{cc}y& a^{}\\ a& (X)_{}\end{array}\right)\mathrm{as}x_4(m_i)_+`$ (4.6)
where $`y_\mu 𝐑`$ and each $`a_\mu `$ is a complex $`n_i`$-vector. The obvious analog of this boundary condition holds when $`n_i=n_{i+1}+1`$.
These boundary conditions obviously descend from the original Nahm boundary conditions for monopoles. Just as the space of Nahm data is isomorphic to the moduli space of magnetic monopoles, we conjecture that the moduli space of linearized Nahm data described above is isomorphic to the moduli space of domain walls. We shall shortly show that it indeed captures the most relevant aspect of domain walls: their ordering. In fact, the linearized Nahm equations (4.2) are rather trivial to solve. We first employ the $`_iU(n_i)`$ gauge transformations to make $`A_4(x_4)`$ a constant in each interval $`m_ix_4m_{i+1}`$. This can be achieved by first diagonalizing $`A_4`$, and subsequently acting with the $`U(1)^{n_i}`$ transformation $`A_4A_4_4\alpha `$ where, in each segment $`m_ix_4m_{i+1}`$, $`\alpha `$ is given by
$`\alpha (x_4)={\displaystyle _{m_i}^{x_4}}A_4(x_4^{})𝑑x_4^{}\left[{\displaystyle _{m_i}^{m_{i+1}}}A_4(x_4^{})𝑑x_4^{}\right]{\displaystyle \frac{m_ix_4}{m_im_{i+1}}}`$ (4.7)
which has the property that $`\alpha (m_i)=\alpha (m_{i+1})=0`$. Further gauge transformations with non-zero winding on the interval ensure that $`A_4`$ is periodic, with each eigenvalue lying in $`A_4[0,2\pi /(m_im_{i+1}))`$. These $`N1`$ “Wilson lines” will play the role of the phases associated to domain wall system. Note that when $`n_i=n_{i+1}`$, the above choice of gauge leaves a residual $`U(n_i)`$ gauge symmetry acting only on the chiral impurity $`\psi `$. In this gauge we can now easily integrate (4.2) in each interval,
$`X_3(x_4)=e^{iA_4x_4}\widehat{X}_3e^{iA_4x_4}`$ (4.8)
where the eigenvalues of $`X_3`$ are independent of $`x_4`$ in each interval. We identify these $`n_i`$ eigenvalues with the positions of the $`n_i`$ $`\stackrel{}{\alpha }_i`$ elementary domain walls.
We are now in a position to derive the linearized Nahm equations (4.2) from the original Nahm equations (3.10) in terms of a fixed point set of a $`\widehat{U}(1)`$ action. Consider first the action of the $`U(1)_RSU(2)_R`$ isometry on the Nahm data, which rotates $`X_1`$ and $`X_2`$ while leaving $`X_3`$ fixed. This rotation also acts on the impurity $`\omega =(\psi ,\stackrel{~}{\psi }^{})`$ by $`(\psi ,\stackrel{~}{\psi })e^{i\alpha }(\psi ,\stackrel{~}{\psi })`$. To retain half of the impurities for the domain wall equations (4.3), we need to compensate for this transformation with the residual $`U(1)U(n_i)`$ transformation acting on the appropriate impurity $`\omega `$ by $`\omega e^{i\beta }\omega `$. By choosing $`\beta =\pm \alpha `$ we can pick a $`\widehat{U}(1)`$ action which leaves either the $`\psi `$ or the $`\stackrel{~}{\psi }`$ impurity invariant. Which we choose to save is correlated with the choice of minus sign in (4.3) which, in turn, dictates the ordering of neighbouring domain walls as we shall now demonstrate.
To summarize, we have shown the that description of domain wall dynamics (4.2) arises from the fixed point of a $`\widehat{U}(1)`$ on the original Nahm equations (3.10). This action descends to a $`\widehat{U}(1)`$ isometry on the monopole moduli space $`_\stackrel{}{g}`$, the fixed points of which coincide with the domain wall moduli space $`𝒲_\stackrel{}{g}`$. A physical explanation for this correspondence follows along the lines of : in theories in the Higgs phase, confined magnetic monopoles with charge $`\stackrel{}{g}`$ exist, emitting $`k`$ multiple vortex strings. When these vortex strings coincide, the worldvolume theory is of the form described in Section 2 and the monopoles appear as charge $`\stackrel{}{g}`$ kinks.
### 4.2 The Ordering of Domain Walls Revisited
As explained in Section 2, in contrast to monopoles, domain walls must satisfy a specific ordering on the $`x^3`$ line. We will now show that this ordering is encoded in the boundary conditions described above. Suppose first that $`n_i=n_{i+1}`$. The positions of the $`\stackrel{}{\alpha }_i`$ domain walls are given by the eigenvalues of $`X_3`$ restricted to the interval $`m_ix_4m_{i+1}`$. Let us denote this matrix as $`X_3^{(i)}`$ and the eigenvalues as $`\lambda _m^{(i)}`$, where $`m=1,\mathrm{}n_i`$. The impurity (4.3) relates the two sets of eigenvalues by the jumping condition
$`X_3^{(i+1)}=X_3^{(i)}+\psi \psi ^{}`$ (4.9)
where we have chosen the positive sign for definiteness. However, from the discussion in Section 2 (see, in particular, figure 3) we know that the domain walls cannot have arbitrary position but must be interlaced,
$`\lambda _1^{(i)}\lambda _1^{(i+1)}\lambda _2^{(i)}\mathrm{}\lambda _{n_i1}^{(i+1)}\lambda _{n_i}^{(i)}\lambda _{n_i}^{(i+1)}`$ (4.10)
We will now show that the ordering of domain walls (4.10) follows from the impurity jumping condition (4.9).
To see this, consider firstly the situation in which $`\psi ^{}\psi \mathrm{\Delta }\lambda _m^{(i)}`$ so that the matrix $`\psi \psi ^{}`$ may be treated as a small perturbation of $`X_3^{(i)}`$. The positivity of $`\psi \psi ^{}`$ ensures that each $`\lambda _m^{(i+1)}\lambda _m^{(i)}`$. Moreover, it is simple to show that the $`\lambda _m^{(i+1)}`$ increase monotonically with $`\psi ^{}\psi `$. This leaves us to consider the other extreme, in which $`\psi ^{}\psi \mathrm{}`$. It this limit $`\psi `$ becomes one of the eigenvectors of $`X_3^{(i+1)}`$ with eigenvalue $`\lambda _{n_i}^{(i+1)}=\psi ^{}\psi \mathrm{}`$ which reflects the fact that this limit corresponds to the situation in which the last domain wall is taken to infinity. What we want to show is that the remaining $`(n_i1)`$ $`\stackrel{}{\alpha }_{i+1}`$ domain walls are trapped between the $`n_i`$ $`\stackrel{}{\alpha }_i`$ domain walls as depicted in Figure 3. Define the $`n_i\times n_i`$ projection operator
$`P=1\widehat{\psi }\widehat{\psi }^{}`$ (4.11)
where $`\widehat{\psi }=\psi /\sqrt{\psi ^{}\psi }`$. The positions of the remaining $`(n_i1)`$ $`\stackrel{}{\alpha }_{i+1}`$ domain walls are given by the (non-zero) eigenvalues of $`PX_3^{(i)}P`$. We must show that, given a rank $`n`$ hermitian matrix $`X`$, the eigenvalues of $`PXP`$ are trapped between the eigenvalues of $`X`$. This elementary property of hermitian matrices can be seen as follows:
$`det(PXP\mu )`$ $`=`$ $`det(XP\mu )`$
$`=`$ $`det(X\mu X\widehat{\psi }\widehat{\psi }^{})`$
$`=`$ $`det(X\mu )det(1(X\mu )^1X\widehat{\psi }\widehat{\psi }^{})`$
Since $`\widehat{\psi }\widehat{\psi }^{}`$ is rank one, we can write this as
$`det(PXP\mu )`$ $`=`$ $`det(X\mu )[1\mathrm{Tr}((X\mu )^1X\widehat{\psi }\widehat{\psi }^{})]`$ (4.12)
$`=`$ $`\mu det(X\mu )\mathrm{Tr}((X\mu )^1\widehat{\psi }\widehat{\psi }^{})`$
$`=`$ $`\mu \left[{\displaystyle \underset{m=1}{\overset{n}{}}}(\lambda _m\mu )\right]\left[{\displaystyle \underset{m=1}{\overset{n}{}}}{\displaystyle \frac{|\widehat{\psi }_m|^2}{\lambda _m\mu }}\right]`$
where $`\widehat{\psi }_m`$ is the $`m^{\mathrm{th}}`$ component of the vector $`\psi `$. We learn that $`PXP`$ has one zero eigenvalue while, if the eigenvalues $`\lambda _m`$ of $`X`$ are distinct, then the eigenvalues of $`PXP`$ lie at the roots the function
$`R(\mu )={\displaystyle \underset{m=1}{\overset{n}{}}}{\displaystyle \frac{|\widehat{\psi }_m|^2}{\lambda _m\mu }}`$ (4.13)
The roots of $`R(\mu )`$ indeed lie between the eigenvalues $`\lambda _m`$. This completes the proof that the impurities (4.3) capture the correct ordering of the domain walls.
The same argument shows that the boundary condition (4.6) gives rise to the correct ordering of domain walls when $`n_{i+1}=n_i+1`$, with the $`\stackrel{}{\alpha }_i`$ domain walls interlaced between the $`\stackrel{}{\alpha }_{i+1}`$ domains walls. Indeed, it is not hard to show that (4.6) arises from (4.3) in the limit that one of the domain walls is taken to infinity.
## Acknowledgement
AH is supported in part by the CTP and LNS of MIT, DOE contract $`\mathrm{\#}`$DE-FC02-94ER40818, NSF grant PHY-00-96515, the BSF American-Israeli Bi-national Science Foundation and a DOE OJI Award. DT is supported by the Royal Society. |
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