id
stringlengths 27
33
| source
stringclasses 1
value | format
stringclasses 1
value | text
stringlengths 13
1.81M
|
|---|---|---|---|
warning/0001/astro-ph0001500.html | ar5iv | text | # The Host Galaxy of the Lensed Quasar Q 0957+5611footnote 11footnote 1Based on Observations made with the NASA/ESA Hubble Space Telescope, obtained at the Space Telescope Science Institute, which is operated by AURA, Inc., under NASA contract NAS 5-26555.
## 1 Introduction
Gravitational lenses offer an attractive independent method of determining the Hubble constant ($`H_0`$) on cosmological scales without the systematic difficulties associated with the local distance ladder (Refsdal 1964, 1966). Six of the more than 50 known gravitational lenses now have time delays that can be used to estimate $`H_0`$: B 0218+357 (Biggs et al. 1999); Q 0957+561 (Schild & Thomson 1995; Kundić et al. 1997; Haarsma et al. 1999); PG 1115+080 (Schechter et al. 1997; Barkana 1997); B 1600+434 (Hjorth et al. 1999); B 1608+656 (Fassnacht et al. 1999); and PKS 1830$``$211 (Lovell et al. 1998); preliminary $`H_0`$ estimates for five of these systems have been compiled by Koopmans & Fassnacht (1999). Once a time delay is accurately determined, the uncertainties in the derived value of the Hubble constant are due almost entirely to the systematic uncertainties in the model for the lensing potential, plus “cosmic variance” due to the effects of weak density inhomogeneities along the line of sight (Seljak 1994; Barkana 1996).
Only two of the time delay lenses have been modeled in sufficient detail to fully understand the systematic uncertainties created by their geometries: Q 0957+561 (Falco, Gorenstein & Shapiro 1985, 1991; Kochanek 1991; Bernstein, Tyson & Kochanek 1993; Grogin & Narayan 1996; Chartas et al. 1998; Barkana et al. 1999; Bernstein & Fischer 1999; Chae 1999; Romanowsky & Kochanek 1999), and PG 1115+080 (Schechter et al. 1997; Courbin et al. 1997; Keeton & Kochanek 1997; Saha & Williams 1997; Impey et al. 1998). In Q 0957+561, the lens consists of a brightest cluster galaxy and its parent cluster, and the value of the Hubble constant depends on the mass balance between the two components. Judgments about the ability of models to determine the correct mass balance vary significantly: optimistic estimates yield $`H_0=61_{15}^{+13}`$ km s<sup>-1</sup> Mpc<sup>-1</sup>at 95% confidence (Grogin & Narayan 1996, using the stellar dynamical models of Romanowsky & Kochanek 1999); while more pessimistic estimates yield $`H_0=77_{24}^{+29}`$ (Kochanek 1991; Bernstein et al. 1993; Bernstein & Fischer 1999). In PG 1115+080, the lens is one of the brighter galaxies in a small group, but the four-image geometry easily determines the relative roles of the primary lens galaxy and the group. Instead, the value of the Hubble constant depends on the assumed radial mass distribution of the primary lens galaxy. Galaxy models with dark matter and mass distributions consistent with the best estimates for early-type galaxies lead to low values for the Hubble constant ($`H_0=44\pm 4`$ km s<sup>-1</sup> Mpc<sup>-1</sup>; Impey et al. 1998).
Impey et al. (1998) discovered an Einstein ring image of the quasar host galaxy in the PG 1115+080 lens, and noted that the geometry of the ring could be used to break the degeneracy in the lens models. Here we report the discovery of the host galaxy in the Q 0957+561 lens and discuss the implications for lens models and the Hubble constant. In §2 we describe our observations. In §3 we discuss the data and methods used to model the system, summarizing previous work and introducing our new data and techniques. In §4 we examine the arcs predicted by previously published lens models and demonstrate that they fail to match the observed arcs. In §5 we present new models that match all data, including the arcs, and discuss their implications for the properties of the lens galaxy and cluster and for $`H_0`$. In §6 we summarize our results and discuss prospects for further improving the constraints on this lens system and on $`H_0`$.
## 2 Observations
We observed Q 0957+561 with the Hubble Space Telescope as part of the CfA–Arizona Space Telescope Lens Survey (CASTLES; Lehár et al. 1999; Falco et al. 1999). Using the NIC2 camera, we obtained a 2800-second F160W (H band) image of Q 0957+561 divided into four dithered exposures. A log of the observations is presented in Table 1. We reduced the images using nicred, a custom reduction package developed for CASTLES (Lehár et al. 1999). We also reanalyzed archival WFPC2 images, including a 32200-second F555W (V band) exposure and a 2620-second F814W (I band) exposure (Bernstein et al. 1997). In the optical images, Q 0957+561 fell in chip WF3 of the WFPC2 camera (with pixel size $``$0$`\stackrel{}{\mathrm{.}}`$1, compared to $``$0$`\stackrel{}{\mathrm{.}}`$076 for NIC2). Because the images A and B of Q 0957+561 were saturated in the WFPC2 images, we derived our astrometry of A, B and the lens galaxy G1 exclusively from our unsaturated NIC2 data.
Figure 1a shows our combined H band image, which prominently shows the lens galaxy and the two quasar images. The quasar host galaxy is visible as a faint arc next to the quasar A image. Table 2 summarizes a photometric model consisting of point sources for the two quasar images and elliptical de Vaucouleurs models for the main lens galaxy (G1) and a neighboring galaxy (G2, not seen in Figure 1).<sup>2</sup><sup>2</sup>2Young et al. (1981) tabulated objects in the field of Q 0957+561 and applied the labels G1–G5 to five of the bright galaxies. We labeled only objects close to the lens galaxy, so other than G1 our labels do not match those of Young et al. Our object G2 appears as object #97 in their Table 1. Our objects G3 and G4 do not appear in their list. The lens galaxy is known to have a small ellipticity gradient and isophote twist (Bernstein et al. 1997), but we used only a simple elliptical model. Figure 1b shows the residuals after subtracting the PSF-convolved photometric model from the original image, and Figure 1c shows these residuals convolved with a $``$0$`\stackrel{}{\mathrm{.}}`$076 (1 pixel) FWHM Gaussian to enhance the visibility of low surface brightness features. The residual image shows the host galaxy image near quasar A to be an arc distorted tangentially relative to the lens galaxy. It also reveals an extended asymmetric arc near the B quasar image. If the residuals near quasar B were created by the error of fitting a simple elliptical surface brightness model to a galaxy with a radially varying ellipticity and orientation, we would expect them to have reflection symmetry through the center of the galaxy. Because the arc near quasar B lacks such symmetry, we conclude that it is the lensed counterpart of the arc near quasar A, and that both are images of the quasar host galaxy.
The residual H band image also shows two additional sources, one east of quasar A and another just west of quasar B (labeled G3 and G4 in Figure 1c). Their positions and photometric properties are given in Table 2. These objects correspond in position and V magnitude to two of the faint sources seen in the deep V band image and labeled Blobs 7 and 1 (respectively) by Bernstein et al. (1997). Their colors (especially V$``$I) are similar to those of G1 and G2 and suggest that they are probably faint cluster members. Whether or not they are associated with the cluster, it is clear that they are not multiply imaged and are not associated with the quasar source. The fact that G4 is not part of the host galaxy arc will be important to remember when examining the arc structure predicted by lens models (see §3.3).
The lensed host galaxy is not detected at significant levels in the shallow I band image, although a visual examination suggests that low-level residuals are present and should be detected at a significant level with a longer integration. The deep V band image does not have large extended arcs, but it does show seven faint “blobs” and a thin arclet (Bernstein et al. 1997). As noted above, at least two of the blobs are probably faint cluster galaxies. However, Blobs 2 and 3 and the arclet have been identified as possible lensed features, and lens models support this interpretation (Barkana et al. 1999; Bernstein & Fischer 1999; Chae 1999). Since the V band corresponds to rest-frame UV at the quasar redshift $`z_s=1.41`$, these features probably correspond to discrete star forming regions in the host galaxy. We follow previous models and include Blobs 2 and 3 plus two “Knots” in the arclet as model constraints (see §3.2).
We estimated the fluxes of the host galaxy images using the IRAF task polyphot, which measures fluxes within polygonal apertures. We traced polygonal apertures along the edges of the distorted images of the host galaxy at the $`3\sigma `$ level above the sky. We found that both lensed images of the host galaxy have H band brightnesses of $``$18.4 mag, and a surface brightness of $``$20.6 mag/arcsec<sup>2</sup> (both images have approximately the same area). The main source of uncertainty for the brightness estimates is the relatively large size of the residuals from the subtraction of the quasar images, compared to the (low) brightness levels of the arcs; we estimate the uncertainties to be $``$0.3 mag.
Given an acceptable lens model we can map the host galaxy images to the source plane to obtain a map of the unlensed host galaxy (see §3.3). In principle we could use the source maps to measure the photometric properties of the host galaxy. However, in practice the imperfect quasar subtraction corrupts the flux in the bright central regions of the source and hinders the measurement. One robust statement we can make is that part of the host galaxy is doubly-imaged like the quasar, but part of it is quadruply-imaged (see §4). This accounts for the shape differences between the A and B arcs: the A arc is a single distorted image of the host galaxy, while the B arc is a composite of three images that straddle the lensing critical line.
## 3 Constraining models of Q 0957+561
Q 0957+561 is the most thoroughly studied gravitational lens. The system comprises a radio-loud quasar at redshift $`z_s=1.41`$ lensed into two images by a brightest cluster galaxy and its parent cluster at redshift $`z_l=0.36`$ (Walsh, Carswell & Weymann 1979; Young et al. 1980). There is a time delay of $`417\pm 3`$ days between the images (Schild & Thomson 1995; Kundić et al. 1997; Haarsma et al. 1999). VLBI observations have resolved each image into a core and $``$80 milli-arcsecond jet (Garrett et al. 1994). Deep optical images have uncovered faint “Blobs” and “Knots” that are probably lensed image pairs of star forming regions in the quasar host galaxy (Bernstein et al. 1997). Our observations have revealed infrared arcs representing distorted images of the quasar host galaxy. In this section we discuss how to use these and other data to constrain models of the system and values for the Hubble constant $`H_0`$. In §3.1 we review the theory of lens modeling and discuss important model degeneracies. In §§3.2 and 3.3 we discuss previous and new observational constraints on models. In §§3.4 and 3.5 we describe classes of models applied to the lens galaxy and cluster in Q 0957+561, summarizing previous classes and introducing a new one. Finally, in §3.6 we discuss how we apply the observational constraints to our new class of models.
### 3.1 Basic lens theory: model degeneracies
In a multiply-imaged gravitational lens system, the light from a distant source is deflected by the gravitational potential of foreground objects so that we observe multiple images of the source. The lensing potential is usually dominated by a single galaxy, although there may be non-negligible perturbations from other objects nearby.<sup>3</sup><sup>3</sup>3There may also be a contribution to the potential from density fluctuations along the line of sight, but it is usually small compared to the contribution from objects at the same redshift as the main lens galaxy (Seljak 1994; Barkana 1996; Keeton, Kochanek & Seljak 1997). In a few cases there are two lens galaxies enclosed by the lensed images (e.g. Jackson, Nair & Browne 1997; Koopmans et al. 1999), but this complication is not an issue in Q 0957+561. Basic lens theory is presented in the book by Schneider, Ehlers & Falco (1992), and we quote the relevant results here. The lensing potential $`\varphi `$ is determined by the two-dimensional Poisson equation $`^2\varphi =2\kappa `$, where $`\kappa =\mathrm{\Sigma }/\mathrm{\Sigma }_{crit}`$ is the surface mass density in units of the critical surface density for lensing (in angular units),
$$\mathrm{\Sigma }_{crit}=\frac{c^2}{4\pi G}\frac{D_{ol}D_{os}}{D_{ls}},$$
(1)
where $`D_{ol}`$ and $`D_{os}`$ are angular diameter distances from the observer to the lens and source, respectively, and $`D_{ls}`$ is the angular diameter distance from the lens to the source. The lensing potential deflects a light ray so that the angular position $`\stackrel{}{u}`$ of the source on the sky and the angular position $`\stackrel{}{x}`$ of an image are related by the lens equation,
$$\stackrel{}{u}=\stackrel{}{x}\stackrel{}{}\varphi (\stackrel{}{x}).$$
(2)
There is an image corresponding to each solution $`\stackrel{}{x}_i`$ of this equation. Lensing introduces a time delay between the ray paths of two images of the same source. The time delay between images at positions $`\stackrel{}{x}_i`$ and $`\stackrel{}{x}_j`$ is
$$\mathrm{\Delta }t_{ij}=\frac{1+z_l}{c}\frac{D_{ol}D_{os}}{D_{ls}}\left[\frac{1}{2}\left(|\stackrel{}{x}_i\stackrel{}{u}|^2|\stackrel{}{x}_j\stackrel{}{u}|^2\right)\left(\varphi (\stackrel{}{x}_i)\varphi (\stackrel{}{x}_j)\right)\right],$$
(3)
where $`z_l`$ is the redshift of the lens. This equation is the basis of attempts to use lensing to determine the Hubble constant $`H_0`$. By measuring light curves of images one can determine the time delay $`\mathrm{\Delta }t_{ij}`$. A lens model gives the term in square brackets in eq. (3). The combination of distances is $`H_0^1`$ and only weakly dependent on other cosmological parameters.
A lens model consists of a description of the lensing potential $`\varphi `$. The observed images provide two principle constraints on such a model. First, roughly speaking an image at projected distance $`R`$ from the main lens galaxy measures the enclosed mass $`M(R)`$. In Q 0957+561 and many 2-image gravitational lenses, the two images lie at different distances $`R_1R_2`$, so they measure two masses $`M(R_1)`$ and $`M(R_2)`$ and hence constrain the mass profile (e.g. Grogin & Narayan 1996). Second, with at least three well-determined positions (two images and the lens galaxy, or four images), the images determine the quadrupole moment of the net potential. However, there are four properties of the mass distribution that must be determined for a complete description of the model: the mass profile of the main lens galaxy; the shape (ellipticity and orientation) of the main lens galaxy; the shear from the the gravitational tidal field induced by objects near the lens galaxy; and the amount of gravitational focusing (or “convergence”) contributed by the environment of the lens galaxy.
With more unknown quantities than constraints, there are two common degeneracies in the lens models. First, to lowest order the convergence $`\kappa `$ contributed by the environment cannot be determined by lens models, which leads to the so-called “mass sheet degeneracy” (Falco et al. 1985). If a lensing potential $`\varphi (\stackrel{}{x})`$ fits the observed data, then any potential
$$\varphi ^{}(\stackrel{}{x})=\frac{1}{2}\kappa |\stackrel{}{x}|^2+(1\kappa )\varphi (x)$$
(4)
will fit the data equally well. The $`\kappa `$ term is equivalent to the potential from a uniform mass sheet with surface density $`\mathrm{\Sigma }=\kappa \mathrm{\Sigma }_{crit}`$. The only potentially observable difference between the lens models represented by $`\varphi `$ and $`\varphi ^{}`$ is in the predicted time delays,
$$\mathrm{\Delta }t_{model}^{}=(1\kappa )\mathrm{\Delta }t_{model}.$$
(5)
However, this effect is not observable if we want to use gravitational lensing to determine the Hubble constant because it simply translates into a scaling of the inferred value for $`H_0`$. If $`H_0`$ and $`H_0^{}`$ are the values inferred from the two lens models, then
$$H_0^{}=(1\kappa )H_0.$$
(6)
This analysis has used only the lowest order term of the contribution from the environment, but the higher order terms cannot necessarily eliminate the degeneracy (e.g. Chae 1999). The mass sheet degeneracy is important for Q 0957+561 because the parent cluster of the lens galaxy contributes a significant convergence. The only way to break this degeneracy is to obtain an independent mass constraint to determine the relative contributions of the main lens galaxy and the environment to the mass enclosed by the images (see §3.2).
The second important degeneracy is between the shape of the main lens galaxy and the shear from the gravitational tidal field of objects in the environment of the lens galaxy. Because the images determine only the quadrupole moment of the total potential, there can be a wide range of parameter space in which the lens galaxy ellipticity and the external shear combine to produce the required joint quadrupole (e.g. Keeton et al. 1997). Consider a 2-image lens with image positions $`\stackrel{}{x}_A`$ and $`\stackrel{}{x}_B`$, a galaxy potential $`\varphi _{gal}`$, and an external shear with amplitude $`\gamma `$ and direction $`\theta _\gamma `$. Fitting the image positions exactly is equivalent to solving the equation
$$\stackrel{}{x}_A\stackrel{}{}\varphi _{gal}(\stackrel{}{x}_A)\mathrm{\Gamma }\stackrel{}{x}_A=\stackrel{}{x}_B\stackrel{}{}\varphi _{gal}(\stackrel{}{x}_B)\mathrm{\Gamma }\stackrel{}{x}_B$$
(7)
where the shear is described by the tensor
$$\mathrm{\Gamma }=\left[\begin{array}{cc}\hfill \gamma \mathrm{cos}2\theta _\gamma & \hfill \gamma \mathrm{sin}2\theta _\gamma \\ \hfill \gamma \mathrm{sin}2\theta _\gamma & \hfill \gamma \mathrm{cos}2\theta _\gamma \end{array}\right].$$
(8)
No matter what galaxy potential $`\varphi _{gal}`$ is used, it is straightforward to solve eq. (7) for the shear parameters $`\gamma `$ and $`\theta _\gamma `$. In other words, one can fit the image positions in a 2-image lens to arbitrary precision for any model lens galaxy, although in some cases the inferred shear $`\gamma `$ may be unphysically strong.
Constraints from the images fluxes or additional images often break this degeneracy, although the $`1\sigma `$ range of models may still be large. In previous models of Q 0957+561, the additional constraints came primarily from the VLBI observations of the quasar images, which show several discrete components with $``$0.1 milli-arcsecond errorbars (Garrett et al. 1994; Barkana et al. 1999). It is easy to be misled about the true structure of the potential, however, when using such strong constraints (see Kochanek 1991; Bernstein et al. 1993; Bernstein & Fischer 1999). Internal structure in the lens galaxy, such as isophote twists or lumpiness in the mass distribution, may affect the images at sub-milli-arcsecond scales (e.g. Mao & Schneider 1998). Neglecting that structure forces the models to adjust large-scale model components in order to fit small-scale constraints, leading them to converge to a best-fit solution that is well defined (i.e. with no apparent degeneracy) but incorrect. In Q 0957+561, the observed ellipticity gradient and isophote twist suggest that the galaxy does have important internal structure (Bernstein & Fischer 1997), and in §4 we show that models neglecting this structure indeed converged to incorrect solutions.
### 3.2 Previous observational constraints
In Q 0957+561 the high-resolution VLBI maps of the quasar images provide strong position constraints, plus somewhat weaker constraints on the relative magnification matrix between the images (Garrett et al. 1994; Barkana et al. 1999). The optical Blobs and Knots offer additional but weaker position constraints. The Knots are particularly useful because they appear to represent a pair of “fold” images and thus require the lensing critical line to pass between them (see Bernstein & Fischer 1999). The constraints from these images leave two of the three model degeneracies discussed in §3.1: the mass sheet degeneracy due to the convergence provided by the cluster; and the degeneracy between the lens galaxy shape and the cluster shear.
Several steps have been taken to break the mass sheet degeneracy by determining the relative masses of the lens galaxy and cluster. First, Rhee (1991), Falco et al. (1997), and Tonry & Franx (1999) measured the lens galaxy’s central velocity dispersion to be $`288\pm 9`$ km s<sup>-1</sup>. Romanowsky & Kochanek (1999) used stellar dynamical models to translate this into constraints on the lens galaxy mass and concluded that $`(80\pm 12)\%`$ of the image separation is contributed by the galaxy (and the rest by the cluster convergence). Second, Fischer et al. (1997) detected a weak lensing signal from the cluster, which Bernstein & Fischer (1999) translated into an estimate of the mean surface density of all mass inside an aperture of radius 30″ centered on the lens galaxy,
$$\kappa _{obs}_{30\mathrm{}}=\mathrm{\Sigma }_{30\mathrm{}}/\mathrm{\Sigma }_{crit}=0.26\pm 0.08.$$
(9)
Furthermore, Bernstein & Fischer (1999) showed that the surface mass density $`\kappa _{clus}=\mathrm{\Sigma }_{clus}/\mathrm{\Sigma }_{crit}`$ of the cluster at the position of the lens galaxy is given by
$$1\kappa _{clus}=\frac{1\kappa _{obs}_{30\mathrm{}}}{1\kappa _{mod}_{30\mathrm{}}},$$
(10)
where $`\kappa _{mod}`$ is the mean mass density of a model lens galaxy. In this way one can combine the weak lensing measurement with a lens model to estimate $`\kappa _{clus}`$ (see §5.3). There have been other attempts to measure the cluster mass, using the velocity dispersion of the galaxies in the cluster ($`715\pm 130`$ km s<sup>-1</sup>; Garrett, Walsh & Carswell 1992; Angonin-Willaime, Soucail & Vanderriest 1994) or the X-rays from the cluster gas (Chartas et al. 1998), but they are still limited by significant systematic uncertainties.
Although measurements of the cluster surface density are important (especially for determining $`H_0`$, see §5.4), they do not break the degeneracy between the lens galaxy shape and the cluster shear. Another approach is to constrain the shape of the lens galaxy’s mass distribution. The orientation of the lens galaxy’s light distribution should offer a guide to the orientation of its mass distribution, because other lenses suggest that mass and light are at least roughly aligned (typically within $``$10°; Keeton, Kochanek & Falco 1998). An important detail, though, is that the lens galaxy has a radially varying ellipticity and orientation (Bernstein et al. 1997): the ellipticity increases from $``$0.1 inside a radius of 1″ to $``$0.4 outside 10″, and the PA varies from $``$40° to $``$60° (albeit with large uncertainties). These variations suggest that the galaxy’s projected mass distribution may not have simple elliptical symmetry, and they may well need to be incorporated into lens models.
### 3.3 New constraints from the host galaxy arcs
The host galaxy arcs act as an extensive set of position constraints that may break the degeneracy between the lens galaxy shape and the cluster shear. Although the intrinsic source structure is unknown, the simplicity of the lensing geometry means that we can model the arcs self-consistently using a modified version of the Ring Cycle algorithm of Kochanek et al. (1989). Specifically, we know that there is a one-to-one mapping between the A image of the host galaxy and the source, so for a given lens model we can project the A arc back to the source plane to obtain the correct source structure for that lens model.<sup>4</sup><sup>4</sup>4A similar projection fails for the B arc because in all reasonable models this arc crosses the lensing critical line (see §4). Thus the map from the B arc to the source is not one-to-one. We can then reproject that source onto the image plane to predict the structure of the B arc. Since the arcs are very large compared to the size of the PSF, there is no need to include the effects of the PSF. For the arc map we use the smoothed residual image (Figure 1c).
This technique produces maps of the model arcs that can be compared visually to the observed arcs (see Figures 2–4 below). However, for modeling we want to quantify the differences, so we define a $`\chi ^2`$ term for the arcs using a pixel-by-pixel comparison of the observed and model surface brightnesses,
$$\chi _{arc}^2=\underset{i=1}{\overset{N_{arc}}{}}\frac{(f_{obs,i}f_{mod,i})^2}{2\sigma _{sky}^2}.$$
(11)
The factor of 2 in the denominator enters because the source is constructed from arc A, so it and the observed B arc each have errorbars of $`\sigma _{sky}`$. To compute $`\chi _{arc}^2`$ we mask the arc map and use only the region near the B arc, omitting regions near the B quasar and the core of the lens galaxy where imperfect subtraction makes the residual flux unreliable. We also omit the region around the object G4 where the flux is not due to the quasar host galaxy (see §2). The masked region is shown in Figure 1c. To count the number of arc constraints, we count the number of pixels inside the mask for which either the observed arc or the model arc is more than $`2\sigma `$ above the sky and take this to be the number of arc constraints $`N_{arc}`$. Typically $`N_{arc}10^4`$. Because the source is fully determined from the A arc, we quote $`\overline{\chi }_{arc}^2\chi _{arc}^2/N_{arc}`$ as an effective $`\chi ^2/\text{DOF}`$ for the arc. The fact that our best models produce $`\overline{\chi }_{arc}^21`$ (see §5) suggests that this is a reasonable counting of the constraints.
### 3.4 Lens galaxy models
Q 0957+561 has been modeled extensively; we summarize results from the recent models by Grogin & Narayan (1996), Barkana et al. (1999), Bernstein & Fischer (1999), and Chae (1999) in Table 3. (Technical details of comparing the models are discussed in Appendix A.) The models fall into two main families based on how they treat the main lens galaxy. First is the “power law” family, in which the lens galaxy is modeled with a circular or elliptical surface density with a softened power law profile: $`\mathrm{\Sigma }\left(s^2+m^2\right)^{\alpha /21}`$, where $`m`$ is an ellipsoidal coordinate, $`s`$ is a core radius, and $`\alpha `$ is the power law exponent such that $`M(R)R^\alpha `$ asymptotically. The best-fit models typically have a small core radius and a power law $`\alpha 1.1`$ corresponding to a profile slightly shallower than isothermal. In the second family, called “FGS” models after being introduced by Falco, Gorenstein & Shapiro (1991), the lens galaxy is treated as a circular or elliptical King model, plus a central point mass to account for a mass deficit in the core of the King model.
In almost all of the previous models the lens galaxy was assumed to have a projected mass distribution with circular or elliptical symmetry. However, as discussed in §§3.1 and 3.2 the observed ellipticity gradient and isophote twist suggest that this assumption may not be correct. The key improvement in the models is the addition of internal structure in the lens galaxy’s mass comparable to that seen in its light. While we do not expect the mass to trace the light exactly, it should be given the same freedoms in order to avoid oversimplifying the models.
Bernstein & Fischer (1999) did add freedom to the galaxy’s radial and angular structure by using independent power laws in different radial zones, but this led to unphysical models with density discontinuities that cannot match normal rotation curves. We introduce models that allow similar radial and angular freedom while keeping the density smooth. We start with the pseudo-Jaffe ellipsoid, whose projected density distribution has elliptical symmetry and a profile that is roughly flat inside a core radius $`s`$, falls as $`\mathrm{\Sigma }R^1`$ out to a cut-off radius $`a`$ (yielding a rotation curve that is approximately flat), and then falls as $`\mathrm{\Sigma }R^3`$ to maintain a finite mass; a detailed definition is given in Appendix B. We then construct “double pseudo-Jaffe” models comprising two concentric pseudo-Jaffe components with different scale lengths, ellipticities, and orientations. We fix the model galaxy to its observed position, which leaves 10 galaxy parameters: a mass parameter ($`b_i`$), ellipticity ($`e_i`$), orientation angle ($`\text{PA}_i`$), core radius ($`s_i`$), and cut-off radius ($`a_i`$) for each component ($`i=1,2`$).
The double pseudo-Jaffe model provides a great deal of freedom in the lens galaxy. First, it is essentially a smooth generalization of both the power law and the FGS models. When the inner pseudo-Jaffe component is compact, it mimics the point mass in FGS models while the outer component is similar to the King model (see Appendix B); and when the cut-off radius of the inner component is comparable to the core radius of the outer component, the two components combine to produce nearly a smooth $`\alpha =1`$ power law. Second, with this model we can mimic the observed internal structure of the lens galaxy: by adjusting the ellipticities and orientations we can produce a smooth ellipticity gradient and isophote twist. An important difference from the broken power law models of Bernstein & Fischer (1999) is that the double pseudo-Jaffe models have smooth density profiles. We discuss the physical properties of sample double pseudo-Jaffe models in §5.3.
### 3.5 Cluster models
Previous models have used several different methods for including the cluster’s contribution to the lensing potential. The simplest approach is to expand the cluster potential in a Taylor series and keep only the lowest significant (2nd order) term, which describes the tidal shear produced by the cluster. However, this approximation is thought to be poor for Q 0957+561 because the cluster is close to the lens and the VLBI errorbars are small, so the 3rd order terms are larger than the errorbars. Barkana et al. (1999) and Chae (1999) included all the higher order terms implicitly by introducing a mass distribution to represent the cluster, but degeneracies related to the cluster shape and density profile forced them to make assumptions about the cluster properties. Kochanek (1991) and Bernstein & Fischer (1999) instead used a Taylor series with general 3rd order terms,
$$\varphi _{clus}=\frac{1}{2}\kappa _{clus}r^2+\frac{1}{2}\gamma r^2\mathrm{cos}2(\theta \theta _\gamma )+\frac{1}{4}\sigma r^3\mathrm{sin}(\theta \theta _\sigma )\frac{1}{6}\delta r^3\mathrm{sin}3(\theta \theta _\delta ).$$
(12)
The 2nd order $`\kappa _{clus}`$ term produces the mass sheet degeneracy, so it is usually omitted from the lens models and constrained independently (see §§3.1 and 3.2). The 2nd order $`\gamma `$ term represents the tidal shear from the cluster. The 3rd order $`\sigma `$ and $`\delta `$ terms arise from the gradient of the cluster density and the $`m=3`$ component of the cluster mass (in a frame centered on the lens galaxy; see Bernstein & Fischer 1999; Keeton 2000). The three direction angles $`(\theta _\gamma ,\theta _\sigma ,\theta _\delta )`$ are written here as position angles (i.e. measured East of North).
We follow Bernstein & Fischer (1999) and use the 3rd order Taylor series (omitting the $`\kappa _{clus}`$ term). Although we simply fit for the Taylor series parameters, it is important to understand how they relate to physical properties of the cluster and what their reasonable ranges are. In popular cluster mass models the cluster amplitudes $`(\gamma ,\sigma ,\delta )`$ and direction angles $`(\theta _\gamma ,\theta _\sigma ,\theta _\delta )`$ have the following properties (Keeton 2000):
* Circularly symmetric mass distribution: the three angles all point to the center of the cluster.
* Singular isothermal ellipsoid: the shear angle $`\theta _\gamma `$ points to the cluster center and the shear amplitude equals the convergence from the cluster, $`\gamma =\kappa _{clus}`$. Both of these results hold for all positions and all values of the axis ratio $`q`$. The gradient amplitude $`\sigma `$ is bounded by $`1\frac{\sigma r_0}{\kappa _{clus}}\frac{1+q^2}{2q}`$, where $`r_0`$ is the distance from the lens galaxy to the cluster.
* Softened isothermal ellipsoid: Outside of the core the results for the singular isothermal ellipsoid are approximately true. For example, for clusters with an axis ratio $`q>0.6`$ the shear angle $`\theta _\gamma `$ points to within $`10\mathrm{°}`$ of the cluster center for $`r_02s`$, where $`s`$ is the core radius. The cluster has $`\kappa _{clus}\gamma `$, but $`\kappa _{clus}>2\gamma `$ only for $`r_03s`$.
* An ellipsoid with the “universal” dark matter profile of Navarro, Frenk & White (1996): For clusters with an axis ratio $`q>0.6`$, the shear angle $`\theta _\gamma `$ points to within $`4\mathrm{°}`$ of the cluster center for $`r_00.5r_s`$, where $`r_s`$ is the scale radius in the “universal” profile. The cluster generally has $`\kappa _{clus}\gamma `$ for $`r_0r_s`$ and $`\kappa _{clus}\gamma `$ for $`r_0r_s`$.
We let the shear and gradient parameters $`(\gamma ,\theta _\gamma ,\sigma ,\theta _\sigma )`$ vary freely, but we keep these general relations in mind when interpreting the results. To avoid an explosion of parameters, we follow Bernstein & Fischer (1999) and use a “restricted” 3rd order cluster with $`\theta _\delta =\theta _\sigma `$ and $`\delta =3\sigma /2`$ (as for a singular isothermal sphere). This model is useful because it requires only four parameters for the cluster, and it can provide model-independent evidence for a non-circular cluster if $`\theta _\sigma \theta _\gamma `$. We also experimented with models using a singular isothermal mass distribution for the cluster, and we discuss these models briefly but do not quote detailed results (§5).
### 3.6 Constraining the new models
Our models have 10 galaxy and 4 cluster parameters, but fortunately we need not examine all of them explicitly. First, the general features of the lens allow us to fix two scale radii. The lack of a central image requires that the galaxy be nearly singular, so we set the core radius of the inner pseudo-Jaffe component ($`s_1`$) to zero. Also, previous lens models suggest that the galaxy needs to have mass extending at least to the distant A image ($`5\mathrm{}`$ from the galaxy), so the cut-off radius of the outer pseudo-Jaffe component ($`a_2`$) should be larger than $`5\mathrm{}`$. The model should not be very sensitive to any particular value larger than this, so we fix $`a_2=30\mathrm{}`$ (similar to Bernstein & Fischer 1999).
Second, we can use linear techniques for the two galaxy mass parameters ($`b_1`$ and $`b_2`$) and the two cluster amplitudes ($`\gamma `$ and $`\sigma `$) because they enter the potential as simple multiplicative factors. (This is a variant of solving for the shear parameters in eq. 7.) If we take the strong VLBI constraints on the quasar core and jet positions to be exact (see Kochanek 1991; Bernstein & Fischer 1999), they lead to four linear constraint equations:
$`\text{cores}:\stackrel{}{x}_{\mathrm{A1}}\varphi (\stackrel{}{x}_{\mathrm{A1}})`$ $`=`$ $`\stackrel{}{x}_{\mathrm{B1}}\varphi (\stackrel{}{x}_{\mathrm{B1}})`$ (13)
$`\text{jets}:\stackrel{}{x}_{\mathrm{A5}}\varphi (\stackrel{}{x}_{\mathrm{A5}})`$ $`=`$ $`\stackrel{}{x}_{\mathrm{B5}}\varphi (\stackrel{}{x}_{\mathrm{B5}})`$ (14)
where A<sub>1</sub> and B<sub>1</sub> denote the two quasar cores while A<sub>5</sub> and B<sub>5</sub> denote the brightest jet components; the positions are given by Barkana et al. (1999). The four constraint equations can be solved explicitly to write the four linear parameters as functions of the remaining (non-linear) parameters. This technique ensures that the VLBI constraints are fitted exactly, while reducing by four the number of parameters that must be examined directly.
We are left with 8 explicit parameters: the ellipticities $`e_1`$ and $`e_2`$, orientation angles $`\text{PA}_1`$ and $`\text{PA}_2`$, and scale lengths $`a_1`$ and $`s_2`$ of the galaxy; and the direction angles $`\theta _\gamma `$ and $`\theta _\sigma `$ of the cluster. We evaluate models in this parameter space using the remaining constraints: the flux ratios of the quasar cores and jets and the positions of the optical Blobs and Knots (taken from Bernstein & Fischer 1999), and the structure of the arcs (see §3.3). Techniques and results for our new models are discussed in §5.
## 4 The failure of existing lens models
Previous models of Q 0957+561 (summarized in Table 3) posited that the VLBI observations and the optical Blobs and Knots provided strong enough constraints to break the degeneracy between the lens galaxy shape and the cluster shear. As we discussed in §3.1, however, it is dangerous to use smooth circular or elliptical lens models with high-precision constraints, because the oversimplified models may be forced to converge on a best-fit solution with the wrong global structure for the potential. Until now we have lacked the constraints to test this concern, although the fact that previous models gave wildly difference shapes for the lens galaxy certainly suggests that the models were not robust.
The host galaxy arcs finally provide an extensive set of constraints for testing the previous models. The first clue comes from comparing the qualitative features of the arcs and the existing models. Despite differences in details, most of the previous models have a lensing potential with a strong tidal shear from the cluster ($`\gamma 0.1`$$`0.4`$). The galaxy dominates the potential for the close image (B), but the strong shear means that the cluster dominates the potential for the distant image (A). The distinction is important because in general a round host galaxy produces an image distorted tangentially relative to the mass that dominates the potential. As a result, for a circular host galaxy most previous models of Q 0957+561 would predict a B arc tangential to the galaxy (as observed), and an A arc tangential to the cluster – i.e. stretched radially relative to the lens galaxy, opposite what is observed. These lens models can produce an A image tangential to the galaxy only if the source galaxy is highly flattened and oriented at just the right angle. In other words, we must either invoke a very special source configuration or conclude that the standard Q 0957+561 models have generic problems.
We can quantify these problems by using the arc modeling technique described in §3.3 to predict the structure of the B arc for each model and compare it to the observed B arc. (Recall that this technique uses the observed A arc to construct the source, so models always reproduce the A arc correctly and hence are evaluated by examining the B arc.) Figures 2 and 3 show the observed and predicted models arcs for previous power law models, while Figure 4 shows the arcs for FGS models; Figure 5 shows the intrinsic source inferred for two of the models. Table 3 includes the quantitative $`\overline{\chi }_{arc}^2`$ estimates for the models.
All of these models are inconsistent with the structure of the B arc, both visually and in terms of $`\overline{\chi }_{arc}^2`$. Although the predicted arcs differ among the various models, there are two common features. First, in some of the models (notably those with a circular lens galaxy) the predicted B arc has a pair of bright ridges not seen in the observed B arc. Models that produce such ridges are strongly inconsistent with the data; they have $`\overline{\chi }_{arc}^2`$ statistics no better than 5.0 and as poor as 10.8. Second, the models generally fail to predict the correct shape for the B arc at low surface brightness levels. The predicted arcs have too much curvature and do not match the northeast extension of the observed arc. The failure to match even the rough shape (extent and curvature) of the B arc illustrates the point explained above: if the potential near the A image is dominated by the cluster shear, reproducing the A arc requires a flattened source (e.g. the Barkana/SPEMD model in Figure 5) that is distorted into a highly curved B arc. The existing models cannot match the arcs better than $`\overline{\chi }_{arc}^2=2.2`$, which given the large number of arc constraints means that these models are formally excluded at an extremely high significance level.
The models arcs also illustrate a crucial limitation of the broken power law models of Bernstein & Fischer (1999), represented here by the DM2+C2 model in Figure 3. These models use independent elliptical power law models in different radial zones to mimic the ellipticity and orientation variations in the observed lens galaxy. The problem is that the density is discontinuous across the zone boundary, which leads to a discontinuity in the predicted B arc. The discontinuity does not affect the quantitative $`\overline{\chi }_{arc}^2`$ statistic because it occurs outside our mask, but it is qualitatively unacceptable. Combining this with the inability of the broken power law models to match normal rotation curves emphasizes that the models are physically unacceptable. The need for a more physically reasonable way to include ellipticity and orientation variations motivates our introduction of the double pseudo-Jaffe models.
Although the previous models do not correctly reproduce the distortions of the host galaxy, they do reveal a model-independent qualitative feature of the lensing. Part of the quasar host galaxy is doubly-imaged like the quasar, but a small region crosses the lensing caustic and is quadruply-imaged (see Figure 5). Arc A is a single distorted image of the host galaxy, but arc B is a combination of the three remaining images of the quadruply-imaged and the one remaining image of the doubly-imaged region. This fact explains why the A and B arcs have such different geometries. It also implies that with the current observations of the arcs, most of the constraints come from the small region of the host galaxy that is quadruply-imaged. A deeper image of the arcs should reveal a complete Einstein ring that would strengthen the constraints by using more of the host galaxy (see §6).
In summary, the host galaxy arcs in Q 0957+561 offer a strong new probe of the global shape of the lensing potential, and they show that existing lens models converged to the wrong potential. Generically, the previous models have too much shear from the cluster. The failure to fit the arcs is a common problem of models that combine strong constraints from the sub-milli-arcsecond structure of the quasar jets with oversimplified circular or elliptical lens models. It rules out existing models of the system, together with the bounds on $`H_0`$ drawn from them.
## 5 Successful new models
Two goals motivate our search for new models of Q 0957+561. First, if we want to use the VLBI constraints with sub-milli-arcsecond precision, we must be sensitive to details of the lens galaxy structure such as its radially varying ellipticity and orientation. We used double pseudo-Jaffe models (see §3.4) to incorporate such structure in a smooth and physically reasonable way. Second, we want to incorporate the powerful constraints from the host galaxy arcs into the modeling process, instead of using them for a posteriori tests as in §4. In this way we hope to break the environment degeneracy in a robust way. When we began to examine new models we found a range of solutions consistent with the data, so we adopted Monte Carlo techniques to sample this range. In this section we first describe the Monte Carlo techniques (§5.1), and then discuss the models and their implications for breaking the degeneracy between the galaxy and cluster (§5.2), the physical properties of the lens galaxy and cluster (§5.3), and the Hubble constant (§5.4).
### 5.1 Monte Carlo techniques
To explore our model space, we picked random values for the non-linear model parameters, using restricted but physically-motivated ranges:
* Galaxy ellipticity: There are no successful models with an outer ellipticity below 0.3. We believed models with an ellipticity larger than $`0.7`$ to be implausible. Hence we considered $`0.3e_20.7`$ for the outer ellipticity. In models with an ellipticity gradient the inner ellipticity can be small, so we considered $`0e_10.7`$.
* Galaxy orientation: In the observed lens galaxy the PA varies from $``$40° to $``$60° (Bernstein et al. 1997). We expect the mass to be roughly aligned with the light (to within $``$10°; Keeton et al. 1998), so we considered $`30\mathrm{°}(\text{PA}_1,\text{PA}_2)70\mathrm{°}`$ for the inner and outer position angles.
* Galaxy scale lengths: We fixed the inner pseudo-Jaffe component to be singular ($`s_1=0`$) and the outer component to have a cut-off radius $`a_2=30\mathrm{}`$ (see §3.6). For the cut-off radius of the inner component ($`a_1`$) and the core radius of the outer component ($`s_2`$) we considered $`0\stackrel{}{\mathrm{.}}1(a_1,s_2)4\mathrm{}`$.
* Cluster direction angles: Galaxy counts and the weak lensing measurement suggest that the cluster has its mass concentration toward $`\theta _{clus}55\mathrm{°}`$ (Fischer et al. 1997). For reasonable cluster models the shear angle $`\theta _\gamma `$ points roughly to the cluster center (see §3.5), so we considered shear angles in the range $`30\mathrm{°}\theta _\gamma 70\mathrm{°}`$. We were less restrictive with the gradient angle $`\theta _\sigma `$ and required only that it be in the same quadrant as the cluster, $`0\mathrm{°}\theta _\sigma 90\mathrm{°}`$.
Given values for these parameters, we fixed the four remaining parameters ($`b_1`$, $`b_2`$, $`\gamma `$, $`\sigma `$) using the constraints from the quasar core and jet positions (see §3.6). Our limited parameter ranges may omit some models that are formally consistent with the data, but they span what we expect for physically reasonable models.
We tabulated models that fit the data at the 95% confidence level. With $`N=(2,4,6)`$ constraints this is equivalent to $`\chi ^2(5.99,9.49,12.59)`$ (e.g. Press et al. 1992, §§6.2 and 15.6). We applied these thresholds to the constraints from the quasars and the optical Blobs and Knots both separately and jointly: we required $`\chi _{flux}^25.99`$ for the flux ratios of the quasar cores and jets (2 constraints), $`\chi _{bk}^29.49`$ for the positions of the optical Blobs and Knots (4 constraints), and $`\chi _{flux}^2+\chi _{bk}^212.59`$ for the 6 joint constraints. We did not impose constraints from the less reliable Blob and Knot flux ratios. We considered different thresholds for the arc constraints, as discussed below. Note that it was not useful to combine all the constraints into a total $`\chi ^2/\text{DOF}`$ because the number of constraints from the arcs is so large.
The Monte Carlo technique is not especially efficient; we examined $`10^7`$ models and found of order 1000 models consistent with the data. However, examining that many models is not prohibitively time consuming. Most of the models can be ruled out quickly because they fail the flux and Blob/Knot $`\chi ^2`$ cuts. For the remaining models, we first computed a fast $`\overline{\chi }_{arc}^2`$ using a 4 times undersampled arc map; only for promising models did we compute $`\overline{\chi }_{arc}^2`$ using the full arc map. The benefit of this brute force approach is the ability to identify and sample the wide range of models consistent with the data, and thereby estimate the full range of possible $`H_0`$ values.
We first considered models in which the inner and outer components of the model galaxy were fixed to have the same ellipticity and orientation; we found $`(3,275,832,1575)`$ models with $`\overline{\chi }_{arc}^2<(1.2,1.3,1.4,1.5)`$, whose properties are summarized in Figure 6. Next we allowed the inner and outer galaxy components to have different shapes, mimicking ellipticity and orientation gradients. For simplicity, we refer to such models as having a “twist,” even though that term formally describes only an orientation gradient. We found $`(25,134,286,477)`$ twist models with $`\overline{\chi }_{arc}^2<(1.2,1.3,1.4,1.5)`$, whose properties are summarized in Figure 7. Table 4 gives parameters and Figure 8 shows the predicted arcs for sample models of both types. We note that the number of successful models with a twist is smaller than the number of successful models without a twist. Introducing the twist enlarges the parameter space without necessarily increasing the space of successful models by the same amount; hence sampling the same total number of models yields fewer acceptable models. In the following sections we discuss models with $`\overline{\chi }_{arc}^2<1.5`$, but our conclusions would not change substantially if we lowered this threshold.
### 5.2 Breaking the galaxy/cluster degeneracy
The host galaxy arcs finally break the degeneracy between the lens galaxy shape and the cluster shear in Q 0957+561. First, they constrain the lens galaxy shape. Among the previous models, the three with the best fits to the arcs have $`62\mathrm{°}<\text{PA}<67\mathrm{°}`$ (Table 3). In the new models, the bounds on the (outer) orientation of the lens galaxy are $`55\mathrm{°}<\text{PA}<66\mathrm{°}`$ in models without a twist (Figure 6d) and $`52\mathrm{°}<\text{PA}_2<68\mathrm{°}`$ in models with a twist (Figure 7f). These ranges are consistent with the observed galaxy’s orientation of about $`56\pm 8\mathrm{°}`$ (Bernstein et al. 1997). In other words, the arcs not only constrain the model galaxy’s orientation, they require that its mass distribution be at least roughly aligned with the light distribution; such alignment is also seen in other lens systems (Keeton et al. 1998). Also, in the new models the bounds on the (outer) ellipticity of the galaxy are $`0.38<e<0.61`$ in models without a twist (Figure 6c) and $`0.35<e_2<0.63`$ in models with a twist (Figure 7e). These ranges are slightly higher than the ellipticity seen in the outer parts of the galaxy (Bernstein et al. 1997), but there is no a priori reason to expect the ellipticities of the dark halo and the light to be similar. These constraints on the lens galaxy shape span models with a wide range of lens galaxy mass profiles (controlled by the scale radii $`a_1`$ and $`s_2`$; see Figures 6b and 7b), ellipticity profiles (with or without a twist), and cluster contributions. Thus we believe that they are robust, and that the ability of the host galaxy arcs to finally constrain the galaxy shape is the most important result of our new models.
Second, the arcs place limits on the lensing contribution of the cluster. The new models have moderate or even small cluster shears (typically $`\gamma 0.05`$; Figures 6e and 7g), which contrasts with the strong shears in previous models. A small shear is required by the fact that both arcs are tangential to the lens galaxy (see §4). Also, including the 3rd order cluster term is important, as expected from the proximity of the cluster to the lens (Kochanek 1991). Most models have a cluster gradient amplitude in the range $`0.006<\sigma <0.016`$ (Figures 6g and 7i), so at the position of image A the 3rd order term is comparable to or even larger than the 2nd order term (see eq. 12). Because the cluster amplitudes are small, the corresponding direction angles $`\theta _\gamma `$ and $`\theta _\sigma `$ are poorly constrained (Figures 6f, 6h, 7h, and 7j). Nevertheless, the fact that they usually differ implies that the cluster cannot be spherical (see §3.5).
Unfortunately, even the arcs cannot uniquely determine the model. First, the cluster angles can adjust to accommodate the narrow but finite range of lens galaxies given above. Second, in models with a twist the inner component can take on all values of ellipticity and orientation in the ranges we allowed (Figures 7c and 7d). This is because the inner component is usually compact, so like the point mass in FGS lens models it serves mainly to correct a central mass deficit and the model is largely insensitive to the distribution of that mass. Both of these degeneracies affect the potential enough to produce a disappointingly wide range of Hubble constant values (see §5.4).
These conclusions are drawn from models using a Taylor series for the cluster potential. We also examined models treating the cluster using an elongated isothermal model with an arbitrary position, axis ratio, and orientation.<sup>5</sup><sup>5</sup>5For technical reasons, the isothermal cluster models we used had only approximate elliptical symmetry. The above conclusions about the galaxy properties still pertain, with one modification: in models with a twist the distribution of outer orientation angles gains a tail down to $`\text{PA}_240\mathrm{°}`$. The above conclusions about the cluster contributions also hold, except that the new bounds on the shear amplitude are $`0.03<\gamma <0.11`$. Since an isothermal cluster has $`\gamma \kappa _{clus}`$, it is hard to obtain a small shear without essentially eliminating the cluster.
### 5.3 Properties of the galaxy and cluster
We introduced the double pseudo-Jaffe models in order to smoothly include ellipticity and orientation variations in the lens galaxy, but we must ask whether the models are physically plausible. Figure 9 shows the estimated density and circular velocity profiles for the models in Table 4. These profiles are only estimates because they require the full 3-d mass distribution, while lens models give only the projected distribution. We assumed that the intrinsic distributions are oblate spheroids viewed edge-on, and we computed the density and circular velocity profiles in the equatorial plane, neglecting any twist. The high inferred circular velocities ($`v_c600`$ km s<sup>-1</sup>) are misleading, because our models omit any convergence from the cluster and $`v_c(1\kappa _{clus})^{1/2}`$, and because converting to an observed velocity dispersion involves systematic uncertainties in the stellar dynamics (see Romanowsky & Kochanek 1999). The qualitative features, however, are plausible. The inner and outer galaxy components generally combine to produce a rotation curve that is approximately flat, although in several models it turns up at small radii because the inner component of the pseudo-Jaffe models is compact. This effect, which also appears in FGS models because of the point mass, occurs because lensing constrains only the total mass enclosed by the B image ($`r_B=1\stackrel{}{\mathrm{.}}0`$), and if that mass is compact then it produces a rising rotation curve. We conclude that the fitted double pseudo-Jaffe models are fairly reasonable.
As for the cluster, recall that the host galaxy arcs require the shear to be small. At the same time, two independent estimates of the cluster mass suggest that the convergence from the cluster is $`\kappa _{clus}0.2`$. First, stellar dynamical modeling of the lens galaxy constrains its mass such that explaining the image separation requires a cluster convergence of $`\kappa _{clus}=0.20\pm 0.12`$ (Romanowsky & Kochanek 1999). Second, Fischer et al. (1997) and Bernstein & Fischer (1999) used weak lensing to measure the total mass inside an aperture centered on the lens galaxy. Subtracting the mass of a model lens galaxy then gives a model-dependent estimate of the remaining cluster mass (see eq. 10). With this technique all of our models yield $`0.20\kappa _{clus}0.23`$.
These results are surprising because with popular cluster models it is difficult to produce a shear that is much smaller than the convergence (Keeton 2000). Specifically, with a singular isothermal ellipsoid, $`\gamma =\kappa _{clus}`$ for all cluster positions, orientations, and axis ratios. Introducing a core radius allows $`\gamma \kappa _{clus}`$, but the shear is substantially smaller than the convergence only if the lens galaxy is near or within the cluster core. If the cluster halo has the “universal” or NFW density profile (Navarro et al. 1996), $`\gamma \kappa _{clus}/2`$ only inside $`0.5r_s`$, where $`r_s`$ is the NFW scale length. For the observed cluster at redshift $`z_l=0.36`$ with a velocity dispersion of $``$700 km s<sup>-1</sup> (Angonin-Willaime et al. 1994), the NFW scale length would be $`r_s100`$ $`h^1`$ kpc or $`30\mathrm{}`$ (e.g. Navarro et al. 1996). Combining these results with the estimates of the cluster contribution ($`\gamma 0.1`$ or even $`0.05`$, and $`\kappa _{clus}0.2`$) suggests two possible conclusions about the cluster. On the one hand, if the cluster is ellipsoidal it must have a core or scale radius large enough to encompass the galaxy. Having the lens galaxy be in the cluster core would not be too surprising since it is the brightest galaxy in the cluster. On the other hand, the assumption of an ellipsoidal model for the cluster may not be correct, perhaps because of substructure.
### 5.4 Implications for the Hubble constant
Finally, we want to combine the lens models and the observed time delay and use eq. (3) to determine the Hubble constant $`H_0`$. As discussed in §3.1, the inferred value of $`H_0`$ depends weakly on the other cosmological parameters; we quote results assuming $`\mathrm{\Omega }_0=1`$ and $`\mathrm{\Lambda }_0=0`$ and note that they would increase by 5.8% (4.5%) for an $`\mathrm{\Omega }_0=0.3`$ open (flat) cosmology. Also, the mass sheet degeneracy implies $`H_0(1\kappa _{clus})`$ (see eq. 6), and $`\kappa _{clus}`$ cannot be constrained by lens models. Thus if we write $`H_0=100h\text{km s}\text{-1}\text{ Mpc}\text{-1}`$, lensing directly measures only the combination $`h/(1\kappa _{clus})`$, which is what we quote in Tables 3 and 4 and Figures 6 and 7. We need an independent estimate of $`\kappa _{clus}`$ in order to constrain $`H_0`$ itself.
We saw in §5.2 that the host galaxy arcs help break the degeneracy between the lens galaxy shape and the cluster shear, but they do not eliminate it entirely. The remaining freedom in the models is small in terms of the properties of the lens galaxy but large in terms of the Hubble constant. Models without a twist yield $`1.1h/(1\kappa _{clus})1.4`$, with a tail down to $`h/(1\kappa _{clus})0.95`$ (Figure 6). Allowing a twist in the lens galaxy broadens the distribution down to $`h/(1\kappa _{clus})0.85`$ (Figure 7). In other words, within the lens models there is a $`\pm 25\%`$ variation in inferred values for the Hubble constant; this is independent of uncertainties in $`\kappa _{clus}`$. Most of this variation is related to a strong correlation between the lens galaxy ellipticity and the Hubble constant (Figures 6c, 7c, and 7e): flatter galaxies have deeper central potential wells when normalized to produce the same images.
Although the model uncertainties limit our ability to use Q 0957+561 to constrain the Hubble constant, we might invert our thinking and ask what properties of the lens model are required to be consistent with local distance ladder determinations of $`H_0`$ (e.g. Mould et al. 1999). Assuming a cluster convergence $`\kappa _{clus}0.2`$ (see §5.3), obtaining $`H_080`$ km s<sup>-1</sup> Mpc<sup>-1</sup> requires a lens galaxy that is relatively round in the center ($`e_10.30`$) and moderately flattened in the outer parts ($`0.35e_20.56`$). These bounds bolster the suggestion that the lens galaxy must have a mass distribution with an ellipticity gradient fairly similar to that seen in the light distribution.
## 6 Conclusions
We have detected large, distorted images of the quasar host galaxy in the gravitational lens Q 0957+561. In the H band (rest frame R), the host galaxy appears as two long ($``$5″) arcs stretched tangentially relative to the lens galaxy. Previously published models of Q 0957+561 fail to predict the correct shape for the host galaxy arcs. This failure rules out those models and any conclusions about the value of the Hubble constant drawn from them.
The problem with the previous models is that they oversimplified the mass distribution for the lens galaxy, so they had to adjust the large-scale features of the model (the galaxy shape and the cluster shear) in order to fit small-scale constraints (the VLBI jets, with $``$0.1 milli-arcsecond errorbars). Without strong constraints on the global shape of the lensing potential, each class of models happily converged to a best-fit solution with the wrong global structure. The failure of these models presents two important lessons (also see Kochanek 1991; Bernstein et al. 1993; Mao & Schneider 1998; Bernstein & Fischer 1999). First, when the constraints are extremely precise it is necessary to include the full complexity of the structure of the lens galaxy. Second, it is important to carefully explore the full range of possible models before leaping to conclusions about the value of the Hubble constant and its uncertainties.
Two improvements are crucial to finding better models. First, the host galaxy arcs finally provide enough constraints to separately determine both the lens galaxy shape and the cluster shear, and they must be incorporated into the modeling process. Second, models for the lens galaxy must allow internal structure similar to that seen in the observed galaxy, which for Q 0957+561 means a radially varying ellipticity and orientation. Although the mass need not exactly trace the light, it must be given the same freedoms lest we oversimplify the models. We introduced a new class of lens models, the double pseudo-Jaffe models, that smoothly incorporate an ellipticity gradient and isophote twist.
Our models lead to several new conclusions about the system. First, the shape of the lens galaxy’s mass distribution must be surprisingly similar to the shape of its luminosity distribution. The mass distribution must be moderately flattened and roughly aligned with the light distribution. Such alignment between the mass and the light is also seen in other gravitational lenses (Keeton et al. 1998), but was not seen in previous models of Q 0957+561. Second, the shear from the cluster must be small ($`\gamma 0.1`$), in marked contrast with the moderate or strong shears seen in previous models ($`\gamma 0.1`$$`0.4`$). The small shear combined with an estimated convergence $`\kappa _{clus}0.2`$ implies that the cluster potential must be approximately centered on the lens galaxy. Unfortunately, the current observations of the host galaxy cannot fully determine the lens model. There are still freedoms related to several direction angles that describe the cluster and to the amplitudes of the ellipticity and orientation gradients, and they leave a 25% uncertainty in the inferred value of the Hubble constant.
However, there are two promising prospects for further improving the constraints on this systems. First, the substantial progress we have found is based on a 2800-second H band image. A deeper image with a refurbished NICMOS camera should show a complete Einstein ring image. Since an Einstein ring probes the potential all the way around the lens galaxy, it is extremely useful for determining both the lens model and the intrinsic shape of the host galaxy (see Keeton, Kochanek & McLeod 2000). Filling in the gaps in the Q 0957+561 ring should eliminate the remaining uncertainties in the models. Second, X-rays from the cluster gas have been detected (Chartas et al. 1998), but the image resolution was poor and the signal was dominated by X-rays from the quasar images. New high-resolution X-ray observations to map the cluster gas would constrain the cluster potential and help determine the cluster angles that are still unknown (namely the angle to the cluster center and the angle of the cluster’s density gradient at the lens). Combining deeper infrared and X-ray imaging would thus dramatically improve the constraints on the models and allow consistency checks of the cluster potential. Improved X-ray imaging would also further improve the mass estimates of the cluster for breaking the mass sheet degeneracy.
Distorted images of the host galaxy have now been observed in four of the time delay lenses (Q 0957+561, PG 1115+080, B 1600+434, and B 1608+656; see Impey et al. 1998; Kochanek et al. 1999). As we have illustrated, images of the host galaxy are a powerful constraint on models of the system and hence on the uncertainties in the value of $`H_0`$ derived from the time delay measurements. To date, host galaxies have been thought of mainly as pleasant bonuses in relatively shallow images targeting the lens galaxies. Taking full advantage of the host galaxies will require deeper images focused on the host galaxies themselves. Such images hold great promise for breaking common degeneracies in lens models and allowing gravitational lensing to map in detail the mass distributions of distant galaxies, to probe the potentials of lens galaxy environments, and to determine a robust and independent measurement of the Hubble constant at cosmological distances.
###### Acknowledgements.
Acknowledgements: We thank Gary Bernstein for discussions and comparisons with new optical data. Support for the CASTLES project was provided by NASA through grant numbers GO-7495 and GO-7887 from the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555.
## Appendix A Appendix: Comparing lens models
Table 3 summarizes results from models of Q 0957+561 by Grogin & Narayan (1996), Barkana et al. (1999), Bernstein & Fischer (1999), and Chae (1999). The various authors reported results in different ways. Sorting out the differences can be confusing, so we converted all the results to a standard form and explain the conversions here.
First, Bernstein & Fischer (1999) and Chae (1999) quoted an ellipticity parameter $`ϵ`$ that is related to the axis ratio $`q`$ by $`q^2=(1ϵ)/(1+ϵ)`$. We report the true ellipticity $`e=1q`$. (Bernstein & Fischer called their ellipticity parameter $`e`$, but it played the same role as $`ϵ`$ here.)
Second, different authors quoted results for the Hubble constant in different ways. In general, lens models can determine only the combination $`h/(1\kappa _{clus})`$ where $`h=H_0/(100\text{km s}\text{-1}\text{ Mpc}\text{-1})`$ and $`\kappa _{clus}=\mathrm{\Sigma }_{clus}/\mathrm{\Sigma }_{crit}`$ is the surface mass density of the cluster (in critical units) at the position of the lens. Models that use a multipole expansion for the cluster (an external shear or 3rd order cluster, see eq. 12) offer no constraints on $`\kappa _{clus}`$. Models that use an actual mass distribution for the cluster do predict $`\kappa _{clus}`$: an isothermal sphere that has core radius $`s`$ and is located a distance $`d`$ from the lens produces $`\kappa _{clus}=b/(2\xi )`$ and shear $`\gamma _{clus}=(bd^2)/[2\xi (s+\xi )^2]`$, where $`\xi =\sqrt{s^2+d^2}`$ and $`b`$ is a mass parameter such that $`M(R)=\pi \mathrm{\Sigma }_{crit}b(\sqrt{s^2+R^2}s)`$. However, the $`\kappa _{clus}`$ prediction is not unique because varying the cluster’s shape and profile can change $`\kappa _{clus}`$ (and $`\gamma _{clus}`$) without changing the goodness of fit (see Chae 1999 for examples). Hence we describe these cluster models in terms of an equivalent external shear with magnitude $`\gamma _{\mathrm{eff}}=\gamma _{clus}/(1\kappa _{clus})`$ (see Grogin & Narayan 1996; Barkana et al. 1999). In Table 3 we quote the magnitude of the external shear ($`\gamma `$, or $`\gamma _{\mathrm{eff}}`$ in the isothermal sphere cluster models) and the Hubble constant combination $`h/(1\kappa _{clus})`$. We note that if the cluster is a singular isothermal sphere, $`(1\kappa _{clus})=(1+\gamma _{\mathrm{eff}})^1`$ (see Barkana et al. 1999). The estimates of $`h/(1\kappa _{clus})`$ can be translated into actual estimates for $`H_0`$ using the weak lensing measurement of the cluster mass (see §3.2).
Third, different authors used slightly different definitions of $`\chi ^2`$. We did not convert $`\chi ^2`$ values; we used what the authors reported.
Finally, we corrected several apparent typographical errors. Barkana et al. (1999) claimed that their galaxy orientation angle $`\theta `$ was a position angle, but in fact the position angle was $`90\mathrm{°}\theta `$ and this is what we quote. Bernstein & Fischer (1999) quoted angles measured from the positive $`x`$-axis. Their formalism is self-consistent, except that their eq. (16) needs a minus sign in front of the $`\gamma `$ term. We report the Bernstein & Fischer (1999) angles as position angles. Finally, for the Bernstein & Fischer (1999) 3rd order cluster models there appears to be a difference of a factor of 2 between the amplitude of 3rd order cluster term as written in their eq. (16) and as reported in their Table 3. We quote values consistent with their eq. (16).
## Appendix B Appendix: The pseudo-Jaffe ellipsoid
A standard Jaffe (1983) model has a 3-dimensional density distribution $`\rho m^2(m+a)^2`$ where $`m`$ is an ellipsoidal coordinate and $`a`$ is the break radius. For lensing it is more convenient to use a modified density distribution $`\rho (m^2+s^2)^1(m^2+a^2)^1`$ where $`a`$ is again the break radius, and we have added a core radius $`s<a`$. For this model the projected surface mass density, in units of the critical surface density for lensing, is
$$\frac{\mathrm{\Sigma }}{\mathrm{\Sigma }_{crit}}=\frac{b}{2}\left[\left(m^2+s^2\right)^{1/2}\left(m^2+a^2\right)^{1/2}\right],$$
(B1)
where the mass normalization parameter $`b`$ is chosen to match the lensing critical radius in the limit of a singular isothermal sphere ($`s0`$, $`a\mathrm{}`$, $`q1`$). Eq. (B1) defines what we call the pseudo-Jaffe ellipsoid. Its projected surface density is roughly constant for $`Rs`$, falls as $`R^1`$ for $`sRa`$, and falls as $`R^3`$ for $`Ra`$; its total mass is $`M=\pi \mathrm{\Sigma }_{crit}qb(as)`$. The ellipsoid coordinate $`m`$ can be written in terms of the projected axis ratio $`q`$ and the position angle PA (measured East of North) as
$$m^2=\frac{R^2}{2q^2}\left[\left(1+q^2\right)+\left(1q^2\right)\mathrm{cos}2\left(\theta \text{PA}\right)\right].$$
(B2)
The pseudo-Jaffe lens model is easy to compute because it is written as the difference of two softened isothermal ellipsoids, whose analytic lensing properties are known (Kassiola & Kovner 1993; Kormann, Schneider & Bartelmann 1994; Keeton & Kochanek 1998). This model has been used previously by de Zeeuw & Pfenniger (1988), Brainerd, Blandford & Smail (1996), and Keeton & Kochanek (1998).
In the limit of a singular model ($`s=0`$), the pseudo-Jaffe model is an example of a general class of “cuspy” lens models with $`\rho m^\gamma (m^2+a^2)^{(\gamma 4)/2}`$. These models are a more realistic family for real galaxies than the softened power law models, and are discussed by Muñoz, Kochanek & Keeton (2000).
The pseudo-Jaffe model is related to the King model used in FGS lens models. The standard approach is to approximate the King model with a combination of isothermal models (Young et al. 1980),
$$\frac{\mathrm{\Sigma }}{\mathrm{\Sigma }_{crit}}=\frac{2.12b}{\sqrt{m^2+0.75r_s^2}}\frac{1.75b}{\sqrt{m^2+2.99r_s^2}}.$$
(B3)
The King model has a single scale radius $`r_s`$, while the pseudo-Jaffe model has independent core and break radii. Also, in the King model the coefficients of the two terms differ, while in the pseudo-Jaffe model they are the same. |
warning/0001/nlin0001051.html | ar5iv | text | # Semiclassical Trace Formulas for Two Identical Particles
## I Introduction
Semiclassical physics has experienced a resurgence of interest, largely due to the work of Gutzwiller , Balian and Bloch and Berry and Tabor . (For recent reviews see .) These works showed that if we separate the density of states into smooth and oscillatory components, then the oscillatory part is related to the dynamics of the underlying classical system via periodic orbits. This complements the earlier work of Weyl, Wigner, Kirkwood and others who showed that the smooth component is related to the geometry of the classical phase space. Actually, the two components are related in a subtle way since the complete geometry imparts the full dynamics and vice-versa.
Most of the theoretical work has concentrated on the single particle density of states, however, there are some notable exceptions, namely . In the focus is on the average level density and its extension to systems of identical particles. Specifically, the authors consider a system of $`N`$ fermions in one dimension. Their Weyl formula for fermions works well for attractive two body interactions, but overestimates the quantum staircase function when there are repulsive two body interactions. The author of develops a generalization of the canonical periodic orbit sum for the special case of $`N`$ interacting spinless fermions in one dimension. It is assumed the periodic orbits are isolated and therefore it is most applicable to fully chaotic systems. The author also considers a system of noninteracting fermions and writes the many body level density as a convolution integral involving one body level densities. Finally, we mention which presents an expansion of the periodic orbit sum in terms of the particle number using ideas from .
Similarly, most of the applications of semiclassical theory have been to systems which involve only single particle dynamics. Here, we mention some exceptions. The authors of extend the study of scars to classically chaotic few body systems of identical particles. A study of the eigenfunctions of an interacting two particle system can be found in . The semiclassical approach to the helium atom, which can be understood as two interacting electrons in the presence of a helium nucleus, has been studied in . We also mention the novel applications of semiclassical theory to mesoscopic physics . For example, orbital magnetism has been studied semiclassically for diffusive systems in and for ballistic systems in . For reviews see .
Ultimately, one would like to study an arbitrary number of interacting particles in any kind of potential. In the present work, we begin by exploring the structure of the trace formula for two noninteracting particles including an examination of the decomposition into bosonic and fermionic spaces. This sets the stage for the interacting $`N`$body problem to be explored in a subsequent publication . The method employed here uses the fact that the two particle density of states is the autoconvolution of the single particle density of states. Subsequently, we decompose the semiclassical two particle density of states into three distinct contributions, each of which corresponds to specific dynamical properties of the system. Of particular interest is the contribution which corresponds to two particle dynamics.
Billiards have served as prominent model systems in quantum chaos. They combine conceptual simplicity (the model of a free particle in a box) while allowing the full range of classical dynamics, from integrable to chaotic. Therefore, as initial applications of the formalism, we study two noninteracting identical particles in a disk and in a cardioid. The former problem is integrable while the second is chaotic so these two examples provide a direct test of the formalism in the two limiting cases of classical motion. In both cases, we find the semiclassical formalism does a good job of reproducing the quantum density of states.
## II Background Theory
### A Single Particle Semiclassical Theory
In this section, we review the formalism for the semiclassical decomposition of the single particle density of states. Let $`\left\{ϵ_i\right\}`$ be the single-particle energies so that the single particle density of states is
$$\rho _1(ϵ)=\underset{i}{}\delta (ϵϵ_i),$$
(1)
where the subscript $`1`$ indicates that it is a single particle density. A fundamental property of the quantum density of states is that it can be exactly decomposed into an average smooth part and an oscillatory part
$$\rho _1(ϵ)=\overline{\rho }_1(ϵ)+\stackrel{~}{\rho }_1(ϵ).$$
(2)
There are various approaches for calculating these quantities . For example, in systems with analytic potentials, the smooth part may be obtained from an extended Thomas-Fermi calculation which is an asymptotic expansion in powers of $`\mathrm{}`$. In billiard systems, where the particle is confined to a spatial domain by the presence of infinitely steep potential walls, the smooth part may be obtained from the Weyl expansion. In two dimensional billiards with piecewise smooth boundaries and Dirichlet boundary conditions, the first three terms of the Weyl expansion is
$$\overline{\rho }_1(ϵ)=\left(\frac{\alpha 𝒜}{4\pi }\frac{\alpha ^{1/2}}{8\pi }\frac{}{\sqrt{ϵ}}\right)\theta (ϵ)+𝒦\delta (ϵ)+\mathrm{}$$
(3)
where $`\alpha =2m/\mathrm{}^2`$, $`𝒜`$ is the area, $``$ is the perimeter and
$$𝒦=\frac{1}{12\pi }dl\kappa (l)+\frac{1}{24\pi }\underset{i}{}\frac{\pi ^2\theta _i^2}{\theta _i}$$
(4)
is the average curvature integrated along the boundary with corrections due to corners with angles $`\theta _i`$. The oscillating part is obtained from semiclassical periodic orbit theory, and in particular the various trace formulas for $`\stackrel{~}{\rho }_1(ϵ)`$ of the form
$$\stackrel{~}{\rho }_1(ϵ)\frac{1}{\pi \mathrm{}}\underset{\mathrm{\Gamma }}{}A_\mathrm{\Gamma }(ϵ)\mathrm{cos}\left(\frac{1}{\mathrm{}}S_\mathrm{\Gamma }(ϵ)\sigma _\mathrm{\Gamma }\frac{\pi }{2}\right).$$
(5)
$`\mathrm{\Gamma }`$ denotes topologically distinct periodic orbits, $`S_\mathrm{\Gamma }(ϵ)`$ is the classical action integral along the orbit $`\mathrm{\Gamma }`$. The amplitude $`A_\mathrm{\Gamma }(ϵ)`$ depends on energy, the period of the corresponding primitive orbit, the stability of the orbit, and whether it is isolated or non-isolated. The index $`\sigma _\mathrm{\Gamma }`$ depends on the topological properties of each orbit. For isolated orbits, it is just the Maslov index. For nonisolated orbits, there may be additional phase factors in the form of odd multiples of $`\pi /4`$ which we account for, in a slight abuse of notation, by allowing $`\sigma _\mathrm{\Gamma }`$ to be half-integer. In the case of non-isolated orbits, $`\mathrm{\Gamma }`$ denotes distinct families of degenerate orbits. The amplitude of an isolated orbit is given by the Gutzwiller trace formula
$$A_\mathrm{\Gamma }(ϵ)=\frac{T_\gamma (ϵ)}{\sqrt{\left|det(\stackrel{~}{M}_\mathrm{\Gamma }I)\right|}}$$
(6)
where $`T_\gamma (ϵ)`$ is the period of the primitive orbit $`\gamma `$, corresponding to $`\mathrm{\Gamma }`$ (i.e. $`\mathrm{\Gamma }`$ is an integer repetition of $`\gamma `$) and $`\stackrel{~}{M}_\mathrm{\Gamma }`$ is the stability matrix of that orbit.
### B Quantum Two Particle Density of States
Now suppose we have a system of two identical noninteracting particles. The total Hamiltonian is the sum of the single particle Hamiltonians and it follows that the energies of the composite system are just the sums of the single particle energies. The analogue of (1) is then
$$\rho _2(E)=\underset{i,j}{}\delta (E(ϵ_i+ϵ_j)).$$
(7)
A useful relation is that the two particle density of states is the autoconvolution of the single particle density of states:
$$\rho _2(E)=_0^Edϵ\rho _1(ϵ)\rho _1(Eϵ)=\rho _1\rho _1(E),$$
(8)
as can be verified by direct substitution. In fact, this works even if the particles are not identical, where the full density is still the convolution of the two distinct single-particle densities. This would also apply to a single particle in a separable potential, which is mathematically equivalent. Rather than encumber the notation to explicitly allow for this possibility, we defer this discussion to Appendix A, where some formulas for nonidentical, noninteracting particles are presented.
We can decompose the two particle density of states for a system of two identical particles into a symmetric and an antisymmetric density,
$$\rho _2(E)=\rho _S(E)+\rho _A(E).$$
(9)
We shall use the terms symmetric/antisymmetric and bosonic/fermionic interchangeably. Each partial density may be obtained using a projection operator onto the relevant subspaces resulting in
$$\rho _{S/A}(E)=\frac{1}{2}\left(\rho _2(E)\pm \frac{1}{2}\rho _1\left(\frac{E}{2}\right)\right).$$
(10)
We seek semiclassical approximations to these quantum expressions, a topic which is pursued in the following sections.
## III Semiclassical Calculations for the Two Particle System
Decomposing the single particle density into its smooth and oscillatory components as in (2) gives a decomposition of the two particle density of states into three distinct contributions,
$$\rho _2^{\mathrm{sc}}(E)=\overline{\rho }_1\overline{\rho }_1(E)+2\overline{\rho }_1\stackrel{~}{\rho }_1(E)+\stackrel{~}{\rho }_1\stackrel{~}{\rho }_1(E).$$
(11)
The first term is a smooth function of energy since the convolution of two smooth functions results in a smooth function. This is followed by a cross term and finally by a purely oscillating term. The cross term is also an oscillating function. At first, this may seem incorrect since the convolution of a smooth function with an oscillating function usually yields a smooth function. As we will show, the oscillatory nature of the cross term is due to contributions from the end-points of the convolution integral. Physically, the smooth term does not depend on dynamics since it corresponds to the Weyl formula in the full two-particle space. The cross term depends only on single particle dynamics because it corresponds to the situation where one particle is stationary and the other particle is evolving dynamically on a periodic orbit. It is only the last term which contains two particle dynamics in the sense that both particles are evolving dynamically on periodic orbits. Hence, we will refer to the last term as the dynamical term.
We find a general expression for $`\stackrel{~}{\rho }_1\stackrel{~}{\rho }_1(E)`$ by substituting a generalized trace formula for $`\stackrel{~}{\rho }_1(E)`$ and then evaluating the resulting convolution integral using the method of stationary phase. Using (5), the dynamical term can be written as
$`\stackrel{~}{\rho }_1\stackrel{~}{\rho }_1(E){\displaystyle \frac{1}{(\pi \mathrm{})^2}}{\displaystyle \underset{\mathrm{\Gamma }_1,\mathrm{\Gamma }_2}{}}{\displaystyle _0^E}dϵA_{\mathrm{\Gamma }_1}(ϵ)A_{\mathrm{\Gamma }_2}(Eϵ)`$ (12)
$`\mathrm{cos}\left({\displaystyle \frac{1}{\mathrm{}}}S_{\mathrm{\Gamma }_1}(ϵ)\sigma _{\mathrm{\Gamma }_1}{\displaystyle \frac{\pi }{2}}\right)\mathrm{cos}\left({\displaystyle \frac{1}{\mathrm{}}}S_{\mathrm{\Gamma }_2}(Eϵ)\sigma _{\mathrm{\Gamma }_2}{\displaystyle \frac{\pi }{2}}\right).`$ (13)
To evaluate this asymptotically, we should include all critical points in the integration domain. Specifically, this integral has a stationary phase point within the integration domain and finite valued endpoints. We shall show that the stationary phase point corresponds to the situation where both particles are evolving dynamically with the energy partitioned between the two particles in a prescribed way. The endpoint contributions must be evaluated at energies such that one of the particles has all of the energy while the other has no energy. However, this contradicts our assumption that both particles are evolving — this is the definition of the dynamical term. Moreover, if we were to evaluate this contribution, the result would be meaningless since it involves using the trace formula at zero energy where it is known to fail. So we shall omit the contributions from the endpoints; this is discussed more fully in section IV D and in Appendix B, as well as in reference .
Hence, we evaluate the integral in (12) using only the stationary phase point. To leading order, we can extend the integration limits over an infinite domain. Writing the cosine functions as complex exponentials yields four integrals; the first is
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}dϵA_{\mathrm{\Gamma }_1}(ϵ)A_{\mathrm{\Gamma }_2}(Eϵ)\mathrm{exp}\left({\displaystyle \frac{i}{\mathrm{}}}\left(S_{\mathrm{\Gamma }_1}(ϵ)+S_{\mathrm{\Gamma }_2}(Eϵ)\right)\right)`$ (14)
$``$ $`A_{\mathrm{\Gamma }_1}(E_0)A_{\mathrm{\Gamma }_2}(EE_0)\sqrt{{\displaystyle \frac{2\pi \mathrm{}}{\left|\mathrm{{\rm Y}}(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2,E)\right|}}}`$ (16)
$`\mathrm{exp}\left({\displaystyle \frac{i}{\mathrm{}}}(S_{\mathrm{\Gamma }_1}(E_0)+S_{\mathrm{\Gamma }_2}(EE_0))+i\nu {\displaystyle \frac{\pi }{4}}\right)`$
where
$`\mathrm{{\rm Y}}(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2,E)`$ $`=`$ $`\left({\displaystyle \frac{^2S_{\mathrm{\Gamma }_1}(ϵ)}{ϵ^2}}+{\displaystyle \frac{^2S_{\mathrm{\Gamma }_2}(Eϵ)}{ϵ^2}}\right)|_{E_0}`$ (17)
$`\nu `$ $`=`$ $`\mathrm{sign}\left(\mathrm{{\rm Y}}(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2,E)\right).`$ (18)
$`E_0`$ is determined from the stationary phase condition
$`\left({\displaystyle \frac{S_{\mathrm{\Gamma }_1}(ϵ)}{ϵ}}+{\displaystyle \frac{S_{\mathrm{\Gamma }_2}(Eϵ)}{ϵ}}\right)|_{E_0}`$ $`=`$ $`0`$ (19)
$`T_{\mathrm{\Gamma }_1}(E_0)`$ $`=`$ $`T_{\mathrm{\Gamma }_2}(EE_0)`$ (20)
where we have used the fact that the derivative of the action with respect to energy is the period. $`E_0`$ is the energy of particle 1, $`EE_0`$ is the energy of particle 2 and $`E`$ is the total energy of the composite system. The saddle energy $`E_0`$ has a precise physical interpretation; Eq.(19) says that the energies of the two particles are partitioned so that the periods of both periodic orbits are the same. In other words, at $`E_0`$, we have orbits which are periodic in the full two particle phase space since after the period $`T`$ both particles return to their initial conditions.
The next integral has the same stationary phase condition as the first integral and is its complex conjugate. The third integral is
$$_{\mathrm{}}^{\mathrm{}}dϵA_{\mathrm{\Gamma }_1}(ϵ)A_{\mathrm{\Gamma }_2}(Eϵ)\mathrm{exp}\left(\frac{i}{\mathrm{}}\left(S_{\mathrm{\Gamma }_1}(ϵ)S_{\mathrm{\Gamma }_2}(Eϵ)\right)\right)$$
(21)
and has no stationary phase point since setting the first derivative of the action to zero yields the stationary phase condition
$$T_{\mathrm{\Gamma }_1}(E_0)=T_{\mathrm{\Gamma }_2}(EE_0).$$
(22)
The trace formula only involves orbits with positive period, so we ignore this possibility. The last integral is the complex conjugate of the third and will also be ignored.
Adding the contributions from the first two integrals, we arrive at the two particle trace formula:
$`\stackrel{~}{\rho }_1\stackrel{~}{\rho }_1(E){\displaystyle \frac{2}{(2\pi \mathrm{})^{3/2}}}{\displaystyle \underset{\mathrm{\Gamma }_1,\mathrm{\Gamma }_2}{}}{\displaystyle \frac{A_{\mathrm{\Gamma }_1}(E_0)A_{\mathrm{\Gamma }_2}(EE_0)}{\sqrt{\left|\mathrm{{\rm Y}}(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2,E)\right|}}}`$ (23)
$`\mathrm{cos}\left({\displaystyle \frac{1}{\mathrm{}}}(S_{\mathrm{\Gamma }_1}(E_0)+S_{\mathrm{\Gamma }_2}(EE_0))(\sigma _{\mathrm{\Gamma }_1}+\sigma _{\mathrm{\Gamma }_2}){\displaystyle \frac{\pi }{2}}+\nu {\displaystyle \frac{\pi }{4}}\right).`$ (24)
This result possesses the intuitive properties that, other than factors arising from the stationary phase analysis, the actions and Maslov indices are additive and the amplitudes are multiplicative. We note that this saddle-point analysis fails for the simplest problem in physics, the harmonic oscillator, where $`\mathrm{{\rm Y}}=0`$. This is because the two-particle harmonic oscillator has a higher degree of symmetry than we are accounting for here. This is a nongeneric property specific to the harmonic oscillator. We also stress that we have made no assumption about the stability or structure of the orbits. They can be isolated, stable or unstable or come in families. There are also problems with coexisting isolated orbits and families, such as those of the equilateral triangle billiard .
Note that the overall $`\mathrm{}`$ dependence is not multiplicative but picks up an additional factor of $`\mathrm{}^{1/2}`$ from the stationary phase integral. For isolated orbits, the amplitudes $`A`$ are independent of $`\mathrm{}`$ and the expression (23) has a $`1/\mathrm{}^{3/2}`$ prefactor as opposed to the $`1/\mathrm{}`$ in the amplitude of the single particle trace formula. The fact that the $`\mathrm{}`$ dependence is different implies that the periodic orbits of the full system come in continuous degenerate families rather than isolated trajectories, which in turn implies that there exists a continuous symmetry in the problem . This is an important point which we will address in a companion paper . (It was also noted in .) Nonetheless, it may be helpful to give a brief explanation here. Imagine the full phase space periodic orbit $`\mathrm{\Gamma }`$ consists of particle $`1`$ on a periodic orbit $`\mathrm{\Gamma }_1`$ with energy $`E_0`$ and particle $`2`$ on a distinct periodic orbit $`\mathrm{\Gamma }_2`$ with energy $`EE_0`$. We can define $`t=0`$ to be when particle $`2`$ is at some prescribed point on $`\mathrm{\Gamma }_2`$. Keeping particle $`2`$ fixed, we can change the position of particle $`1`$ on $`\mathrm{\Gamma }_1`$ to generate the initial condition of a distinct but congruent periodic orbit in the full phase space. Continuous time translation of the initial condition on $`\mathrm{\Gamma }_1`$ generates a continuous family of congruent periodic orbits in the full phase space. Since the time translational symmetry can be characterized by a single independent symmetry parameter, the $`\mathrm{}`$ dependence is $`𝒪\left(1/\sqrt{\mathrm{}}\right)`$ stronger than for a system with isolated periodic orbits .
## IV Two Particle Quantum Billiards
As an application of the formalism developed in section III, we consider the quantum billiard problem. Billiards are two dimensional enclosures that constrain the motion of a free particle. Classically, a particle has elastic collisions with the walls and depending on the geometric properties of the domain, the dynamics are either regular or chaotic. For the noninteracting problem, the two particles move independently of each other. In a billiard system, classical orbits possess simple scaling properties. For instance, the action of an orbit $`\mathrm{\Gamma }`$, $`S_\mathrm{\Gamma }(ϵ)=\sqrt{2mϵ}L_\mathrm{\Gamma }`$ and the period of the orbit is
$$T_\mathrm{\Gamma }(ϵ)=\frac{S_\mathrm{\Gamma }(ϵ)}{ϵ}=\frac{\sqrt{2m}L_\mathrm{\Gamma }}{2\sqrt{ϵ}}=\frac{\mathrm{}\sqrt{\alpha }}{2\sqrt{ϵ}}L_\mathrm{\Gamma }.$$
(25)
The parameter $`\alpha =2m/\mathrm{}^2`$ already appeared in Eq. (3); it will recur often. For example, in all final expressions, the energy occurs with $`\alpha `$; this is a result of the scaling property (the quantity $`\alpha E`$ having the units of $`1/\text{length}^2`$.) In the theoretical development, it will be convenient to retain $`\alpha `$ and use it to keep track of relative orders in the semiclassical expansions (since it contains $`\mathrm{}`$). However, once we have the final expressions, we are free to set it to unity for the purposes of numerical comparisons.
We also mention that for billiards, it is common to express the density of states in terms of the wave number $`k`$, where $`ϵ=k^2/\alpha `$ so that $`\rho (k)=2k\rho (ϵ)/\alpha `$. This is convenient since $`k`$ is conjugate to the periodic orbit lengths $`L`$. Therefore, many of our results will be quoted as a function of $`k`$, although it should be stressed that all convolution integrals must be done in the energy domain. Thus, we shall write
$$\rho _2^{\mathrm{sc}}(k)=\left(\overline{\rho }_1\overline{\rho }_1\right)(k)+2\left(\overline{\rho }_1\stackrel{~}{\rho }_1\right)(k)+\left(\stackrel{~}{\rho }_1\stackrel{~}{\rho }_1\right)(k).$$
(26)
Here, it is understood that each of the functions in brackets is first evaluated in the energy domain and then converted to the $`k`$ domain through the Jacobian relation above. This will always be the case when the argument is $`k`$, so that we will not always write brackets around the various functions. In terms of the wavenumber $`k`$, the decomposition (10) becomes
$$\rho _{S/A}(k)=\frac{1}{2}\left(\rho _2(k)\pm \frac{1}{\sqrt{2}}\rho _1\left(\frac{k}{\sqrt{2}}\right)\right).$$
(27)
### A Smooth Term
The smooth part is defined by the convolution integral
$$\overline{\rho }_1\overline{\rho }_1(E)=_0^Edϵ\overline{\rho }_1(ϵ)\overline{\rho }_1(Eϵ),$$
(28)
where $`\overline{\rho }_1`$ is given by the Weyl expansion. The expansion in (3) is taken only to order $`\mathrm{}^0`$. Hence, after expanding the integrand in (28), it is formally meaningless to include terms that are $`𝒪\left(1/\mathrm{}\right)`$ since there are corrections of the same order in $`\mathrm{}`$ that have not been calculated. Ignoring these terms and performing the necessary integrations, the smooth term is found to be
$$\overline{\rho }_1\overline{\rho }_1(E)\frac{\alpha ^2𝒜^2}{16\pi ^2}E\frac{\alpha ^{3/2}𝒜}{8\pi ^2}\sqrt{E}+\frac{\alpha ^2}{64\pi }+\frac{\alpha 𝒜𝒦}{2\pi }.$$
(29)
### B Cross Term
We next convolve $`\overline{\rho }_1`$ term by term with $`\stackrel{~}{\rho }_1`$. Asymptotically, each convolution integral receives contributions from the upper and lower endpoints. However, we shall only include one of these, namely the endpoint for which the trace formula is not evaluated at zero energy. As in section III, we neglect the other endpoint for reasons explained in IV D and Ref.. This is also discussed in Appendix B, where we evaluate the various integrals for the cross term exactly using isolated billiard orbits and show explicitly that an appropriate asymptotic expansion of the exact expression leads to consistent results.
After convolution, we find the area term involves the integral
$$\mathrm{Re}\left\{_0^EdϵA_\mathrm{\Gamma }(Eϵ)\mathrm{exp}\left(i\sqrt{\alpha (Eϵ)}L_\mathrm{\Gamma }i\sigma _\mathrm{\Gamma }\frac{\pi }{2}\right)\right\}.$$
(30)
The lower endpoint $`ϵ=0`$ corresponds to the physically meaningful situation while the upper endpoint is spurious in the sense mentioned above and discussed in detail below. Hence, to leading order, we can remove the amplitude factor from inside the integral, Taylor expand the argument of the exponential and extend the upper limit to infinity. This leads to
$$I_𝒜(E)\frac{\alpha 𝒜}{4\pi ^2}\underset{\mathrm{\Gamma }}{}\frac{A_\mathrm{\Gamma }}{T_\mathrm{\Gamma }}\mathrm{cos}\left(\sqrt{\alpha E}L_\mathrm{\Gamma }\sigma _\mathrm{\Gamma }\frac{\pi }{2}\frac{\pi }{2}\right).$$
(31)
By similar logic, the perimeter term and curvature terms are
$`I_{}(E)`$ $``$ $`{\displaystyle \frac{\sqrt{\alpha }}{8\pi ^{3/2}\sqrt{\mathrm{}}}}{\displaystyle \underset{\mathrm{\Gamma }}{}}{\displaystyle \frac{A_\mathrm{\Gamma }}{\sqrt{T_\mathrm{\Gamma }}}}\mathrm{cos}\left(\sqrt{\alpha E}L_\mathrm{\Gamma }\sigma _\mathrm{\Gamma }{\displaystyle \frac{\pi }{2}}{\displaystyle \frac{\pi }{4}}\right)`$ (32)
$`I_𝒦(E)`$ $``$ $`{\displaystyle \frac{𝒦}{\pi \mathrm{}}}{\displaystyle \underset{\mathrm{\Gamma }}{}}A_\mathrm{\Gamma }\mathrm{cos}\left(\sqrt{\alpha E}L_\mathrm{\Gamma }\sigma _\mathrm{\Gamma }{\displaystyle \frac{\pi }{2}}\right).`$ (33)
Note that all amplitudes and periods in (31) and (32) are evaluated at the system energy $`E`$. Recall $`\alpha 1/\mathrm{}^2`$ so that after convolution the sequence is an expansion in powers of $`\sqrt{\mathrm{}}`$ and not in powers of $`\mathrm{}`$ as for the original Weyl series (3). We also note that the first correction to $`I_𝒜`$ may be of the same order as $`I_𝒦`$ (as happens for the disk ) and should be included if this is the case. We then have $`\overline{\rho }\stackrel{~}{\rho }I_𝒜+I_{}+I_𝒦`$.
### C Dynamical Term
In this section, we derive a general expression for the dynamical term that is valid for any billiard problem. To this end, the first task is to determine the saddle energy from the stationary phase condition. Inserting (25) into (19) yields
$$\frac{L_{\mathrm{\Gamma }_1}}{\sqrt{E_0}}=\frac{L_{\mathrm{\Gamma }_2}}{\sqrt{EE_0}}$$
(34)
which implies
$$\frac{E_0}{E}=\frac{L_{\mathrm{\Gamma }_1}^2}{L_{\mathrm{\Gamma }_1}^2+L_{\mathrm{\Gamma }_2}^2},\frac{EE_0}{E}=\frac{L_{\mathrm{\Gamma }_2}^2}{L_{\mathrm{\Gamma }_1}^2+L_{\mathrm{\Gamma }_2}^2}$$
(35)
and
$$\mathrm{{\rm Y}}(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2,E)=\frac{\sqrt{2m}}{4E^{3/2}}\frac{\left(L_{\mathrm{\Gamma }_1}^2+L_{\mathrm{\Gamma }_2}^2\right)^{5/2}}{L_{\mathrm{\Gamma }_1}^2L_{\mathrm{\Gamma }_2}^2}.$$
(36)
Clearly $`\nu =1`$. We then substitute these results into (23) to obtain the two particle trace formula for billiards
$`\stackrel{~}{\rho }_1\stackrel{~}{\rho }_1(E)`$ $``$ $`{\displaystyle \frac{4E^{3/4}}{\sqrt{\mathrm{}}\alpha ^{1/4}(2\pi \mathrm{})^{3/2}}}`$ (39)
$`{\displaystyle \underset{\mathrm{\Gamma }_1,\mathrm{\Gamma }_2}{}}{\displaystyle \frac{L_{\mathrm{\Gamma }_1}L_{\mathrm{\Gamma }_2}}{\left(L_{\mathrm{\Gamma }_1}^2+L_{\mathrm{\Gamma }_2}^2\right)^{5/4}}}A_{\mathrm{\Gamma }_1}(E_0)A_{\mathrm{\Gamma }_2}(EE_0)`$
$`\mathrm{cos}\left(\sqrt{\alpha E}\sqrt{L_{\mathrm{\Gamma }_1}^2+L_{\mathrm{\Gamma }_2}^2}\left(\sigma _{\mathrm{\Gamma }_1}+\sigma _{\mathrm{\Gamma }_2}\right){\displaystyle \frac{\pi }{2}}{\displaystyle \frac{\pi }{4}}\right).`$
If the single particle periodic orbits are not isolated, then one must make direct use of the corresponding single particle amplitudes in (5) evaluated at the appropriate energies. We will show an explicit example of this when we analyze the disk billiard. Note the amplitudes $`A_\mathrm{\Gamma }`$ typically have an energy dependence so one cannot make any general statements about the energy dependence of this term except that the greater the dimensionality of the periodic orbit families, the greater the energy prefactor. For example, for the disk, it turns out to be $`E^{1/4}`$.
If the single particle periodic orbits are isolated, the amplitudes are given by (6), which for billiards is
$$A_\mathrm{\Gamma }(ϵ)=\frac{\sqrt{\alpha }\mathrm{}}{2\sqrt{ϵ}}\frac{L_\gamma }{\sqrt{\left|det\left(\stackrel{~}{M}_\mathrm{\Gamma }I\right)\right|}}.$$
(40)
In this case, the Gutzwiller amplitudes are evaluated at $`E_0`$ and $`EE_0`$, so we again make use of (35). After some algebra and simplification, we find
$`\stackrel{~}{\rho }_1\stackrel{~}{\rho }_1(E)`$ $``$ $`{\displaystyle \frac{\alpha ^{3/4}}{(2\pi )^{3/2}E^{1/4}}}`$ (43)
$`{\displaystyle \underset{\mathrm{\Gamma }_1,\mathrm{\Gamma }_2}{}}{\displaystyle \frac{L_{\gamma _1}L_{\gamma _2}(L_{\mathrm{\Gamma }_1}^2+L_{\mathrm{\Gamma }_2}^2)^{1/4}}{\sqrt{\left|det\left(\stackrel{~}{M}_{\mathrm{\Gamma }_1}I\right)\right|\left|det\left(\stackrel{~}{M}_{\mathrm{\Gamma }_2}I\right)\right|}}}`$
$`\mathrm{cos}\left(\sqrt{\alpha E}\sqrt{L_{\mathrm{\Gamma }_1}^2+L_{\mathrm{\Gamma }_2}^2}\left(\sigma _{\mathrm{\Gamma }_1}+\sigma _{\mathrm{\Gamma }_2}\right){\displaystyle \frac{\pi }{2}}{\displaystyle \frac{\pi }{4}}\right).`$
Note the $`E^{1/4}`$ prefactor which implies that the amplitude decays weakly with energy. This is the same prefactor that occurs in the single particle disk problem. This is not a coincidence, but arises from the fact that in both problems the periodic orbits come in one parameter families. Also, one must be careful to distinguish between $`L_\mathrm{\Gamma }`$, the length of a periodic orbit and $`L_\gamma `$, the length of the corresponding primitive periodic orbit. In general $`L_\mathrm{\Gamma }=n_\mathrm{\Gamma }L_\gamma `$ where $`n_\mathrm{\Gamma }`$ is the repetition index of that orbit.
### D Spurious Endpoint Contributions
As mentioned above, when confronted with convolution integrals, it is natural to analyse them asymptotically. This involves identifying the critical points and doing appropriate expansions in their neighbourhoods. In our work, these critical points are either stationary phase points or endpoints. The power of semiclassical methods is that each critical point can be given an immediate physical interpretation. For example, the stationary phase point in the dynamical term is found to be that energy such that the two particles have the same period so that the motion is periodic in the full two-particle phase space. This is intuitively reasonable. However, the same integral also has endpoints with finite valued contributions. We could do an asymptotic calculation in the vicinity of these points, but we can argue immediately that the result is spurious and not physically meaningful.
Recall the trace formulas are asymptotic in $`\mathrm{}`$ which typically also means asymptotic in energy. At the endpoints, one of the trace formulas is evaluated at small energy where it is known to be invalid. Alternatively, we can substitute for the trace formula any expression which is asymptotically equivalent to it and expect all meaningful results to be invariant to leading order. If we do this, we will find the endpoint contribution changes while the stationary phase contribution remains invariant, to leading order.
A further argument is that the structure of the endpoint contribution will be incorrect. Typically, it will be a sinusoid with an argument which does not depend on energy, but only depends on the properties of one of the orbits. Hence, it will have the same asymptotic structure as the cross term. However, we know that the cross term completely describes all such contributions and any further contribution with the same structure must be spurious.
Similarly, when we evaluate the cross term, we have two endpoint contributions. At one of these, we are evaluating the trace formula at some finite energy, which is reasonable. This endpoint corresponds to orbits which are periodic in the full phase space and in which one particle evolves on a single particle periodic orbit with all the energy, while the other remains fixed at some point in phase space with zero energy. At the other endpoint, we are evaluating the trace formula at zero energy, which is problematic. This corresponds to the contradictory situation in which the evolving particle has zero energy while the fixed particle has all the energy. In addition, upon inspection of this endpoint contribution, we find a function which is not oscillatory in energy and therefore has the same asymptotic structure as the smooth term. However, the smooth term already completely describes the average behaviour of the two particle density of states and any further contributions with the same structure must be spurious.
These situations are further examples of a general situation described in Ref. where it was shown that when integrating over the trace formula to obtain physical quantities, one should include all critical points except ones at which the trace formula is evaluated at zero energy. Such contributions should simply be ignored as spurious. In , the application was to thermodynamic calculations, but the principle is precisely the same. In Appendix B, we show the result of evaluating the cross term exactly for isolated orbits. An asymptotic analysis of this result leads to two terms which we can identify as coming from the two endpoints. One has the form used in this paper while the other is clearly spurious.
## V Two Particle Disk Billiard
In this section, we apply our results to the problem of two identical noninteracting particles moving in a two dimensional disk billiard of radius $`R`$. Quantum mechanically, this problem is a simple extension of the one body problem. Nevertheless, the spectrum has some interesting features which we discuss below.
### A Quantum Mechanics
For the disk billiard, a general two particle state can be written as
$$|m_1n_1,m_2n_2=|m_1n_1|m_2n_2$$
(44)
where the azimuthal quantum numbers $`m_1`$, $`m_2=0,\pm 1,\pm 2,\mathrm{}`$ and the radial quantum numbers $`n_1`$, $`n_2=1,2,3,\mathrm{}`$. We shall also use a more compact notation $`|N_1,N_2=|m_1n_1,m_2n_2`$ where $`N`$ denotes a pair of integers ($`m,n`$). We can immediately write down the wave numbers of the two particle system as
$$k_{N_1N_2}=\sqrt{\left(\frac{Z_{N_1}}{R}\right)^2+\left(\frac{Z_{N_2}}{R}\right)^2},$$
(45)
where $`Z_N`$ denotes the $`n`$th zero of the $`m`$th Bessel function $`J_m(z)`$. The set of all two particle states is given by $`\left\{|N_1,N_2\right\}`$.
The spectrum is highly degenerate. A typical state $`|N_1,N_2`$ is 8-fold degenerate since we can reverse the sign of either $`m_1`$ or $`m_2`$ or interchange the two particles and the resultant state has the same energy. However, if either $`m_1`$ or $`m_2`$ is zero or if $`N_1=N_2`$ then the state is 4-fold degenerate. If $`m_1=m_2=0`$ and $`N_1N_2`$, then the state is 2-fold degenerate whereas if $`m_1=m_2=0`$ and $`N_1=N_2`$, then the state is nondegenerate. If the particles are in distinct states, the degenerate multiplets divide evenly between the symmetric and antisymmetric spaces. However, if the particles are in the same state, $`N_1=N_2`$, it is somewhat less trivial. If $`N_1=N_2`$ and $`m_1=m_20`$, there is a 4-fold degenerate set of states: $`|mn,mn`$, $`|mn,mn`$, $`|mn,mn`$ and $`|mn,mn`$. The first two states belong to the symmetric space. From the second two states, we can construct one symmetric and one antisymmetric combination. (This is analogous to coupling two spin 1/2 states to construct a 3-fold symmetric $`S=1`$ state and a nondegenerate antisymmetric $`S=0`$ state.) If $`N_1=N_2`$ and $`m_1=m_2=0`$, this yields the state $`|0n,0n`$, which is singly degenerate and belongs to the symmetric space.
The quantum density of states
$$\rho _2(k)=\underset{N_1,N_2}{}\delta (kk_{N_1N_2})$$
(46)
and the corresponding symmetric and antisymmetric densities are shown in Fig. 1 as a function of the wavenumber $`k`$. Note that in this figure some of the peaks have different degeneracies in the symmetric and antisymmetric densities, as discussed above.
### B Semiclassical Density of States
We first review the semiclassical decomposition of the single particle density of states. The smooth part of the density of states may be obtained using the general result for two dimensional billiards (3). In fact, many higher order terms have been calculated . But, for our purposes, it suffices to use the first three terms as in (3) with $`𝒜=\pi R^2`$, $`=2\pi R`$ and $`𝒦=1/6`$.
The oscillating part of the level density can be obtained using trace formulas for systems with degenerate families of orbits. The periodic orbit families may be uniquely labelled by two integers ($`\mathrm{v}`$, $`\mathrm{w}`$) where $`\mathrm{v}`$ is the number of vertices and $`\mathrm{w}`$ is the winding number around the center. The two integers must satisfy the relation $`\mathrm{v}2\mathrm{w}`$. The length of an orbit with vertex number $`\mathrm{v}`$ and winding number $`\mathrm{w}`$ is given by $`L_{\mathrm{vw}}=2\mathrm{v}R\mathrm{sin}\left(\pi \mathrm{w}/\mathrm{v}\right)`$. With this notation, the trace formula for the oscillating part of the density of states is
$`\stackrel{~}{\rho }_1(ϵ)`$ $``$ $`{\displaystyle \frac{\alpha ^{3/4}}{2\sqrt{2\pi }ϵ^{1/4}}}{\displaystyle \underset{\mathrm{vw}}{}}{\displaystyle \frac{𝒟_{\mathrm{vw}}L_{\mathrm{vw}}^{3/2}}{\mathrm{v}^2}}`$ (48)
$`\mathrm{cos}\left(\sqrt{\alpha E}L_{\mathrm{vw}}3\mathrm{v}{\displaystyle \frac{\pi }{2}}+{\displaystyle \frac{\pi }{4}}\right)`$
where the sum goes from $`\mathrm{w}=1\mathrm{}\mathrm{}`$ and $`\mathrm{v}=2\mathrm{w}\mathrm{}\mathrm{}`$ and the degeneracy factor $`𝒟_{\mathrm{vw}}`$, which accounts for negative windings, is $`1`$ for $`\mathrm{v}=2\mathrm{w}`$ and $`2`$ for $`\mathrm{v}>2\mathrm{w}`$. Comparing (48) with the general form (5), we identify
$`A_{\mathrm{vw}}(ϵ)`$ $`=`$ $`{\displaystyle \frac{\sqrt{2\pi }\alpha ^{3/4}\mathrm{}𝒟_{\mathrm{vw}}L_{\mathrm{vw}}^{3/2}}{4\mathrm{v}^2ϵ^{1/4}}}`$ (49)
$`\sigma _{\mathrm{vw}}`$ $`=`$ $`3\mathrm{v}{\displaystyle \frac{1}{2}}`$ (50)
Adding the smooth and oscillating terms gives the semiclassical approximation to the single particle density of states which we denote by $`\rho _1^{\mathrm{sc}}(ϵ)`$.
To evaluate the semiclassical approximation to the two particle density of states, we must evaluate the smooth, cross and dynamical terms. The smooth term can be taken from Eq. (29) to be
$$\overline{\rho }_1\overline{\rho }_1(E)\frac{\alpha ^2R^4}{16}E\frac{\alpha ^{3/2}R^3}{4}\sqrt{E}+\left(\frac{3\pi +4}{48}\right)\alpha R^2.$$
(51)
The arguments of the previous section and in particular Eqs. (31) and (32) lead to the cross term
$`\overline{\rho }_1\stackrel{~}{\rho }_1(E){\displaystyle \frac{\alpha ^{5/4}R^2E^{1/4}}{4\sqrt{2\pi }}}{\displaystyle \underset{\mathrm{vw}}{}}{\displaystyle \frac{\sqrt{L_{\mathrm{vw}}}𝒟_{\mathrm{vw}}}{\mathrm{v}^2}}(\mathrm{cos}(\mathrm{\Phi }_{\mathrm{vw}}{\displaystyle \frac{\pi }{2}})`$ (52)
$`\sqrt{{\displaystyle \frac{\pi }{2}}}\chi _{\mathrm{vw}}\mathrm{cos}(\mathrm{\Phi }_{\mathrm{vw}}{\displaystyle \frac{\pi }{4}})+({\displaystyle \frac{1}{3}}+{\displaystyle \frac{R^2}{2L_{\mathrm{vw}}^2}})\chi _{\mathrm{vw}}^2\mathrm{cos}\mathrm{\Phi }_{\mathrm{vw}})`$ (53)
where $`\mathrm{\Phi }_{\mathrm{vw}}=\sqrt{\alpha E}L_{\mathrm{vw}}3\mathrm{v}\pi /2+\pi /4`$ and $`\chi _{\mathrm{vw}}=\sqrt{L_{\mathrm{vw}}}/(\alpha E)^{1/4}R`$. We have also included the first correction to the area term, $`I_𝒜(E)`$ which appears in the third term above . The dynamical term can be obtained using (39). Noting that $`\mathrm{\Gamma }_i`$ in (39) corresponds to the pair of integers ($`\mathrm{v}_i`$,$`\mathrm{w}_i`$), the result is
$`\stackrel{~}{\rho }_1\stackrel{~}{\rho }_1(E)`$ $``$ $`{\displaystyle \frac{\alpha ^{5/4}E^{1/4}}{4\sqrt{2\pi }}}{\displaystyle \underset{\mathrm{v}_1\mathrm{w}_1,\mathrm{v}_2\mathrm{w}_2}{}}\left({\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \frac{𝒟_iL_i^2}{\mathrm{v}_i^2}}\right)L_{12}^{3/2}`$ (55)
$`\mathrm{cos}\left(\sqrt{\alpha E}L_{12}3\left(\mathrm{v}_1+\mathrm{v}_2\right){\displaystyle \frac{\pi }{2}}+{\displaystyle \frac{\pi }{4}}\right)`$
where we defined $`L_i=L_{\mathrm{v}_i}\mathrm{w}_i`$, $`𝒟_i=𝒟_{\mathrm{v}_i}\mathrm{w}_i`$ and $`L_{12}=\sqrt{L_1^2+L_2^2}`$.
### C Numerics
For numerical purposes, we take $`\alpha =\mathrm{}=1`$ and $`R=1`$ so the single-particle energies are just the squares of the zeros of Bessel functions. Since we can only include a finite number of orbits, the periodic orbit sums must be truncated. As a representative case, we truncate the sum in (52) at $`\mathrm{w}_{\mathrm{max}}=50,\mathrm{v}_{\mathrm{max}}=100`$ (see Fig. 2) and use the same limits to truncate the quadruple sum in (55). This is a relatively small set of orbits, yet it does very well in reproducing the peaks of the quantum density of states. As an illustration, we show the first few peaks of (26) in Fig. 3. We calculated the semiclassical density of states (26) on the interval $`0k11`$. After doing so, we found only two sets of two peaks which were not resolved. These are shown in Fig. 4. Obviously, using more orbits will produce better results, but this increases the computation time because of the quadruple sum in (55). (Although one can reduce the computational overhead by limiting the sum to orbits whose amplitude exceeds some prescribed threshold ).
As an additional test, we want to determine whether (26) gives the correct degeneracies. We could do this by integrating the area under each of the peaks. However, a simpler procedure is to do a Gaussian smoothing by convolving $`\rho _2^{\mathrm{sc}}(k)`$ with an unnormalized Gaussian of variance $`\sigma `$:
$$\rho _2^{\mathrm{sc}}(k)G_\sigma (k)=_0^{\mathrm{}}dk^{}\rho _2^{\mathrm{sc}}(k^{})G_\sigma (kk^{})$$
(56)
where
$$G_\sigma (k)=\mathrm{exp}(k^2/2\sigma ^2).$$
(57)
and $`\sigma `$ is the smoothing width. The reason for this is that if the variance $`\sigma `$ of the Gaussian is larger than the intrinsic width of a peak in the semiclassical spectrum, then each peak acts like $`d\delta (kk_n)`$ with respect to the Gaussian. Thus, the integral in (56) becomes $`dG_\sigma (kk_n)`$ or $`d`$ at $`k=k_n`$. Of course, this is invalid when the spacing between two adjacent peaks is smaller than about $`\sigma `$. Some examples are discussed in the next section.
We also studied the symmetrised densities by using the expression (27) for both the quantum and semiclassical densities and convolving as above. The periodic orbit sums in the oscillating parts of the one and two body densities were truncated in the standard manner as before. The result of this numerical procedure is shown in Fig. 5. Clearly, the semiclassical approximations reproduce the correct degeneracies of the quantum spectrum as well as the approximate positions.
### D Discussion
In , it was noted that the trace formula replicates the single particle EBK spectrum obtained from torus quantisation more precisely than it duplicates the exact single particle quantum spectrum. After inspection of Figs. 3 and 5, we notice the same effect in the two particle spectrum. This property of the trace formulas also accounts for the unresolved peaks in the semiclassical spectrum. When the spacing of two levels of the EBK spectrum is very small, our truncated trace formulas may not resolve them, regardless of the spacing of the corresponding levels in the quantum spectrum (cf. Fig.4).
Comparing Figs. 1 and 5, we observe generally good agreement between the quantum and semiclassical spectra. Still, there are some apparent inconsistencies, for example, the two tall peaks in Fig. 5. These are the two sets of unresolved levels in Fig. 4, in each case an octet and a quartet. The reason for the discrepancy is the level spacings are smaller than the smoothing width $`\sigma `$, in contradiction to the assumption above, so that the peak height does not equal the degeneracy. In fact, the peak heights observed are rather close to 12 since the octets and quartets are very nearly degenerate on the scale of $`\sigma `$ and act almost like a 12-fold degenerate set. It is not perfectly $`12`$ due to the fact that the degeneracy is not perfect. However, we also observe that the integrated weight under the peak is consistent with a set of 12 energy levels. Other inconsistencies in Fig. 5 occur for the same reason. Of course, overall improvements can be made by including more orbits. As well, the artifacts of the single particle spectrum (some examples are marked by an “X” in Fig. 5) which arise from errors in the cross term, presumably decrease when corrections to the single particle trace formula are incorporated into the cross term. These preliminary numerical findings support our analytical results, which we now test in the rather different context of a chaotic billiard.
## VI Two Particle Cardioid Billiard
In this section, we study the problem of two identical noninteracting particles evolving in the cardioid billiard, which is fully chaotic . Since the billiard has a reflection symmetry, all the quantum states are either even or odd (this symmetry should not be confused with the symmetric/antisymmetric symmetry due to particle exchange.) In the subsequent analysis, we will exclusively use the odd spectrum. The reason for this is to avoid the additional complication of diffractive orbits which strike the vertex. Classically, these orbits are undefined and are therefore not included in the standard Gutzwiller theory. Studies of diffractive effects in trace formulas can be found in . The latter reference explores the specific application to the cardioid and shows that diffractive orbits are important in describing the even spectrum but are largely absent from the odd spectrum.
We could proceed as before by doing an explicit semiclassical analysis of each term in the decomposition of the two particle semiclassical density of states (11). However, we can simplify the analysis by removing single particle dynamics from the discussion. That is, we will focus exclusively on those quantum mechanical and semiclassical quantities that inherently describe two particle dynamics. More specifically, we compare the Fourier transform of the dynamical term
$$\stackrel{~}{F}_2^{\mathrm{sc}}(L)=\{\stackrel{~}{\rho }_1\stackrel{~}{\rho }_1(k)\}$$
(58)
with its quantum mechanical analogue which we define to be
$$\stackrel{~}{F}_2^{\mathrm{qm}}(L)=\{\rho _2(k)\overline{\rho }_1\overline{\rho }_1(k)2\overline{\rho }_1\stackrel{~}{\rho }_1(k)\}.$$
(59)
The integral operator $``$ will be defined precisely below. In the semiclassical transform (58), we use (43) expressed in terms of the wavenumber. Here, $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ are periodic orbits in the fundamental domain (i.e. the half-cardioid.) Orbit properties are discussed in and some representative orbits are shown in Fig. 6. The stability matrices in the denominators of the single particle Gutzwiller amplitudes are computed using the standard prescription for the stability of free flight billiards (see, for example, .)
In the quantum mechanical analogue (59), $`\rho _2(k)`$ is the quantum two particle density of states $`\rho _2(k)=_I\delta \left(kk_I\right)`$ where the superindex $`I`$ denotes the pair of integers ($`i,j`$) and $`k_I=\sqrt{k_i^2+k_j^2}`$. In (59), we subtract the smooth average part and the part which contains single particle dynamics. Using $`𝒜=3\pi /4`$, $`=6`$, and $`𝒦=3/16`$ in (29), the smooth term is
$$\overline{\rho }_1\overline{\rho }_1(E)\frac{9}{256}\alpha ^2E\frac{9}{16\pi }\alpha ^{3/2}\sqrt{E}+\left(\frac{9}{16\pi }+\frac{9}{128}\right)\alpha .$$
(60)
The cross term is given by the general expressions (31) and (32).
### A Numerics for the Unsymmetrised Cardioid
As before, we take $`\alpha =1`$ and use a standard sized cardioid as in to obtain the single particle spectrum. In this section, we numerically compare the two particle quantum mechanics with the two particle semiclassics. We do this by making a direct comparison of the Fourier transforms in the reciprocal space of orbit lengths, $`L`$. In this space, we expect peaks at lengths which correspond to the Euclidean lengths of the full periodic orbits of the two particle system. For instance, if the full orbit $`\mathrm{\Gamma }`$ is comprised of particle 1 travelling on the orbit $`\mathrm{\Gamma }_1`$ and particle 2 traversing a distinct orbit $`\mathrm{\Gamma }_2`$, we expect a peak at $`L_\mathrm{\Gamma }=\sqrt{L_{\mathrm{\Gamma }_1}^2+L_{\mathrm{\Gamma }_2}^2}`$. In the event that both particles are on the same orbit $`\mathrm{\Gamma }`$, we expect a peak at $`\sqrt{2}L_\mathrm{\Gamma }`$. In this way, any peak in the two particle spectrum can be attributed to the dynamics of a particular periodic orbit of the full classical phase space.
We construct the two particle spectrum by adding the energies of the single particle spectrum. We include the first 1250 single particle energies which allows us to construct the first $`766794`$ two particle energy levels representing all two particle energies less than $`6.8856\times 10^3`$.
For a precise numerical comparison, we define the Fourier transform
$$\{f(k)\}=_{\mathrm{}}^{\mathrm{}}dk\mathrm{w}(k)e^{ikL}f(k)$$
(61)
as a function of the conjugate variable $`L`$. Here, $`\mathrm{w}(k)`$ is the three term Blackman-Harris window function
$$\mathrm{w}(k)=\{\begin{array}{ccc}_{j=0}^2a_j\mathrm{cos}\left(2\pi j\frac{kk_0}{k_fk_0}\right)\hfill & & k_0<k<k_f\hfill \\ 0\hfill & & \text{otherwise}\hfill \end{array}$$
(62)
with $`(a_0,a_1,a_2)=(0.42323,0.49755,0.07922)`$. We choose $`k_0`$ and $`k_f`$ so that the window function goes smoothly to zero at the first and last eigenvalues of the two particle spectrum. Numerical integration of (58) and (59) using this definition of $``$ is displayed in Fig. 7. In the semiclassical transform, a total of 100 periodic orbits including multiple repetitions were used.
In Fig. 7, we observe good agreement between the quantum and semiclassical results for $`L<6.5`$ and $`L>10.3`$. In the region $`6.5<L<10.3`$, there are appreciable discrepancies for the following reason. Recall that the amplitudes of the two particle trace formula (43) apply only to billiard systems whose single particle periodic orbits are isolated. In the single particle cardioid problem, there exist orbits which are not well isolated in phase space, in fact two geometric orbits and a diffractive orbit are sometimes very close in phase space. For example, the two geometric orbits $`4a`$ and $`{}_{}{}^{}10b`$ together with the similar looking diffractive orbit $`4a^{}`$ (not shown) . In this event, the stationary phase approximation underlying the Gutzwiller formalism fails as does the argument in that diffractive orbits do not affect the odd spectrum. As a result, whenever a two particle orbit in the full space is comprised of one or both particles on one of these problematic single particle periodic orbits, the resulting two particle amplitude is inaccurate. (There is recent work on uniform approximations to account for such effects , unfortunately it seems not to apply to the cardioid which has the additional curious feature that the boundary curvature is infinite at the vertex.)
We now consider some specific examples. Consider first the peak structures “o” and “\*”. In this region, the single particle trace formula is erroneous and these errors propagate through to the cross term and inevitably to the quantum mechanical transform. We have also computed the cross term using quantum mechanics, that is, using $`\stackrel{~}{\rho }_1=\rho _1\overline{\rho }_1`$ in (59) and confirmed that these discrepancies do not arise (cf. Figs. 8 and 9). Thus, these discrepancies are due to errors in the semiclassical approximations of the cross term.
For the rest of the discussion, $`\mathrm{\Lambda }`$ refers to a particular periodic orbit family in the full phase space with each two particle orbit in this family comprised of the single particle orbits $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ and $`L_\mathrm{\Lambda }`$ are the lengths of the orbits in each family. Next, consider the peak structure at $`L7.5`$ (+). There are two families of orbits, $`\mathrm{\Lambda }_1`$ and $`\mathrm{\Lambda }_2`$ that are responsible for these peaks. The underlying structures of these orbits are shown in Table I. Bearing in mind the two single particle orbits $`\mathrm{\Gamma }_2`$ are not well isolated (cf. Fig. 6), the Gutzwiller amplitude of each $`\mathrm{\Gamma }_2`$ is incorrect. Consequently, the two particle Gutzwiller amplitude will also be incorrect, as Fig. 7 demonstrates. Let us look at the next peak structure. Clearly, the quantum peak heights are underestimated at $`L8.1`$ (X). We account for this by recognizing the two particle orbit structure involves the single particle orbit $`\frac{1}{2}(^{}8\mathrm{b})`$ which is an orbit which passes close to the vertex. More specifically, the orbit family $`\mathrm{\Lambda }=3`$ is composed of single particle orbits $`\mathrm{\Gamma }_1=\frac{1}{2}(^{}4\mathrm{b})`$ and $`\mathrm{\Gamma }_2=\frac{1}{2}(^{}8\mathrm{b})`$ and the lengths of these two particle orbits are $`L_\mathrm{\Lambda }=8.109`$. As a final illustration, we consider the region $`9.6<L<10.3`$. In this neighbourhood, the semiclassics are particularly bad. This can be accounted for by inspection Table II.
As the Table II and Fig. 6 show, there are many instances where both of the single particle orbits constituting the full orbit are poorly isolated. In view of this, both single particle Gutzwiller amplitudes are incorrect making the product even worse. This accounts for the gross inconsistencies in this region of the reciprocal space. The other discrepancies can be accounted for in a similar manner.
### B Symmetry Decomposition
In this section, we explore the symmetry decomposition of the two particle problem and in particular the comparison of the symmetrized two particle quantum mechanics with the corresponding two particle semiclassics. We start by defining the smooth and oscillating symmetry reduced densities of states from (27)
$$\overline{\rho }_{S/A}(k)=\frac{1}{2}\left((\overline{\rho }_1\overline{\rho }_1)(k)\pm \frac{1}{\sqrt{2}}\overline{\rho }_1\left(\frac{k}{\sqrt{2}}\right)\right)$$
(63)
and
$$\stackrel{~}{\rho }_{S/A}^{\mathrm{dyn}}(k)=\frac{1}{2}\left((\stackrel{~}{\rho }_1\stackrel{~}{\rho }_1)(k)\pm \frac{1}{\sqrt{2}}\stackrel{~}{\rho }_1\left(\frac{k}{\sqrt{2}}\right)\right).$$
(64)
While the second term in (64) is a single particle density, in a future paper we will demonstrate that this term describes the physical situation in which two particles are traversing the same periodic orbit, with the same energy and are exactly half a period out of phase. It describes the effect of particle exchange on the spectrum and for this reason affects the symmetric and antisymmetric spaces differently and is based on the theory of Ref. . Therefore, this second term also belongs to the two-particle dynamical term and we identify (64) as being a purely dynamical term. We want to compare it with the corresponding term in the symmetrized quantum densities of states. Hence, in analogy with the previous subsection, we compare
$$\stackrel{~}{F}_{S/A}^{\mathrm{dyn}}(L)=\left\{\stackrel{~}{\rho }_{S/A}^{\mathrm{dyn}}(k)\right\}$$
(65)
and
$$\stackrel{~}{F}_{S/A}^{\mathrm{qm}}(L)=\left\{\rho _{S/A}(k)\overline{\rho }_{S/A}(k)\overline{\rho }_1\stackrel{~}{\rho }_1(k)\right\}.$$
(66)
where $`\rho _{S/A}(k)`$ is the quantum bosonic (S) or fermionic (A) density of states.
### C Numerics for the Symmetrized Cardioid
In this section, we numerically compare the symmetrized quantum mechanics with the corresponding semiclassical quantities. In particular, we compute the transforms (65) and (66) . The symmetrized quantum densities are
$`\rho _S(k)`$ $`=`$ $`{\displaystyle \underset{i<j}{}}\delta \left(k\sqrt{k_i^2+k_j^2}\right)+{\displaystyle \underset{i}{}}\delta \left(k\sqrt{2}k_i\right),`$ (67)
$`\rho _A(k)`$ $`=`$ $`{\displaystyle \underset{i>j}{}}\delta \left(k\sqrt{k_i^2+k_j^2}\right)`$ (68)
using the same constraint on the energies as above. Of course, the sum of these symmetrized densities is the total density of states.
Before presenting our numerical results, we describe what we expect. First, all the peaks of the unsymmetrised two particle density should be present. In addition, for each periodic orbit $`\mathrm{\Gamma }`$, there should also be peaks at lengths $`L_\mathrm{\Gamma }/\sqrt{2}`$ arising from the oscillating part of the single particle density of states. The results are shown in Fig. 10. For the two particle density term of (65), we used the same 100 two particle orbits of section VI A while in the single density term we used all single particle orbits with length $`L<11`$. As well, we included the single particle orbit $`\frac{1}{2}(^{}10\mathrm{h})`$ (not shown in Fig. 6) which has a length $`L=10.477`$. Fig.10 displays the peak structure in the reciprocal space up to $`L=6.75`$.
We notice that most of the amplitude divides evenly between the symmetric and the antisymmetric densities. Nonetheless, there are exceptions such as the peak at $`L3.6`$ (o). Here, both terms of (64) contribute and the difference in the sign of the second term accounts for the uneven amplitude division between the two symmetrized densities. Semiclassically, we account for the peak structure by noting that two different physical situations are responsible for the peak structure “o”. First, there is the situation in which both particles are on the orbit $`\mathrm{\Gamma }=\frac{1}{2}(^{}2\mathrm{a})`$ with no restrictions on the time phase difference between the two particles. This contribution comes from the two particle density term resulting in a peak at a length $`\sqrt{2}L_\mathrm{\Gamma }=3.673`$ and produces identical structures in both densities. The second situation occurs when both particles are on the orbit $`\mathrm{\Gamma }=(^{}2\mathrm{a})`$ exactly half a period out of phase. This contribution comes from the single density term at $`L_\mathrm{\Gamma }/\sqrt{2}=3.673`$ and is explained more fully in . Since the second contribution comes with a different sign in the two symmetries, the amplitudes are different for the symmetric and antisymmetric spaces. In this particular case, it is stronger in the symmetric density and weaker in the antisymmetric density, although the opposite may be true in other cases.
In closing, we remark that the overall agreement between the quantum and semiclassical calculations is good. The poorly reproduced peak just above $`L=5`$ (+) comes from the single density term. This is just the poorly reproduced peak of the single particle density at $`L7`$ shifted down by a factor of $`\sqrt{2}`$.
## VII Conclusion
Initially, we developed a semiclassical formalism to describe the two particle density of states. After deriving a trace formula describing two particle dynamics, we investigated its structure and noted intuitive properties such as the additivity of the actions and topological phase factors. As well, we briefly explained the structure of the full two particle orbits which come in degenerate families. As a first application, we wrote down a two particle trace formula for two identical particles in a billiard. The semiclassical symmetry decomposition involved formal substitution of the semiclassical quantities into the quantum mechanical expressions for the symmetrized densities. In a future paper , we show how these formal expressions emerge directly from the classical structures.
Following these general considerations, we studied two identical noninteracting particles in a disk and in a cardioid. In each case, we find that the formalism correctly reproduced the full and symmetrized densities of states. In the integrable problem, we found that our formalism replicates the two body EBK spectrum more precisely than the quantum spectrum, suggesting a deep connection between periodic orbit theory and EBK quantization for integrable systems. In the chaotic cardioid billiard, we note that the single particle orbits which pass close to the vertex lead to inconsistencies in the Fourier transform of the semiclassical density of states. Clearly, our formalism fails here because the Gutzwiller theory itself fails for these “semi-diffractive” orbits. For all other orbits, the two particle trace formula works very well.
The techniques employed here involve the classical phase space of each particle. In a future paper , we derive the same results by working in the full two particle phase space. This approach has the advantage of being more general than what we have presented here. Nonetheless, it is conceptually useful to see how the same structure emerges from these two distinct points of view. We would also like to incorporate interactions between the particles. Such a project would undoubtedly require working in the full phase space since it is no longer true that the full density of states is the convolution of the single particle level densities. This provides an additional motivation for working out the noninteracting problem in the full phase space as a first step towards the more ambitious goal. This full phase space analysis also generalises more readily to more particles. Finally, it has the conceptual advantage that the spurious endpoint contributions discussed in IV D and Appendix B do not arise and therefore need not be explained away.
It may be argued that interacting many body systems are too complex to be accessible to the semiclassical method. However, given the intractability of the many body problem, there may well be questions which semiclassical theory can answer. In particular, we have in mind the applications of semiclassical theory to mesoscopic physics . Here, our seemingly academic study of billiard systems finds physical applications in the context of nanostructures. For example, the disk billiard can serve as a realistic lowest-order approximation to the mean field of the electrons in a circular quantum dot . In fact, many phenomena in ballistic mesoscopic systems can, at least qualitatively, be described by using quantum billiards with independent particles as physical models.
###### Acknowledgements.
We thank Rajat Bhaduri, Matthias Brack and Randy Dumont for useful discussions.
## VIII Appendix A: Nonidentical Particles
As we have mentioned, most of the discussion still applies if the two particles are not identical. Another situation is a single particle in a separable potential. For example, in two dimensions, one could have $`V(x,y)=V_a(x)+V_b(y)`$ in which case, the dynamics in the $`x`$ direction are completely uncoupled from the dynamics in the $`y`$ direction so that the system is formally the same as if there were distinct particles executing the $`x`$ and $`y`$ motions. The formalism presented above follows in a natural way. The main differences are that one no longer considers the symmetrised density of states since the symmetry of particle exchange no longer exists and secondly there are two distinct cross terms so that (11) is replaced by
$`\rho _2(E)`$ $`=`$ $`\overline{\rho }_{1a}\overline{\rho }_{1b}(E)+\overline{\rho }_{1a}\stackrel{~}{\rho }_{1b}(E)`$ (69)
$`+`$ $`\stackrel{~}{\rho }_{1a}\overline{\rho }_{1b}(E)+\stackrel{~}{\rho }_{1a}\stackrel{~}{\rho }_{1b}(E),`$ (70)
where the indices $`a`$ and $`b`$ refer to the two distinct particles, while the indices $`1`$ and $`2`$ still refer to one or two particle densities of states.
Imagine, for example, that we have two nonidentical particles in distinct billiard enclosures. We introduce two parameters, $`\alpha _a=2m_a/\mathrm{}^2`$ and $`\alpha _b=2m_b/\mathrm{}^2`$. The smooth term (29) is replaced by
$`\overline{\rho }_{1a}\overline{\rho }_{1b}(E)`$ $``$ $`{\displaystyle \frac{\alpha _a\alpha _b𝒜^2}{16\pi ^2}}E`$ (71)
$``$ $`\left(\alpha _a^{1/2}+\alpha _b^{1/2}\right){\displaystyle \frac{\sqrt{\alpha _a\alpha _b}}{16\pi ^2}}𝒜\sqrt{E}`$ (72)
$`+`$ $`{\displaystyle \frac{\alpha _a^{1/2}\alpha _b^{1/2}^2}{64\pi }}+{\displaystyle \frac{(\alpha _a+\alpha _b)𝒜𝒦}{4\pi }}.`$ (73)
The cross terms each separately have the same structure as the cross term for identical particles. Obviously, they are no longer equal to each other, but functionally little has changed. It is just a question of inserting the relevant information from the different smooth and oscillating densities of states of the two particles. Following the same logic as before, we find
$`I_𝒜(E)`$ $``$ $`{\displaystyle \frac{\alpha _a𝒜_a}{4\pi ^2}}{\displaystyle \underset{\mathrm{\Gamma }_b}{}}{\displaystyle \frac{A_{\mathrm{\Gamma }}^{}{}_{b}{}^{}}{T_{\mathrm{\Gamma }}^{}{}_{b}{}^{}}}\mathrm{cos}\left(\mathrm{\Phi }_{\mathrm{\Gamma }_b}{\displaystyle \frac{\pi }{2}}\right),`$ (74)
$`I_{}(E)`$ $``$ $`{\displaystyle \frac{\sqrt{\alpha _a}_a}{8\pi ^{3/2}\sqrt{\mathrm{}}}}{\displaystyle \underset{\mathrm{\Gamma }_b}{}}{\displaystyle \frac{A_{\mathrm{\Gamma }}^{}{}_{b}{}^{}}{\sqrt{T_\mathrm{\Gamma }}_b}}\mathrm{cos}\left(\mathrm{\Phi }_{\mathrm{\Gamma }_b}{\displaystyle \frac{\pi }{4}}\right),`$ (75)
$`I_𝒦(E)`$ $``$ $`{\displaystyle \frac{𝒦_a}{\pi \mathrm{}}}{\displaystyle \underset{\mathrm{\Gamma }_b}{}}A_{\mathrm{\Gamma }}^{}{}_{b}{}^{}\mathrm{cos}\left(\mathrm{\Phi }_{\mathrm{\Gamma }_b}\right),`$ (76)
where $`\mathrm{\Phi }_{\mathrm{\Gamma }_b}=\sqrt{\alpha _bE}L_{\mathrm{\Gamma }_b}\sigma _{\mathrm{\Gamma }_b}\pi /2`$. For $`\overline{\rho }_{1b}\stackrel{~}{\rho }_{1a}(E)`$, we just interchange $`a`$ and $`b`$.
The formula for $`\stackrel{~}{\rho }_{1a}\stackrel{~}{\rho }_{1b}(E)`$ still has the same basic structure, but should obviously use the distinct periodic orbits for particles $`a`$ and $`b`$. In particular, Eqs. (12) and (23) still apply, but with two important differences. Firstly, the double sums over periodic orbits are now labelled by the distinct periodic orbits of the two particles. Secondly, the energy partition will change due to differing masses. The criterion of stationary phase will still specify that the two particles have the same period, but relations such as (34) and (35) do not apply since they assume equal masses. The generalisations are rather straight-forward to determine. For example, the saddle energies (35) are replaced by
$`{\displaystyle \frac{E_0}{E}}={\displaystyle \frac{m_aL_{\mathrm{\Gamma }_a}^2}{m_aL_{\mathrm{\Gamma }_a}^2+m_bL_{\mathrm{\Gamma }_b}^2}},{\displaystyle \frac{EE_0}{E}}={\displaystyle \frac{m_bL_{\mathrm{\Gamma }_b}^2}{m_aL_{\mathrm{\Gamma }_a}^2+m_bL_{\mathrm{\Gamma }_b}^2}}`$ (77)
while the general dynamical expression for billiards (39) is replaced by
$`\stackrel{~}{\rho }_{1a}\stackrel{~}{\rho }_{1b}(E){\displaystyle \frac{(2E)^{3/4}\sqrt{\alpha _a\alpha _b\mathrm{}}}{(2\pi )^{3/2}}}`$ (78)
$`{\displaystyle \underset{\mathrm{\Gamma }_a,\mathrm{\Gamma }_b}{}}{\displaystyle \frac{L_{\mathrm{\Gamma }_a}L_{\mathrm{\Gamma }_b}}{\left(m_aL_{\mathrm{\Gamma }_a}^2+m_bL_{\mathrm{\Gamma }_b}^2\right)^{5/4}}}A_{\mathrm{\Gamma }_a}(E_0)A_{\mathrm{\Gamma }_b}(EE_0)`$ (79)
$`\mathrm{cos}\left(\sqrt{\alpha _aL_{\mathrm{\Gamma }_a}^2+\alpha _bL_{\mathrm{\Gamma }_b}^2}\sqrt{E}\left(\sigma _{\mathrm{\Gamma }_a}+\sigma _{\mathrm{\Gamma }_b}\right){\displaystyle \frac{\pi }{2}}{\displaystyle \frac{\pi }{4}}\right).`$ (80)
In the special case of identical particles, it is simple to check that this expression reduces to (39). For lack of an immediate physical context, we do not explore this case any further.
## IX Appendix B: Spurious End-point Contributions for the Cardioid
Here we evaluate the cross term integrals exactly for isolated periodic orbits. This allows us to do an asymptotic expansion to explicitly demonstrate that the additional endpoint contributions not included are spurious. We must evaluate the integral
$$\overline{\rho }_1\stackrel{~}{\rho }_1(E)=_0^Edϵ\overline{\rho }_1(ϵ)\stackrel{~}{\rho }_1(Eϵ)$$
(81)
where $`\overline{\rho }_1(ϵ)`$ is given by the Weyl expansion (3) and $`\stackrel{~}{\rho }_1(ϵ)`$ for a billiard with isolated orbits is given by
$$\stackrel{~}{\rho }_1(ϵ)\frac{\alpha ^{1/2}}{2\pi \sqrt{ϵ}}\underset{\mathrm{\Gamma }}{}\frac{L_\gamma }{\sqrt{\left|det\left(\stackrel{~}{M}_\mathrm{\Gamma }I\right)\right|}}\mathrm{cos}\left(\sqrt{\alpha ϵ}L_\mathrm{\Gamma }\sigma _\mathrm{\Gamma }\frac{\pi }{2}\right).$$
(82)
This gives
$`\overline{\rho }_1\stackrel{~}{\rho }_1(E)`$ $``$ $`{\displaystyle \underset{\mathrm{\Gamma }}{}}{\displaystyle \frac{L_\gamma }{\sqrt{\left|det\left(\stackrel{~}{M}_\mathrm{\Gamma }I\right)\right|}}}`$ (84)
$`\left(\alpha ^{3/2}{\displaystyle \frac{𝒜}{8\pi ^2}}I_1\alpha {\displaystyle \frac{}{16\pi ^2}}I_2+\alpha ^{1/2}{\displaystyle \frac{𝒦}{2\pi }}I_3\right)`$
where
$`I_1`$ $`=`$ $`{\displaystyle _0^E}dϵ{\displaystyle \frac{1}{\sqrt{Eϵ}}}\mathrm{cos}\left(\sqrt{\alpha (Eϵ)}L_\mathrm{\Gamma }\sigma _\mathrm{\Gamma }{\displaystyle \frac{\pi }{2}}\right),`$ (85)
$`I_2`$ $`=`$ $`{\displaystyle _0^E}dϵ{\displaystyle \frac{1}{\sqrt{ϵ}}}{\displaystyle \frac{1}{\sqrt{Eϵ}}}\mathrm{cos}\left(\sqrt{\alpha (Eϵ)}L_\mathrm{\Gamma }\sigma _\mathrm{\Gamma }{\displaystyle \frac{\pi }{2}}\right),`$ (86)
$`I_3`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{E}}}\mathrm{cos}\left(\sqrt{\alpha E}L_\mathrm{\Gamma }\sigma _\mathrm{\Gamma }{\displaystyle \frac{\pi }{2}}\right).`$ (87)
If we evaluate the first two integrals exactly, we get
$$I_1=\frac{2}{\alpha ^{1/2}L_\mathrm{\Gamma }}\left(\mathrm{cos}\left(\mathrm{\Phi }_\mathrm{\Gamma }+\varphi _\mathrm{\Gamma }\frac{\pi }{2}\right)\mathrm{cos}\left(\varphi _\mathrm{\Gamma }\frac{\pi }{2}\right)\right)$$
(88)
and
$`I_2`$ $`=`$ $`\pi \mathrm{cos}\varphi _\mathrm{\Gamma }J_0\left(\mathrm{\Phi }_\mathrm{\Gamma }\right)\pi \mathrm{sin}\varphi _\mathrm{\Gamma }𝐇_0\left(\mathrm{\Phi }_\mathrm{\Gamma }\right)`$ (89)
$``$ $`\sqrt{{\displaystyle \frac{2\pi }{\sqrt{\alpha E}L_\mathrm{\Gamma }}}}\mathrm{cos}\left(\mathrm{\Phi }_\mathrm{\Gamma }+\varphi _\mathrm{\Gamma }{\displaystyle \frac{\pi }{4}}\right)`$ (91)
$`{\displaystyle \frac{2}{\sqrt{\alpha E}L_\mathrm{\Gamma }}}\mathrm{cos}\left(\varphi _\mathrm{\Gamma }\right)+\mathrm{}`$
where $`\mathrm{\Phi }_\mathrm{\Gamma }=\sqrt{\alpha E}L_\mathrm{\Gamma }`$, $`\varphi _\mathrm{\Gamma }=\sigma _\mathrm{\Gamma }\pi /2`$, $`J_0`$ is a zero-order Bessel function and $`𝐇_0`$ is a zero-order Struve function. In the second line of (89), we have used the asymptotic expansions of these two functions.
In both $`I_1`$ and $`I_2`$, we note that asymptotically there are terms with two distinct structures. The first are terms which are sinusoidal in $`\sqrt{E}`$ and correspond exactly to what was used as the cross term for the cardioid (i.e. Eqs. (31) and 32)). There are also terms which are nonsinusoidal in $`E`$. In $`I_1`$, this comes directly from the upper endpoint of the integral while in $`I_2`$ it comes from the expansion of the Struve function. In each term, the nonsinusoidal terms arise from the endpoint around $`ϵ=E`$ which, as we argued in section IV D, corresponds to an unphysical situation. Therefore, keeping only the asymptotically appropriate term (i.e. the oscillatory one) yields the correct behaviour for the cross term.
A similar analysis would yield similar results for the spurious endpoint contributions in the cross term of the disk billiard and the dynamical term of either billiard. |
warning/0001/math0001024.html | ar5iv | text | # HyperKähler Potentials in Cohomogeneity Two
## 1. Introduction
HyperKähler metrics are special Ricci-flat structures that are known to arise in many physical theories. For example, moduli spaces of magnetic monopoles often carry such metrics. For good choices of boundary conditions, these moduli spaces can be identified with more familiar mathematical objects. In this way, hyperKähler metrics have been shown to exist on the adjoint orbits of a complex semi-simple Lie group $`G^{}`$ . In , it was shown that these examples include all hyperKähler metrics of cohomogeneity one.
Some of the earliest examples of hyperKähler metrics were found by Calabi . His method was to take a complex symplectic manifold, such as the cotangent bundle $`T^{}\mathrm{P}(n)`$, and find a potential for a Kähler structure that would combine with the complex symplectic structure to give a hyperKähler metric. This approach has been applied to certain semi-simple nilpotent orbits by a number of authors. Biquard & Gauduchon gave a beautiful construction for a potential on those semi-simple orbits that are the cotangent bundle of a Hermitian symmetric space. At the other extreme, Hitchin used spectral theory to describe a potential for the biggest semi-simple orbit in $`𝔰𝔩(n,)`$ in terms of theta functions (the special case of $`n=2`$ may be found in ).
Much attention has been paid to the semi-simple orbits, because one can show that they are the only orbits to admit hyperKähler metrics that are complete. However, the incomplete metrics on nilpotent orbits still have much interest. One reason, is that each such orbit admits a *hyperKähler potential*, a function that is a Kähler potential for each complex structure compatible with the hyperKähler structure, and so these metrics on nilpotent orbits determine quaternionic Kähler metrics of positive scalar curvature on a certain quotient manifold .
The structures considered on coadjoint orbits are invariant under the action of the compact group $`G`$. For nilpotent orbits, there is a natural partial order given by inclusions of closures. When $`G`$ is simple, the smallest non-trivial orbits in this order are unique and they are distinguished by being of cohomogeneity one under the action of $`G`$. In , it was shown that the nilpotent orbits of cohomogeneity two also fit nicely in to the partial order: except when $`G=\text{SU}(3)`$, they are exactly the next-to-minimal orbits. Given that the nilpotent orbits of cohomogeneity one are understood (see also ), it is natural to look at those of cohomogeneity two.
In this paper, we consider cohomogeneity-two nilpotent orbits and find all compatible $`G`$-invariant hyperKähler metrics on them that admit Kähler potentials. Our approach is that of Calabi’s and we obtain the hyperKähler potentials explicitly. The hyperKähler potentials are unique, but in a few cases we find that they lie in a one-parameter family of hyperKähler metrics with Kähler potential. These families may be regarded as generalisations of the Eguchi-Hanson metrics in dimension four.
Combining our results with , means that hyperKähler potentials are now known for all next-to-minimal orbits. One feature of the cohomogeneity-two case that makes the calculations possible, is that each element of the orbit lies in a small rank $`2`$ real subalgebra which determines much of the hyperKähler structure. In fact, unless $`G`$ is the exceptional Lie group $`\text{G}_2`$, that subalgebra is $`𝔰𝔬(4,)=𝔰𝔩(2,)𝔰𝔩(2,)`$ and the geometry is the product of the structures from each factor.
For some cohomogeneity-two orbits the hyperKähler potential may also be obtained by one of three other methods: a hyperKähler quotient construction, a finite-cover by a minimal orbit for another group, or a limit of a family of semi-simple orbits. The first two methods will be described elsewhere; the first only succeeds if the hyperKähler quotient is sufficiently simple and the second only covers orbits on the list of “shared orbits” of Brylinski & Kostant . The third is contained in Biquard & Gauduchon’s work . However, there are orbits for which the approach of this paper is the only one known to give the result and our approach is uniform for all orbits of cohomogeneity two.
###### Acknowledgements.
It is a pleasure to thank Brian Dupée for advice on Maple, Alastair King and Francis Burstall for many useful conversations and Claude LeBrun for enlightenment. This research was supported by the epsrc of Great Britain. The first named author is also grateful for partial support from the KBN of Poland.
## 2. Preliminaries
### 2.1. HyperKähler Structures
Let $`M`$ be a manifold with endomorphisms $`I`$, $`J`$ and $`K`$ of the tangent bundle $`TM`$ satisfying the quaternion identities
$$I^2=J^2=1\text{and}IJ=K=JI.$$
This gives $`T_xM`$ the structure of an $``$-module and so implies that the dimension of $`M`$ is a multiple of $`4`$. If $`g`$ is a Riemannian metric on $`M`$ preserved by $`I`$, $`J`$ and $`K`$, in the sense that $`g(IX,IY)=g(X,Y)`$, etc., for all tangent vectors $`X,Y`$, then we can define two-forms $`\omega _I`$, $`\omega _J`$ and $`\omega _K`$ by
$$\omega _I(X,Y)=g(X,IY),\text{etc.}$$
If these three two-forms are closed, the structure $`(M,g,I,J,K)`$ is said to be *hyperKähler*.
Hitchin showed that on a hyperKähler manifold, the almost complex structures $`I`$, $`J`$ and $`K`$ are integrable, and thus $`(M,g)`$ is a Kähler manifold in three distinct ways. The restricted holonomy group $`\mathrm{Hol}_g`$ of $`(M,g)`$ is then contained in $`\text{Sp}(n)`$. As $`\text{Sp}(n)`$ is a subgroup of $`\text{SU}(2n)`$, this implies that any hyperKähler metric $`g`$ is Ricci-flat.
A function $`\rho :M`$ is a Kähler potential for the complex structure $`I`$ if $`\omega _I=i_I\overline{_I}\rho `$. This may be reformulated as
(2.1)
$$\begin{array}{cc}\hfill \omega _I& =i_I\overline{_I}\rho =id\overline{_I}\rho =\frac{i}{2}d(diId)\rho \hfill \\ & =\frac{1}{2}dId\rho .\hfill \end{array}$$
The function $`\rho `$ is a *hyperKähler potential* if it is simultaneously a Kähler potential for $`I`$, $`J`$ and $`K`$. HyperKähler potentials are defined up to an additive constant. The existence of a hyperKähler potential implies strong restrictions on the geometry of $`M`$ : the metric $`g`$ and potential $`\rho `$ satisfy $`^2\rho =g`$; the manifold $`M`$ admits an infinitesimal action of $`^{}`$, with $`\text{Sp}(1)^{}`$ preserving $`g`$ and permuting $`I`$, $`J`$ and $`K`$; the $`^{}`$-orbits are flat and totally geodesic; locally $`M`$ fibres over a quaternionic Kähler orbifold of positive scalar curvature.
We will be considering hyperKähler structures that are invariant under the action of a compact group $`G`$. It is therefore worth noting that if we have a Kähler potential then this may be taken to be $`G`$-invariant. Indeed, if $`\rho `$ is any Kähler potential, then since the $`G`$-action preserves $`I`$, the expression $`_I\overline{_I}\rho `$ is equivariant for the action of $`G`$. However, $`\omega _I=i_I\overline{_I}\rho `$, is assumed to be $`G`$-invariant, so averaging $`\rho `$ over the group action produces an invariant Kähler potential.
### 2.2. Lie Algebras and Orbits
On the semi-simple complex Lie algebra $`𝔤^{}`$, let $`,=,_𝔤`$ be the negative of the Killing form and let $`\sigma `$ be a real structure giving a compact real form $`𝔤`$ of $`𝔤^{}`$.
At a point $`X`$ of a nilpotent orbit $`𝒪`$, the vector field generated by $`A`$ in $`𝔤^{}`$ is $`\xi _A=[A,X]`$. These vector fields satisfy $`[\xi _A,\xi _B]=\xi _{[A,B]}`$, for all $`A,B𝔤^{}`$.
The orbit $`𝒪`$ carries a complex structure $`I`$ defined by
(2.2)
$$I\xi _A=i\xi _A=\xi _{iA}.$$
There is also a complex symplectic form, known as the Kirillov-Kostant-Souriau form, on $`𝒪`$ which we take to be given by
(2.3)
$$\omega _c^𝒪(\xi _A,\xi _B)_X=X,[A,B]=\xi _A,B.$$
We will be looking for hyperKähler structures on $`𝒪`$ with $`I`$ given by (2.2) and $`\omega _J+i\omega _K=\omega _c^𝒪`$. We will call these *compatible hyperKähler structures* on $`𝒪`$.
## 3. Potentials Depending on Two Invariants
Consider the following two functions on a nilpotent orbit $`𝒪`$:
$$\eta _1(X)=X,\sigma X\text{and}\eta _2(X)=[X,\sigma X],[X,\sigma X].$$
Note that $`\eta _2(X)=Y,\sigma Y`$ with $`Y=[X,\sigma X]`$, so is positive, and that both $`\eta _1`$ and $`\eta _2`$ are invariant under the action of the compact group $`G`$. Suppose $`\rho `$ is a Kähler potential for $`I`$ depending only on $`\eta _1`$ and $`\eta _2`$, i.e.,
(3.1)
$$\rho =\rho (\eta _1,\eta _2).$$
###### Lemma 3.1.
At $`X𝒪`$, the two-form $`\omega _I`$ defined by $`\rho `$ in formula (2.1) is
(3.2)
$$\begin{array}{cc}\hfill \omega _I(\xi _A,\xi _B)_X& =2\rho _1\mathrm{Im}\xi _A,\sigma \xi _B\hfill \\ & 4\rho _2\mathrm{Im}\xi _A,[\sigma \xi _B,[X,\sigma X]]+[\sigma X,[X,\sigma \xi _B]]\hfill \\ & +2\rho _{11}\mathrm{Im}\left(\xi _A,\sigma X\sigma \xi _B,X\right)\hfill \\ & 4\rho _{12}\begin{array}{cc}\hfill \mathrm{Im}(& \xi _A,[\sigma X,[X,\sigma X]]\sigma \xi _B,X\hfill \\ & +\xi _A,\sigma X\sigma \xi _B,[X,[\sigma X,X]])\hfill \end{array}\hfill \\ & +8\rho _{22}\mathrm{Im}\left(\xi _A,[\sigma X,[X,\sigma X]]\sigma \xi _B,[X,[\sigma X,X]]\right),\hfill \end{array}$$
where $`\rho _i=\rho /\eta _i`$, etc.
###### Proof.
Expanding (2.1), we have
(3.3)
$$\begin{array}{c}2\omega _I=\rho _1dId\eta _1+\rho _2dId\eta _2+\rho _{11}d\eta _1Id\eta _1\hfill \\ \hfill +\rho _{12}(d\eta _2Id\eta _1+d\eta _1Id\eta _2)+\rho _{22}d\eta _2Id\eta _2.\end{array}$$
The exterior derivative of $`\eta _1`$ is given by
$$d\eta _1(\xi _A)_X=[A,X],\sigma X+X,\sigma [A,X]=2\mathrm{Re}\xi _A,\sigma X.$$
Hence $`Id\eta _1(\xi _A)=2\mathrm{Im}\xi _A,\sigma X`$ and $`dId\eta (\xi _A,\xi _B)=4\mathrm{Im}\xi _A,\sigma \xi _B`$, at $`X𝒪`$.
For $`\eta _2`$, the initial computation is similar and gives
$$d\eta _2(\xi _A)_X=4\mathrm{Re}\xi _A,[\sigma X,[X,\sigma X]].$$
The second derivative, however, is slightly more involved:
$$\begin{array}{cc}\hfill dId& \eta _2(\xi _A,\xi _B)_X\hfill \\ & =\xi _A(Id\eta _2(\xi _B))\xi _B(Id\eta _2(\xi _A))Id\eta _2([\xi _A,\xi _B])\hfill \\ & =4\begin{array}{cc}\hfill \mathrm{Im}\{& \xi _B,[\sigma \xi _A,[X,\sigma X]]+\xi _B,[\sigma X,[\xi _A,\sigma X]]\hfill \\ & +\xi _B,[\sigma X,[X,\sigma \xi _A]]+[B,\xi _A],[\sigma X,[X,\sigma X]]\hfill \\ & \xi _A,[\sigma \xi _B,[X,\sigma X]]\xi _A,[\sigma X,[\xi _B,\sigma X]]\hfill \\ & \xi _A,[\sigma X,[X,\sigma \xi _B]][A,\xi _B],[\sigma X,[X,\sigma X]]\hfill \\ & +[[A,B],X],[\sigma X,[X,\sigma X]]\}\hfill \end{array}\hfill \\ & =4\begin{array}{cc}\hfill \mathrm{Im}\{& [\xi _A,\sigma \xi _B],[X,\sigma X]+[\xi _B,\sigma \xi _A],[X,\sigma X]\hfill \\ & +[\sigma X,\xi _A],[X,\sigma \xi _B][X,\sigma \xi _A],[\sigma X,\xi _B]\}\hfill \end{array}\hfill \\ & =8\mathrm{Im}\xi _A,[\sigma \xi _B,[X,\sigma X]]+[\sigma X,[X,\sigma \xi _B]].\hfill \end{array}$$
Combining these formulæ gives the claimed result. ∎
The two-form $`\omega _I`$ is our candidate for a Kähler form on $`𝒪`$.
###### Remark 3.2.
The corresponding symmetric bilinear form is given by $`g(\xi _A,\xi _B)=\omega _I(I\xi _A,\xi _B)`$ and is simply the right-hand side of equation (3.2) with ‘$`\mathrm{Im}`$’ replaced by ‘$`\mathrm{Re}`$’ throughout.
We will eventually require $`g`$ to be positive definite. However for now simply assume that $`g`$ is non-degenerate and define an endomorphism $`J`$ of $`T_X𝒪`$ by
(3.4)
$$g(\xi _A,\xi _B)=\mathrm{Re}\omega _c^𝒪(J\xi _A,\xi _B).$$
###### Lemma 3.3.
The endomorphism $`J`$ of $`T_X𝒪`$ is given by
(3.5)
$$\begin{array}{cc}\hfill J\xi _A& =2\rho _1[X,\sigma \xi _A]\hfill \\ & +4\rho _2(2[X,[\sigma X,[X,\sigma \xi _A]]][X,[X,[\sigma X,\sigma \xi _A]]])\hfill \\ & 2\rho _{11}\sigma \xi _A,X[X,\sigma X]\hfill \\ & +4\rho _{12}\begin{array}{cc}\hfill (& \sigma \xi _A,[X,[\sigma X,X]][X,\sigma X]\hfill \\ & +\sigma \xi _A,X[X,[\sigma X,[X,\sigma X]]])\hfill \end{array}\hfill \\ & 8\rho _{22}\sigma \xi _A,[X,[\sigma X,X]][X,[\sigma X,[X,\sigma X]]].\hfill \end{array}$$
###### Proof.
Equation (3.4) implies $`g(\xi _A,\xi _B)=\mathrm{Re}J\xi _A,B`$, and then (3.2) gives the above formula for $`J`$, except that the coefficient of $`\rho _2`$ is
(3.6)
$$4\left([X,[\sigma \xi _A,[X,\sigma X]]]+[X,[\sigma X,[X,\sigma \xi _A]]]\right).$$
Using the Jacobi identity, we have
$$[\sigma \xi _A,[X,\sigma X]]=[X,[\sigma X,\sigma \xi _A]]+[\sigma X,[X,\sigma \xi _A]].$$
Applying this to the first term in (3.6) gives the result. ∎
At this stage there is no guarantee that $`J^2=1`$. It is imposing this condition that severely restricts the possibilities for $`\rho `$.
## 4. Small Nilpotent Orbits and Real Subalgebras
The nilpotent orbits in $`𝔤^{}`$ are partially ordered by saying $`𝒪_1𝒪_2`$ if and only if $`𝒪_1\overline{𝒪_2}`$. When $`𝔤^{}`$ is simple, there is a unique non-zero orbit $`𝒪_{\text{min}}`$ which is minimal for this partial order. This orbit is of cohomogeneity one with respect to the action of the compact group $`G`$, and for each $`X𝒪_{\text{min}}`$, the subalgebra spanned by $`\{X,\sigma X\}`$ is isomorphic to $`𝔰𝔩(2,)`$ and is the complexification of an $`𝔰𝔲(2)`$-subalgebra of $`𝔤`$. Note that $`𝒪_{\text{min}}`$ is the orbit of a root vector for the longest root.
In general, the Jacobsen-Morosov Theorem says that each nilpotent element $`X`$ lies in an $`𝔰𝔩(2,)`$-subalgebra (see e.g. ). However, in general this subalgebra is not $`\sigma `$-invariant. The following result is usually attributed to Borel .
###### Proposition 4.1 (Borel).
Each nilpotent orbit $`𝒪`$ contains an element $`X`$ such that the linear span of $`\{X,\sigma X,[X,\sigma X]\}`$ is a real subalgebra isomorphic to $`𝔰𝔩(2,)`$.
###### Proof.
Fix $`X^{}`$ in $`𝒪`$ and take any $`𝔰𝔩(2,)`$ containing $`X`$. There are $`H`$ and $`Y`$ in $`𝔰𝔩(2,)`$ such that $`[H,X^{}]=2X^{}`$, $`[X^{},Y]=H`$ and $`[H,Y]=2Y`$. The element $`H`$ is thus semi-simple in $`𝔰𝔩(2,)`$ and hence in $`𝔤^{}`$, so we find a Cartan subalgebra $`𝔱`$ of $`𝔤^{}`$ containing $`H`$ and choose a system of positive roots $`\mathrm{\Delta }^+`$ so that $`X`$ lies in a sum of positive root spaces. The pair $`(𝔱,\mathrm{\Delta }^+)`$ has an associated real structure $`\sigma ^{}`$, which maps $`\mathrm{\Delta }^+`$ to $`\mathrm{\Delta }^{}`$ and defines a compact real form of $`𝔤^{}`$. Now all compact real forms of $`𝔤^{}`$ are conjugate, so there is a $`gG^{}`$ such that $`\mathrm{Ad}_g(\sigma ^{}A)=\sigma \mathrm{Ad}_gA`$, for all $`A𝔤^{}`$. Taking $`X=\mathrm{Ad}_gX^{}`$ gives an element of $`𝒪`$ of the desired type. ∎
Let us recall the Morse theory picture of the nilpotent variety described in (see also ). Each nilpotent orbit $`𝒪`$ admits a certain free action of $`^{}/\{\pm 1\}`$. The quotient $`𝔐(𝒪)=𝒪/^{}`$ may be described as a submanifold of the Grassmannian $`\stackrel{~}{\mathrm{Gr}}_3(𝔤)`$ of oriented three-planes in the real Lie algebra $`𝔤`$. One defines a functional $`\psi :\stackrel{~}{\mathrm{Gr}}_3(𝔤)`$ by $`\psi (V)=e_1,[e_2,e_3]`$, where $`\{e_1,e_2,e_3\}`$ is an oriented orthonormal basis for $`V`$. Away from zero, $`\psi `$ is a non-degenerate $`G`$-equivariant Morse function in the sense of Bott. The points on the non-zero critical sets correspond to subalgebras of $`𝔤`$ isomorphic to $`𝔰𝔲(2)`$. The set of real $`𝔰𝔲(2)`$-subalgebras associated to $`𝒪`$ via Proposition 4.1, oriented so that $`\psi `$ is positive, forms a non-zero critical manifold $`𝒞(𝒪)`$. The manifold $`𝔐(𝒪)`$ is the stable manifold attached to $`𝒞(𝒪)`$. The partial order on stable manifolds for the gradient flow induces the partial order $``$ on nilpotent orbits. In particular, the maximum of $`\psi `$ is achieved on $`𝔐(𝒪_{\text{min}})`$.
We are interested in orbits of cohomogeneity two. These were computed in and are the orbits listed in Table 1. We say that a nilpotent orbit $`𝒪`$ is *next-to-minimal* if $`𝒪𝒪_{\text{min}}`$ and there is no orbit $`𝒪^{}`$ with $`𝒪𝒪^{}𝒪_{\text{min}}`$. It is pleasing to note that the orbits listed in Table 1 are precisely the next-to-minimal orbits in the given algebras. The only next-to-minimal orbit that does not occur is that in $`𝔰𝔩(3,)`$, which is cohomogeneity four.
Recall that according to Proposition 5.1 elements of cohomogeneity-one nilpotent orbits lie in a real $`𝔰𝔩(2,)`$, i.e., in a $`\sigma `$-invariant rank one Lie algebra. It is remarkable that the elements of cohomogeneity-two orbits lie in $`\sigma `$-invariant rank two Lie algebras. The following can be thought of as a cohomogeneity-two version of Borel’s result.
###### Theorem 4.2.
Suppose $`G`$ is a compact simple Lie group and that $`𝒪`$ is a nilpotent orbit in $`𝔤^{}`$ of cohomogeneity two. Suppose $`X`$ is an element of $`𝒪`$ that does not lie in a real $`𝔰𝔩(2,)`$-subalgebra.
Let $`𝔥_X^{}`$ be the subalgebra of $`𝔤^{}`$ generated by $`X`$ and $`\sigma X`$. Then $`𝔥_X^{}`$ is isomorphic to $`𝔰𝔬(4,)`$, unless $`𝔤=𝔤_2`$, in which case $`𝔥^{}𝔤_2^{}`$. In all cases, the embedding $`𝔥^{}𝔤^{}`$ is a homothety with respect to the Killing forms.
###### Proof.
Consider the Morse theory picture. Firstly, in $`𝔤^{}`$, the closure of $`𝒪`$ is $`𝒪𝒪_{\text{min}}\{0\}`$. In $`\stackrel{~}{\mathrm{Gr}}_3(𝔤)`$ we have $`\overline{𝔐(𝒪)}=𝔐(𝒪)𝔐(𝒪_{\text{min}})`$. For the orbits of cohomogeneity two, $`𝔐(𝒪)`$ is a manifold of cohomogeneity one; the usual scaling by $`_{>0}`$, which is also part of the $`^{}/\{\pm 1\}`$-action, is transverse to the $`G`$-orbits on $`𝒪`$. Suppose $`\mathrm{}`$ is a curve in $`\overline{𝔐(𝒪)}`$ joining a point of $`𝒞(𝒪)`$ to a point of $`𝔐(𝒪_{\text{min}})`$. Then the fact that $`𝔐(𝒪)`$ is the stable manifold for the gradient flow of the $`G`$-invariant functional $`\psi `$, implies that the image of $`\mathrm{}`$ in $`𝔐(𝒪)/G`$ is the whole (one-dimensional) quotient space.
Now to parameterise $`\overline{𝒪}/G`$ it is enough to find a two-dimensional family of elements which is invariant under scaling and contains an element lying over $`𝒞(𝒪)`$ and an element of $`𝒪_{\text{min}}`$. When $`𝔤𝔤_2`$, we will find such a family lying in a $`\sigma `$-invariant $`𝔰𝔬(4,)`$-subalgebra of $`𝔤^{}`$.
The Lie algebra $`𝔰𝔬(4,)`$ splits as $`𝔰𝔩(2,)_+𝔰𝔩(2,)_{}`$. It contains three non-trivial nilpotent orbits: $`𝒪_\pm `$, the non-trivial nilpotent orbits in the factor $`𝔰𝔩(2,)_\pm `$; and $`𝒪_\mathrm{\Delta }=𝒪_+\times 𝒪_{}`$. The orbits $`𝒪_\pm `$ are cohomogeneity one and $`𝒪_\mathrm{\Delta }`$ is cohomogeneity two. Our orbit $`𝒪`$ will meet $`𝔰𝔬(4,)`$ in $`𝒪_\mathrm{\Delta }`$ and $`𝒪_+𝒪_{}`$ will be the intersection $`𝒪_{\text{min}}𝔰𝔬(4,)`$.
For the classical groups, we use the Jordan normal forms for elements of the orbits. For type $`A_n`$, the Jordan normal form is $`(2^21^{n3})`$ and the matrices
$$X_{s,t}=\left(\begin{array}{ccccc}0& s& & & \\ 0& 0& & & \\ & & 0& t& \\ & & 0& 0& \\ & & & & \text{0}\end{array}\right)$$
lie in the orbit unless $`s`$ or $`t`$ is zero. They also lie in the $`𝔰𝔲(2)𝔰𝔲(2)`$-subalgebra contained in the first two $`(2\times 2)`$ diagonal blocks. The matrix $`X_{1,1}`$ lies in a real $`𝔰𝔩(2,)`$-subalgebra and $`X_{1,0}`$ is in $`𝒪_{\text{min}}`$. So this two parameter family is as required. Exactly the same technique works for $`C_n`$.
For types $`B`$ and $`D`$, we are looking at matrices in $`𝔰𝔬(n,)`$. It is convenient to take $`𝔰𝔬(n,)`$ to be the set of complex $`(n\times n)`$ matrices $`A`$ such that $`A^tB+BA=0`$, where $`B`$ is the matrix with $`1`$’s down the anti-diagonal and $`0`$’s elsewhere. For Jordan form $`(31^{n3})`$, we just take an $`𝔰𝔬(4,)`$-subalgebra containing the Jordan block $`(3)`$. When the Jordan type is $`(2^41^{n8})`$, and the Lie algebra type is not $`D_{2n}`$, we have the same situation as for $`A_n`$, but now the blocks come in pairs. Thus the two families one considers are
(4.1)
$$\left(\begin{array}{ccccccccc}0& s& & & & & & & \\ 0& 0& & & & & & & \\ & & 0& t& & & & & \\ & & 0& 0& & & & & \\ & & & & \text{0}& & & & \\ & & & & & 0& t& & \\ & & & & & 0& 0& & \\ & & & & & & & 0& s\\ & & & & & & & 0& 0\end{array}\right)\text{and}\left(\begin{array}{cccccc}\text{0}& & & & & \text{0}\\ & 0& s& t& 0& \\ & & 0& 0& t& \\ & & & 0& s& \\ & & & & 0& \\ \text{0}& & & & & \text{0}\end{array}\right).$$
For $`D_{2n}`$, the matrices of Jordan type $`(2^41^{n8})`$ form a single $`\text{O}(n,)`$-orbit but split into two orbits $`(2^41^{n8})_\pm `$ under the action of $`\text{SO}(n,)`$. We thus obtain $`(2^41^{n8})_{}`$ from $`(2^41^{n8})_+`$ by conjugating by an element $`W`$ of determinant $`1`$ in $`\text{O}(n,)`$. One now considers three representative matrices, two as in (4.1) and
$$\left(\begin{array}{ccccccccc}0& 0& & & & & & & \\ 0& 0& & & & & & & \\ & & 0& t& & & & & \\ & & 0& 0& & & & & \\ & & & & \text{0}& & & & \\ & & & & & 0& t& & \\ & & & & & 0& 0& & \\ s& 0& & & & & & 0& 0\\ 0& s& & & & & & 0& 0\end{array}\right),\text{obtained with }W=\left(\begin{array}{ccccc}0& & & & 1\\ & 1& & & \\ & & \mathrm{}& & \\ & & & 1& \\ 1& & & & 0\end{array}\right).$$
In all cases, the matrix lies over $`𝒞(𝒪)`$ when $`s=t0`$ and is in $`𝒪_{\text{min}}`$ when $`t=0`$ and $`s0`$.
For the exceptional Lie algebras we use the Beauville bundle $`N(𝒪)`$ as a tool for computation. This bundle is defined as follows. Find a real $`𝔰𝔩(2,)`$-subalgebra associated to $`𝒪`$ and let $`\{e,f,h\}`$ be a basis for this subalgebra, with $`f=\sigma e`$, $`h=[e,f]`$ and $`[h,e]=2e`$. The eigenvalues of $`\mathrm{ad}h`$ on $`𝔤^{}`$ are known to be integers (see ). Let $`𝔤(i)`$ be the $`i`$-eigenspace of $`\mathrm{ad}h`$. Put
$$𝔭=\underset{i0}{}𝔤(i)\text{and}𝔫=\underset{i2}{}𝔤(i).$$
Then $`𝔭`$ is a parabolic subalgebra of $`𝔤^{}`$ and the corresponding homogeneous space $`=G^{}/P`$ is a flag manifold. The subalgebra $`𝔫`$ is preserved by the adjoint action of $`P`$ and the Beauville bundle $`N(𝒪)`$ is defined to be the bundle over $``$ associated to $`𝔫`$, i.e.,
$$N(𝒪)=G^{}\times _P𝔫.$$
The important property of $`N(𝒪)`$ is that it contains $`𝒪`$ as an open dense $`G^{}`$-orbit.
Now each flag manifold is a homogeneous manifold for the action of the compact group. So $`=G/K`$ for some compact subgroup $`K`$ of $`G`$. (In fact, the Lie algebra $`𝔨`$ of $`K`$ is given by $`𝔨^{}=𝔤(0)`$.) The Beauville bundle is then $`G\times _K𝔫`$ and the cohomogeneity of $`𝒪`$ is the cohomogeneity of the action of $`K`$ on $`𝔫`$. Choose a Cartan subalgebra in $`𝔤(0)`$ and a root system for $`𝔤(0)`$ with all root spaces in $`𝔭`$. Note that, by definition, the weighted Dynkin diagram for $`𝒪`$ gives the eigenvalues of $`\mathrm{ad}h`$ on the positive simple root spaces, from which all the other eigenvalues are easily computed.
In the case of cohomogeneity-two orbits not in $`𝔤_2`$, we find that $`𝔫^2V`$ as a representation of $`K=\text{SO}(2)L`$, with $`V`$ irreducible and $`L`$ acting two-point transitively on the unit sphere in $`V`$. Thus under the action of $`K`$, we can move any nilpotent element in $`𝒪`$ into any complex two-dimensional subspace (for the complex structure induced by the action of $`\text{SO}(2)`$). We then find root spaces $`𝔤_\alpha `$ and $`𝔤_\beta `$ contained in $`𝔫`$ with $`\alpha `$ and $`\beta `$ orthogonal long roots such that $`\alpha \pm \beta `$ is not a root. The $`\sigma `$-invariant subalgebra containing these root spaces is then the required $`𝔰𝔬(4,)`$. For the relevant four exceptional algebras, this information is given in Table 2.
For $`G_2`$, the isotropy group for $``$ is $`K=\text{U}(1)\text{SU}(2)`$. We have $`𝔫=𝔤(2)+𝔤(3)`$ with $`𝔤(2)L^2`$ and $`𝔤(3)L^3S^1`$, where $`L=`$ and $`S^1=^2`$ are the fundamental representations of $`\text{U}(1)`$ and $`\text{SU}(2)`$, respectively. The orbit $`𝒪`$ in this case is the orbit of short root vectors and the subalgebra generated by $`X`$ and $`\sigma X`$ contains both short and long roots, so is all of $`𝔤_2`$. ∎
## 5. Two Models
In this section we compute Kähler potentials for hyperKähler structures on two particular nilpotent orbits: one in $`𝔰𝔩(2,)`$ and the other in $`𝔰𝔬(4,)`$. These results will be used in the next section to derive the hyperKähler potentials for cohomogeneity-two orbits. In view of Theorem 4.2, we consider these cases with inner products that are multiples of that given by the Killing form.
We start by considering $`𝔤^{}=𝔰𝔩(2,)`$ with inner product $`k^2,_{𝔰𝔩(2)}`$, where $`k>0`$ is constant and $`,_{𝔰𝔩(2)}`$ is negative of the Killing form. This Lie algebra contains only one non-trivial nilpotent orbit $`𝒪`$ consisting of the $`(2\times 2)`$ matrices $`X`$ such that $`X^2=0`$ and $`X0`$. The orbit is the minimal nilpotent orbit in $`𝔰𝔩(2,)`$ and is of cohomogeneity one under the adjoint action of $`\text{SU}(2)`$. In fact two elements of the orbit have the same norm if and only if they are $`\text{SU}(2)`$-conjugate. Thus any $`\text{SU}(2)`$-invariant Kähler potential $`\rho `$ on $`𝒪`$ is a function of just $`\eta =k^2X,\sigma X_{𝔰𝔩(2)}`$. Write
(5.1)
$$e=\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right),f=\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right)\text{and}h=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
Then $`\{e,f,h\}`$ is an $`𝔰𝔩(2,)`$ triple, with $`f=\sigma e`$ and $`h=\sigma h`$. Using the action of $`\text{SU}(2)`$, we may assume that $`X=te`$, for some $`t>0`$. The tangent space $`T_X𝒪=[X,𝔰𝔩(2,)]`$ is spanned by $`e`$ and $`h`$. If we consider the complex symplectic form $`k^2\omega _C^𝒪`$, we may perform the same calculations as in (3.5) and get
$$J_Xe=2t(\rho ^{}+\eta \rho ^{\prime \prime })h\text{and}J_Xh=4t\rho ^{}e,$$
where $`\rho ^{}=d\rho /d\eta `$, etc. As $`\eta (X)=4k^2t^2`$, the condition that $`J^2=1`$ is equivalent to
(5.2)
$$2\eta \rho ^{}(\rho ^{}+\eta \rho ^{\prime \prime })=k^2.$$
The left-hand side is simply the derivative of $`(\eta \rho ^{})^2`$ with respect to $`\eta `$, so
(5.3)
$$\rho _{}^{}{}_{}{}^{2}=(k^2\eta +c)/\eta ^2,$$
for some constant $`c`$. In order to have the potential defined on the whole orbit we need $`c0`$. The corresponding metric may be calculated as in Remark 3.2 and is given by
(5.4)
$$\begin{array}{cc}\hfill g(\xi _A,\xi _B)=\frac{2k^4}{\eta }\mathrm{Re}(& \rho ^{}\left(\xi _A,\sigma \xi _BX,\sigma X\xi _A,\sigma XX,\sigma \xi _B\right)\hfill \\ & +\frac{k^2}{2\eta \rho ^{}}\xi _A,\sigma XX,\sigma \xi _B),\hfill \end{array}$$
which is positive definite provided we take the positive square root in (5.3). Now (5.3) determines $`\rho `$ up to an additive constant, and this is enough to fix the metric structure.
###### Proposition 5.1.
For fixed $`k`$, the nilpotent orbit in $`𝔰𝔩(2,)`$ has a one-parameter family of $`\text{SU}(2)`$-invariant hyperKähler metrics with a Kähler potential and with $`k^2\omega _c^𝒪`$ as the complex symplectic form. The $`G`$-invariant Kähler potential $`\rho `$ is given by
(5.5)
$$\rho ^{}=\frac{1}{\eta }\sqrt{k^2\eta +c},$$
where $`c0`$ is a constant and $`\eta (X)=k^2X,\sigma X_{𝔰𝔩(2)}`$. ∎
Note that if we rewrite everything in terms of the variable $`t`$, we get
(5.6)
$$\frac{d\rho }{dt}(te)=\sqrt{(2k)^4+\frac{4c}{t^2}}.$$
###### Proposition 5.2.
The Kähler potential $`\rho `$ of Proposition 5.1 is a hyperKähler potential if and only if $`c=0`$. In this case, $`\rho =2k\sqrt{\eta }`$ and $`\rho (te)=4k^2t`$.
###### Proof.
Let $`Y=(d\rho )^{\mathrm{}}`$ be the vector field dual to $`d\rho `$. If $`\rho `$ is a hyperKähler potential then $`IY`$ is an isometry preserving $`I`$ and $`\rho `$ is the corresponding moment map \[23, Proposition 5.5\]. Now $`d\rho =\rho ^{}d\eta =2k^2\rho ^{}\mathrm{Re},\sigma X`$, whereas, by Remark 3.2,
$$g(Y,\xi _A)=2\mathrm{Re}(\rho ^{}k^2\xi _A,\sigma Y+\rho ^{\prime \prime }k^4\xi _A,\sigma XX,\sigma Y).$$
So we have
$$\rho ^{}X=\rho ^{}Y+\rho ^{\prime \prime }k^2Y,\sigma XX,$$
which implies that $`Y=\lambda X`$ with
$$\begin{array}{cc}\hfill \lambda & =\frac{\rho ^{}}{\rho ^{}+\eta \rho ^{\prime \prime }}=\frac{2\eta \rho _{}^{}{}_{}{}^{2}}{k^2}=2+\frac{2c}{k^2\eta },\hfill \end{array}$$
using (5.2). The vector field $`X`$ is generated by scaling in the nilpotent orbit, so $`IX`$ preserves the complex structure $`I`$. Now
$$\begin{array}{cc}\hfill (L_{IY}I)(Z)& =[\lambda IX,IZ]I[\lambda IX,Z]\hfill \\ & =\lambda (L_{IX}I)Z((IZ)\lambda )IX+(Z\lambda )X,\hfill \end{array}$$
so with $`L_{IX}I=0`$, we have $`L_{IY}I=0`$ only if $`\lambda `$ is constant. But this is exactly the requirement that $`c=0`$. ∎
###### Remark 5.3.
The substitution $`k^2\eta +c=(\frac{r}{2})^4`$ in equation 5.4 shows that these are the Eguchi-Hanson metrics (cf. ). See for details.
Let us now turn to the regular nilpotent orbit $`𝒪_\mathrm{\Delta }`$ in $`𝔰𝔬(4,)`$. As in the proof of Theorem 4.2, we write $`𝔰𝔬(4,)=𝔰𝔩(2,)_+𝔰𝔩(2,)_{}`$ and note that $`𝒪_\mathrm{\Delta }=𝒪_+\times 𝒪_{}`$ where $`𝒪_\pm `$ is the nilpotent orbit in $`𝔰𝔩(2,)_\pm `$. Let $`\{e_\pm ,f_\pm ,h_\pm \}`$ be bases for $`𝔰𝔩(2,)_\pm `$ as in (5.1). Again we will use the inner product which is $`k^2,_{𝔰𝔬(4)}`$.
Using the action of $`\text{SO}(4)`$, we may take our representative element $`X`$ of $`𝒪_\mathrm{\Delta }`$ to be $`X=X_++X_{}=se_++te_{}`$ with $`s,t>0`$. We have one invariant for each $`𝔰𝔩(2,)`$: we write $`\eta _\pm =k^2X_\pm ,\sigma X_\pm _{𝔰𝔩(2)}`$, so $`\eta _+=4k^2s^2`$, etc. Let $`\rho _+=\rho /\eta _+`$, etc. Then we may calculate the Kähler form $`\omega _I`$ and the candidate almost complex structure $`J`$ as in §3. For the Kähler form we get
$$\begin{array}{cc}\hfill \omega _I(\xi _A,\xi _B)=2k^2\mathrm{Im}(& \rho _+\xi _A^+,\sigma \xi _B^++\rho _{}\xi _A^{},\sigma \xi _B^{}\hfill \\ & +\rho _{++}k^2\xi _A^+,\sigma X_+\sigma \xi _B^+,X_+\hfill \\ & +\rho _+k^2(\xi _A^+,\sigma X_+\sigma \xi _B^{},X_{}\hfill \\ & +\xi _A^{},\sigma X_{}\sigma \xi _B^+,X_+)\hfill \\ & +\rho _{}k^2\xi _A^{},\sigma X_{}\sigma \xi _B^{},X_{}),\hfill \end{array}$$
where $`\xi _A^+=[A,X_+]`$, etc. The endomorphism $`J`$ is given by
$$\begin{array}{c}J_X(\xi _A^+)=2\rho _+[X_+,\sigma \xi _A^+]\hfill \\ \hfill 2k^2\sigma \xi _A^+,X_+\left(\rho _{++}[X_+,\sigma X_+]+\rho _+[X_{},\sigma X_{}]\right)\end{array}$$
and a similar expression for $`\xi _A^{}`$. In particular, $`J_Xh_+=4s\rho _+e_+`$ and
$$J_Xe_+=2s(\rho _++\eta _+\rho _{++})h_++2\frac{t^2}{s}\eta _+\rho _+h_{}.$$
Thus the $`𝔰𝔩(2,)_{}`$-component of $`J_X^2h_+`$ is a constant times $`\eta _+\eta _{}\rho _+\rho _+h_{}`$. For $`J_X^2`$ to be $`1`$, we need $`\rho _+\rho _+=0`$, which implies $`(\rho _+^2)/\rho _{}=0`$ and hence $`\rho _+=0`$. Thus $`J_X`$ preserves the $`𝔰𝔩(2,)`$-summands of $`𝔰𝔬(4,)`$.
###### Proposition 5.4.
Any hyperKähler structure on the regular orbit $`𝒪_\mathrm{\Delta }=𝒪_+\times 𝒪_{}`$ of $`𝔰𝔬(4,)`$ which is $`\text{SO}(4)`$-invariant, admits a Kähler potential and has complex-symplectic form $`k^2\omega _c^{𝒪_\mathrm{\Delta }}`$, is a product of $`\text{SU}(2)`$-invariant structures on the factors $`𝒪_\pm `$, and these are given by Proposition 5.1. ∎
## 6. Potentials for Next-to-Minimal Orbits
We now come to the main result of this paper. We consider next-to-minimal orbits with compatible $`G`$-invariant hyperKähler metrics, except for $`G=\text{SU}(3)`$. We show that such metrics admitting a hyperKähler potential are unique, and we calculate the potential.
If we assume that the potential is only Kähler, we still have uniqueness in some cases, but we get a list of exceptions: orbits which admit a one-parameter family of hyperKähler metrics. These can be thought of as a generalisation of the Eguchi-Hanson metric (cf. Remark 5.3).
###### Theorem 6.1.
Suppose $`G`$ is a compact simple Lie group and $`𝒪`$ is a nilpotent orbit in $`𝔤^{}`$ of cohomogeneity two.
1. $`𝒪`$ admits a unique $`G`$-invariant compatible hyperKähler metric with hyperKähler potential. This potential is given by
(6.1)
$$\rho =2k\sqrt{\eta _1+2\sqrt{\frac{1}{2}\eta _1^2k^2\eta _2}}$$
for $`𝔤𝔤_2`$, where the constant $`k`$ is given in Table 3, and, for $`𝔤_2`$,
(6.2)
$$\rho =\sqrt{8}\sqrt{\eta _1+\sqrt{6}\sqrt{\eta _1^24\eta _2}}.$$
2. The above metric on $`𝒪`$ is in fact a unique $`G`$-invariant compatible hyperKähler metric with a *Kähler* potential unless $`𝔤`$ is one of $`𝔰𝔭(2)𝔰𝔬(5)`$, $`𝔰𝔲(4)𝔰𝔬(6)`$, $`𝔰𝔬(8)`$ or $`𝒪`$ is of Jordan type $`(31^{n3})`$ in $`𝔰𝔬(n)`$. In these cases, the metric lies in a one-parameter family of hyperKähler metrics with Kähler potentials.
###### Remark 6.2.
Note that the Theorem provides hyperKähler potentials for all next-to-minimal orbits, except when $`𝔤=𝔰𝔲(3)`$. However, the potential in this remaining case was computed in , see also .
We divide the proof of the Theorem into three parts.
### 6.1. The General Case
This is when $`𝔤`$ is neither $`𝔰𝔲(3)`$ nor $`𝔤_2`$.
Let $`X`$ be a generic element of $`𝒪`$. By Theorem 4.2, $`X`$ lies in the regular orbit $`𝒪_\mathrm{\Delta }`$ of a real $`𝔰𝔬(4,)`$-subalgebra.
For $`\rho (\eta _1,\eta _2)`$ to be a hyperKähler potential for $`𝒪`$ it is necessary that $`\rho `$ is a Kähler potential for an invariant hyperKähler structure on $`𝒪_\mathrm{\Delta }`$. To see this, first note that equation (2.1) is invariant by pull-back under the inclusion map $`𝒪_\mathrm{\Delta }𝒪`$. Now equation (3.5) shows that $`J\xi _A`$ remains in the subalgebra generated by $`A`$, $`X`$ and $`\sigma X`$. Thus if $`A𝔰𝔬(4,)`$, so is $`JA`$ and thus $`𝒪_\mathrm{\Delta }`$ is a hyperKähler submanifold of $`𝒪`$.
As in §5, write $`𝔰𝔬(4,)=𝔰𝔩(2,)_+𝔰𝔩(2,)_{}`$, $`𝒪_\mathrm{\Delta }=𝒪_+\times 𝒪_{}`$ and $`X=X_++X_{}=se_++te_{}`$. Our two invariants on $`𝒪`$ are given by
$$\begin{array}{cc}\hfill \eta _1(X)& =X,\sigma X_𝔤=se_++te_{},sf_++tf_{}_𝔤\hfill \\ & =(s^2+t^2)k^2e,f_{𝔰𝔲(2)}=4k^2(s^2+t^2),\hfill \end{array}$$
and a similar computation gives
$$\eta _2(X)=8k^2(s^4+t^4),$$
where $`k^2`$ is the constant such that $`,_𝔤|_{𝔰𝔬(4,)}=k^2,_{𝔰𝔬(4)}`$. Now
$$\begin{array}{cc}\hfill d\rho & =\rho _1d\eta _1+\rho _2d\eta _2\hfill \\ & =8k^2\left(s(\rho _1+4s^2\rho _2)ds+t(\rho _1+4t^2\rho _2)dt\right),\hfill \end{array}$$
so $`\rho _s:=\rho /s=8k^2s(\rho _1+4s^2\rho _2)`$, etc., and solving for $`\rho _1`$ and $`\rho _2`$ we get
$$\rho _1=\frac{t^3\rho _ss^3\rho _t}{8k^2st(s^2t^2)},\rho _2=\frac{t\rho _ss\rho _t}{32k^2st(s^2t^2)}.$$
Note that, by Proposition 5.4 and (5.6), $`\rho _{s}^{}{}_{}{}^{2}=16k^4+c_+/s^2`$ and $`\rho _{t}^{}{}_{}{}^{2}=16k^4+c_{}/t^2`$, for some constants $`c_\pm `$.
The elements $`X_+`$ and $`X_{}`$ lie in the closure of $`𝒪_\mathrm{\Delta }`$ and hence of $`𝒪`$; so $`X_\pm `$ lie in the minimal nilpotent orbit of $`𝔤^{}`$. We deduce that $`M_+:=G/N(\text{SU}(2)_+)`$ is a Wolf space and hence, since $`\text{SU}(2)_+`$ corresponds to a highest root ,
(6.3)
$$𝔤^{}=𝔰𝔩(2,)_++𝔨_++S_+^1E_+,$$
where $`𝔨_+`$ commutes with $`𝔰𝔩(2,)`$, $`E_+`$ is a non-trivial representation of $`𝔨_+`$ and $`S_+^1^2`$ is the fundamental representation of $`𝔰𝔩(2,)_+`$. On the other hand, we have a similar decomposition of $`𝔤^{}`$ corresponding to $`𝔰𝔩(2,)_{}`$. As $`𝔰𝔩(2,)_+`$ and $`𝔰𝔩(2,)_{}`$ commute with each other, we deduce that $`𝔰𝔩(2,)_{}𝔨_+`$ and that $`E_+S_{}^1`$. So as an $`𝔰𝔬(4,)`$-module, $`𝔤^{}`$ always contains a copy of $`S_+^1S_{}^1`$.
On the orthogonal complement to $`𝔰𝔬(4,)`$, we have, from (3.5),
(6.4)
$$\begin{array}{cc}\hfill J_X\xi _A& =2\rho _1[X,\sigma \xi _A]\hfill \\ & +4\rho _2\left(2[X,[\sigma X,[X,\sigma \xi _A]]][X,[X,[\sigma X,\sigma \xi _A]]]\right).\hfill \end{array}$$
First suppose that $`E_+`$ contains a trivial $`𝔰𝔩(2,)_{}`$-module $`^r`$; take $`r`$ maximal. The real structure $`\sigma `$ preserves the module $`S_+^1^r`$ and acts on $`S_+^1`$ as $`j`$, so $`^r`$ has a quaternionic structure $`𝔧`$ and is even-dimensional. Choose a basis for $`S_+^1`$ so that $`\mathrm{ad}e_+`$ acts as $`\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right)`$. Then any tangent vector $`\xi _AS_+^1^r`$ has the form $`\left(\begin{array}{c}1\\ 0\end{array}\right)v`$ and we have
$$J_X\xi _A=2s(\rho _1+4s^2\rho _2)\left(\begin{array}{c}1\\ 0\end{array}\right)𝔧v=\frac{1}{4k^2}\rho _s\left(\begin{array}{c}1\\ 0\end{array}\right)𝔧v.$$
Thus $`J^2=1`$ on $`S_+^1^r`$ if and only if $`\rho _{s}^{}{}_{}{}^{2}=16k^4`$. This implies that the constant $`c_+`$ is zero if $`E_+`$ has an trivial $`𝔰𝔩(2,)_{}`$-submodule.
The existence of an trivial $`𝔰𝔩(2,)_{}`$-submodule in $`E_+`$ is not guaranteed. However, we do always have an $`S_{}^1`$-summand, so we now consider the case when $`\xi _A`$ lies in an $`𝔰𝔬(4)`$-module $`S_+^1S_{}^1`$. This is Killing orthogonal to $`𝔰𝔬(4,)`$. We choose bases so that $`\mathrm{ad}X`$ acts as
$$s\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right)\mathrm{Id}+t\mathrm{Id}\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right)$$
and $`\sigma =jj`$ for $`j`$ the standard quaternionic structure on $`S^1`$. The image of $`\mathrm{ad}X`$ is two-dimensional and spanned by
$$\xi _1:=\left(\begin{array}{c}1\\ 0\end{array}\right)\left(\begin{array}{c}1\\ 0\end{array}\right)\text{and}\xi _2:=s\left(\begin{array}{c}1\\ 0\end{array}\right)\left(\begin{array}{c}0\\ 1\end{array}\right)+t\left(\begin{array}{c}0\\ 1\end{array}\right)\left(\begin{array}{c}1\\ 0\end{array}\right).$$
These satisfy
$$[X,\xi _1]=0,[\sigma X,\xi _1]=\sigma \xi _2,$$
$$[X,\xi _2]=2st\xi _1\text{and}[\sigma X,\xi _2]=(s^2+t^2)\sigma \xi _1.$$
So, equation (6.4) gives
$$J\xi _1=2(\rho _1+4\rho _2(s^2+t^2))\xi _2,$$
$$J\xi _1=2(\rho _1(s^2+t^2)+4\rho _2(s^4+t^4))\xi _1.$$
Substituting for $`\rho _1`$ and $`\rho _2`$ in terms of $`\rho _s`$ and $`\rho _t`$, gives
$$J^2\xi _1=\frac{t^2\rho _{t}^{}{}_{}{}^{2}s^2\rho _{s}^{}{}_{}{}^{2}}{16k^4(t^2s^2)}\xi _1.$$
So $`J^2=1`$ on $`S_+^1S_{}^1`$ if and only if $`t^2\rho _{t}^{}{}_{}{}^{2}s^2\rho _{s}^{}{}_{}{}^{2}=16k^4(t^2s^2)`$. But $`\rho _{s}^{}{}_{}{}^{2}=16k^4+c_+/s^2`$, etc., so $`c_+=c_{}`$.
We conclude that if $`E_+`$ contains a trivial $`𝔰𝔩(2,)_{}`$-summand, then $`c_+=c_{}=0`$. This gives $`\rho _s=4k^2`$ and $`\rho _t=4k^2`$, so $`\rho (s,t)=4k^2(s+t)`$. Rewriting this in terms of $`\eta _1`$ and $`\eta _2`$ gives the potential in the Theorem. If $`E_+`$ does not have a trivial summand, we get a one-parameter family of potentials and hyperKähler metrics with $`c_+=c_{}`$.
It remains to determine the constant $`k`$ and when $`E_+`$ contains a trivial $`\text{SU}(2)_{}`$-module. The decomposition (6.3) gives the action of $`\mathrm{ad}e_+`$ and hence the Killing inner product $`e_+,\sigma e_+_𝔤`$ is $`4+dim_{}E_+`$, since $`e_+,\sigma e_+_{𝔰𝔲(2)_+}=4`$. So $`k^2=(4+dim_{}E_+)/4`$. Moreover, $`S_+^1E_+=TM_+`$, so $`dim_{}E_+`$ is half the real dimension of the Wolf space $`M_+`$, which may be found in, e.g., Besse \[3, p. 409\], or read-off from the discussion below. This leads to Table 3.
Finally, we determine the decompositions of $`E_+`$ under the action of $`𝔰𝔩(2,)_{}`$.
If $`G=\text{SU}(n)`$, then $`𝔨_+𝔲(n2)`$, and $`E_+=^{n2}`$ is the fundamental representation twisted by a representation of the central $`𝔲(1)`$. Now $`𝔰𝔩(2,)_{}`$ corresponds to a highest root vector in $`𝔨_+`$, so $`E_+=S_{}^1+^{n4}`$ as a $`𝔰𝔩(2,)_{}`$-module. So for $`n=4`$, we have a one-parameter family of potentials $`c_+=c_{}`$, and for $`n>4`$, the potential is unique.
For $`G=\text{Sp}(n)`$, $`𝔨_+𝔰𝔭(n1,)`$ and $`E_+^{2n2}^{n1}`$ is the fundamental representation. Under the highest root $`𝔰𝔩(2,)`$, this representation splits as $`S_{}^1+^{2n4}`$, so for $`n>1`$, we have a unique potential.
In the case $`G=\text{SO}(n)`$, there are two orbit types to consider. The centraliser $`𝔨_+=𝔰𝔩(2,)+𝔰𝔬(n4,)`$ and there are two choices for $`𝔰𝔩(2,)_{}`$, one in each summand of $`𝔨_+`$. When $`𝔰𝔩(2,)_{}=𝔰𝔩(2,)`$, we get $`E_+S_{}^1^{n4}`$, and there is a one-parameter family of potentials. On the other hand, if $`𝔰𝔩(2,)_{}`$ lies in the summand $`𝔰𝔬(n4,)`$, then $`E_+^2(S_{}^1+^{n8})`$. For $`n>8`$, this gives a unique potential, but for $`n=8`$, we again get a family.
We now come to the four exceptional cases. Firstly, if $`G=\text{F}_4`$, then $`𝔨_+𝔰𝔭(3,)`$ and if $`E=^3`$ is the fundamental representation, then $`E_+\mathrm{\Lambda }_0^3E=\mathrm{\Lambda }^3EE`$, is a $`14`$-dimensional irreducible representation. For a highest root $`𝔰𝔩(2,)_{}`$ in $`𝔰𝔭(3,)`$, we have $`ES_{}^1+^4`$ and hence $`E_+5S_{}^1+^4`$. So $`E_+`$ has a trivial summand and hence the potential is unique.
For $`G=\text{E}_6`$, $`𝔨_+𝔰𝔩(6,)`$ and $`E_+\mathrm{\Lambda }^{3,0}^6`$. Under a highest root $`𝔰𝔩(2,)_{}`$, we have $`\mathrm{\Lambda }^{1,0}^6S_{}^1+^4`$ and hence $`E_+=4S_{}^1+^8`$, giving a unique potential.
When $`G=\text{E}_7`$, $`𝔨_+𝔰𝔬(12,)`$ and $`E_+\mathrm{\Delta }_+^{12}`$, the positive spin representation. For a highest root $`𝔰𝔩(2,)_{}`$, the normaliser in $`𝔰𝔬(12,)`$ is $`𝔰𝔩(2,)_{}+𝔰𝔩(2,)+𝔰𝔬(8,)`$ and the fundamental representation of $`\text{SO}(12)`$ decomposes as $`^{12}S_{}^1S^1+V`$, where $`V^8`$ is the fundamental representation of $`𝔰𝔬(8,)`$. The spin representation splits as $`\mathrm{\Delta }_+^{12}S_{}^1\mathrm{\Delta }_+^8+S^1\mathrm{\Delta }_{}^8`$, and so $`E_+8S_{}^1+^{16}`$ has a trivial summand.
Finally, for $`G=\text{E}_8`$, $`𝔨_+𝔢_7^{}`$ and $`E_+100\stackrel{0}{0}00`$. A highest root $`𝔰𝔩(2,)_{}`$ in $`𝔢_7^{}`$ has centraliser $`𝔰𝔬(12,)`$ and $`E_+12S_{}^1+^{32}`$, where $`^{32}\mathrm{\Delta }_+^{12}`$. So again we get a unique potential.
### 6.2. The Exceptional Case $`\text{G}_2`$
The Dynkin diagram for the next-to-minimal orbit $`𝒪`$ in $`\text{G}_2`$ is $`0\begin{array}{c}>\\ --\hfill \\ --\hfill \\ --\hfill \end{array}1`$. This says that there is a basis $`\{\alpha ,\beta \}`$ for the simple positive roots, with $`\alpha `$ short and $`\beta `$ long, such that $`\mathrm{ad}h`$ acts on $`𝔤_\alpha `$ and $`𝔤_\beta `$ with eigenvalues $`1`$ and $`0`$ respectively. We thus have $`𝔤(2)=𝔤_{\beta +2\alpha }`$ and $`𝔤(3)=𝔤_{\beta +3\alpha }𝔤_{2\beta +3\alpha }`$. From the discussion in §4, the isotropy group $`\text{SU}(2)\text{U}(1)`$ of the Beauville bundle acts transitively on the unit sphere in $`𝔤(3)`$, so using the action of the compact group $`\text{G}_2`$, we can move a typical element of $`𝒪`$ to $`X𝔤_{\beta +2\alpha }𝔤_{2\beta +3\alpha }`$. We may thus write $`X=sE_{\beta +2\alpha }+tE_{2\beta +3\alpha }`$, with $`s,t>0`$, where $`E_i`$ are such that for $`F_i:=\sigma E_i`$ and $`H_i=[E_i,F_i]`$ we have $`[H_i,E_i]=2E_i`$.
At $`X`$, our two invariants are
$$\eta _1(X)=8(s^2+3t^2)\text{and}\eta _2(X)=16(s^4+6s^2t^2+3t^4).$$
As in the previous section, we compute $`J^2`$ on particular tangent vectors using (3.5) and then rewrite the equations in terms of $`s`$ and $`t`$. This is quite hard work to do by hand, and so we used Maple to do the following computations. The code for this is described in .
On $`𝔤_{\alpha +\beta }`$, one finds that $`J^2=1`$ only if
(6.5)
$$\frac{1}{64s}\rho _s(s\rho _s+t\rho _t)=1,$$
where $`\rho _s`$ is $`\rho /s`$, etc. Now $`X=[X,H_\beta H_\alpha ]=[X,3H_{2\beta +3\alpha }5H_{\beta +2\alpha }]`$, so $`X`$ is tangent to the orbit $`𝒪`$. The condition $`J^2X=X`$, gives the following three equations
(6.6a)
$$\begin{array}{cc}\hfill s(2s\rho _s+t\rho _t)\rho _{ss}+t(t\rho _t+3s\rho _s)\rho _{st}+t^2\rho _s\rho _{tt}& \\ \hfill +2(t\rho _t+s\rho _s)\rho _s& =128s,\hfill \end{array}$$
(6.6b)
$$9s\rho _s\rho _{ss}+(9t\rho _s+s\rho _t)\rho _{st}+t\rho _t\rho _{tt}+9\rho _s^2+\rho _t^2=576$$
(6.6c)
$$\begin{array}{cc}\hfill 3st(9t\rho _s+s\rho _t)\rho _{ss}st(s\rho _t3t\rho _s)\rho _{tt}& \\ \hfill +(3t(s^2+9t^2)\rho _s+s(3t^2s^2)\rho _t)\rho _{st}& =(s\rho _t9t\rho _s)(s\rho _t+3t\rho _s)\hfill \end{array}$$
by considering the components in $`𝔤_{\beta +2\alpha }`$, $`𝔤_{2\beta +3\alpha }`$ and $`𝔤_{2\beta +\alpha }`$. Considering $`9s\rho _s`$ times (6.6a) minus $`s(2\rho _ss+t\rho _t)`$ times $`(\text{6.6b})`$ gives a new equation not involving $`\rho _{ss}`$. In a similar way, we may eliminate $`\rho _{ss}`$ form the pair of equations (6.6a) and (6.6c). Eliminating $`\rho _{st}`$ from these two new equations not involving $`\rho _{ss}`$, we get the following equation which does not involve $`\rho _{tt}`$:
$$s^3t(2s\rho _s+t\rho _t)(s\rho _t9t\rho _s)^2=0.$$
Thus either
(6.7)
$$\text{(i)}\rho _t=2\frac{s}{t}\rho _s,\text{or}\text{(ii)}\rho _t=9\frac{t}{s}\rho _s.$$
In case (i), substituting into (6.5) one gets $`\rho _s^2=64`$, which has no (real) solutions. In case (ii), we have
$$\rho _s=\epsilon \frac{8s}{\sqrt{s^2+9t^2}},\rho _t=\epsilon \frac{72t}{\sqrt{s^2+9t^2}},$$
where $`\epsilon \{\pm 1\}`$. Integrating we find that
(6.8)
$$\rho =\epsilon 8\sqrt{s^2+9t^2}.$$
To get a positive-definite metric, take $`\epsilon =+1`$. Rewriting (6.8) in terms of $`\eta _1`$ and $`\eta _2`$ gives the claimed result. One may check directly that the resulting $`J`$ satisfies $`J^2=1`$ on the whole tangent space.
### 6.3. Uniqueness of HyperKähler Potentials
The only statement left to verify in the proof of Theorem 6.1, is that equations (6.1) and (6.2) give the unique compatible hyperKähler potentials on the orbits. In the cases, when the Kähler potential is unique there is nothing to prove, because the general theory gives the existence of such a potential. We may therefore assume we are in the general case and that our generic element $`X`$ lies in a real $`𝔰𝔬(4,)`$-subalgebra. Now $`𝒪_\mathrm{\Delta }`$ is a hyperKähler submanifold of $`𝒪`$, and so by (2.1), $`\rho `$ is a hyperKähler potential for $`𝒪`$ only if it restricts to a hyperKähler potential for $`𝒪_\mathrm{\Delta }`$. However, the hyperKähler structure on $`𝒪_\mathrm{\Delta }`$ is the product of two hyperKähler structures on $`𝔰𝔩(2,)`$-orbits and on each of these factors the hyperKähler potential is unique by Proposition 5.2. Thus there is only one hyperKähler potential compatible with the structure of $`𝒪`$. ∎
###### Remark 6.3.
The hyperKähler metrics constructed in Theorem 6.1 have an extra $`\text{U}(1)`$-symmetry given by $`Xe^{i\theta }X`$ which preserves the complex structure $`I`$ but moves $`J`$. In the case of $`𝔰𝔩(2,)`$, the metrics are of Bianchi type $`IX`$ and it is known, e.g., from , that there are triaxial hyperKähler metrics that do not have $`\text{U}(2)`$-symmetry. Thus concentrating on metrics admitting a Kähler potential is a genuine restriction.
###### Remark 6.4.
The one-parameter families in Theorem 6.1 occur exactly when $`E_+^sS_{}^1`$. Considering the weights of the action of a semi-simple element in the diagonal $`𝔰𝔩(2,)`$-subalgebra of $`𝔰𝔬(4,)`$ on $`𝔤^{}`$, we see that this exactly the case when $`𝔤(1)=0`$. This says that the Beauville bundle coïncides with the cotangent bundle $`T^{}`$, rather than being a proper subbundle. In the case of the one-parameter families $`=\stackrel{~}{\mathrm{Gr}}_2(^n)`$ and for $`c0`$, the Kähler potentials extend to give non-singular metrics on $`T^{}`$, generalising the Eguchi-Hanson metrics on $`T^{}\mathrm{P}(1)`$. As $``$ is Hermitian symmetric, this is one of the cases considered by Biquard & Gauduchon . |
warning/0001/math0001057.html | ar5iv | text | # Random walks and electric networks
## Preface
Probability theory, like much of mathematics, is indebted to physics as a source of problems and intuition for solving these problems. Unfortunately, the level of abstraction of current mathematics often makes it difficult for anyone but an expert to appreciate this fact. In this work we will look at the interplay of physics and mathematics in terms of an example where the mathematics involved is at the college level. The example is the relation between elementary electric network theory and random walks.
Central to the work will be Polya’s beautiful theorem that a random walker on an infinite street network in $`d`$-dimensional space is bound to return to the starting point when $`d=2`$, but has a positive probability of escaping to infinity without returning to the starting point when $`d3`$. Our goal will be to interpret this theorem as a statement about electric networks, and then to prove the theorem using techniques from classical electrical theory. The techniques referred to go back to Lord Rayleigh, who introduced them in connection with an investigation of musical instruments. The analog of Polya’s theorem in this connection is that wind instruments are possible in our three-dimensional world, but are not possible in Flatland (Abbott ).
The connection between random walks and electric networks has been recognized for some time (see Kakutani , Kemeny, Snell, and Knapp , and Kelly ). As for Rayleigh’s method, the authors first learned it from Peter’s father Bill Doyle, who used it to explain a mysterious comment in Feller (, p. 425, Problem 14). This comment suggested that a random walk in two dimensions remains recurrent when some of the streets are blocked, and while this is ticklish to prove probabilistically, it is an easy consequence of Rayleigh’s method. The first person to apply Rayleigh’s method to random walks seems to have been Nash-Williams . Earlier, Royden had applied Rayleigh’s method to an equivalent problem. However, the true importance of Rayleigh’s method for probability theory is only now becoming appreciated. See, for example, Griffeath and Liggett , Lyons , and Kesten .
Here’s the plan of the work: In Section 1 we will restrict ourselves to the study of random walks on finite networks. Here we will establish the connection between the electrical concepts of current and voltage and corresponding descriptive quantities of random walks regarded as finite state Markov chains. In Section 2 we will consider random walks on infinite networks. Polya’s theorem will be proved using Rayleigh’s method, and the proof will be compared with the classical proof using probabilistic methods. We will then discuss walks on more general infinite graphs, and use Rayleigh’s method to derive certain extensions of Polya’s theorem. Certain of the results in Section 2 were obtained by Peter Doyle in work on his Ph.D. thesis.
To read this work, you should have a knowledge of the basic concepts of probability theory as well as a little electric network theory and linear algebra. An elementary introduction to finite Markov chains as presented by Kemeny, Snell, and Thompson would be helpful.
The work of Snell was carried out while enjoying the hospitality of Churchill College and the Cambridge Statistical Laboratory supported by an NSF Faculty Development Fellowship. He thanks Professors Kendall and Whittle for making this such an enjoyable and rewarding visit. Peter Doyle thanks his father for teaching him how to think like a physicist. We both thank Peter Ney for assigning the problem in Feller that started all this, David Griffeath for suggesting the example to be used in our first proof that 3-dimensional random walk is recurrent (Section 2.2.9), and Reese Prosser for keeping us going by his friendly and helpful hectoring. Special thanks are due Marie Slack, our secretary extraordinaire, for typing the original and the excessive number of revisions one is led to by computer formatting.
## 1 Random walks on finite networks
### 1.1 Random walks in one dimension
#### 1.1.1 A random walk along Madison Avenue
A *random walk*, or *drunkard’s walk*, was one of the first chance processes studied in probability; this chance process continues to play an important role in probability theory and its applications. An example of a random walk may be described as follows:
A man walks along a 5-block stretch of Madison Avenue. He starts at corner $`x`$ and, with probability 1/2, walks one block to the right and, with probability 1/2, walks one block to the left; when he comes to the next corner he again randomly chooses his direction along Madison Avenue. He continues until he reaches corner 5, which is home, or corner 0, which is a bar. If he reaches either home or the bar, he stays there. (See Figure 1.)
The problem we pose is to find the probability $`p(x)`$ that the man, starting at corner $`x`$, will reach home before reaching the bar. In looking at this problem, we will not be so much concerned with the particular form of the solution, which turns out to be $`p(x)=x/5`$, as with its general properties, which we will eventually describe by saying “$`p(x)`$ is the unique solution to a certain Dirichlet problem.”
#### 1.1.2 The same problem as a penny matching game
In another form, the problem is posed in terms of the following game: Peter and Paul match pennies; they have a total of 5 pennies; on each match, Peter wins one penny from Paul with probability 1/2 and loses one with probability 1/2; they play until Peter’s fortune reaches 0 (he is ruined) or reaches 5 (he wins all Paul’s money). Now $`p(x)`$ is the probability that Peter wins if he starts with $`x`$ pennies.
#### 1.1.3 The probability of winning: basic properties
Consider a random walk on the integers $`0,1,2,\mathrm{},N`$. Let $`p(x)`$ be the probability, starting at $`x`$, of reaching $`N`$ before 0. We regard $`p(x)`$ as a function defined on the points $`x=0,1,2,\mathrm{},N`$. The function $`p(x)`$ has the following properties:
(a) $`p(0)=0`$.
(b) $`p(N)=1`$.
(c) $`p(x)=\frac{1}{2}p(x1)+\frac{1}{2}p(x+1)`$ for $`x=1,2,\mathrm{},N1`$.
Properties (a) and (b) follow from our convention that 0 and N are traps; if the walker reaches one of these positions, he stops there; in the game interpretation, the game ends when one player has all of the pennies. Property (c) states that, for an interior point, the probability $`p(x)`$ of reaching home from $`x`$ is the average of the probabilities $`p(x1)`$ and $`p(x+1)`$ of reaching home from the points that the walker may go to from $`x`$. We can derive (c) from the following basic fact about probability:
Basic Fact. Let $`E`$ be any event, and $`F`$ and $`G`$ be events such that one and only one of the events $`F`$ or $`G`$ will occur. Then
$$𝐏(E)=𝐏(F)𝐏(E\text{ given }F)+𝐏(G)𝐏(E\text{ given }G).$$
In this case, let $`E`$ be the event “the walker ends at the bar”, $`F`$ the event “the first step is to the left”, and $`G`$ the event “the first step is to the right”. Then, if the walker starts at $`x`$, $`𝐏(E)=p(x)`$, $`𝐏(F)=𝐏(G)=\frac{1}{2}`$, $`𝐏(E\text{ given }F)=p(x1)`$, $`𝐏(E\text{ given }G)=p(x+1)`$, and (c) follows.
#### 1.1.4 An electric network problem: the same problem?
Let’s consider a second apparently very different problem. We connect equal resistors in series and put a unit voltage across the ends as in Figure 2.
Voltages $`v(x)`$ will be established at the points $`x=0,1,2,3,4,5`$. We have grounded the point $`x=0`$ so that $`v(0)=0`$. We ask for the voltage $`v(x)`$ at the points $`x`$ between the resistors. If we have $`N`$ resistors, we make $`v(0)=0`$ and $`v(N)=1`$, so $`v(x)`$ satisfies properties (a) and (b) of Section 1.1.3. We now show that $`v(x)`$ also satisfies (c).
By Kirchhoff’s Laws, the current flowing into $`x`$ must be equal to the current flowing out. By Ohm’s Law, if points $`x`$ and $`y`$ are connected by a resistance of magnitude $`R`$, then the current $`i_{xy}`$ that flows from $`x`$ to $`y`$ is equal to
$$i_{xy}=\frac{v(x)v(y)}{R}.$$
Thus for $`x=1,2,\mathrm{},N1`$,
$$\frac{v(x1)v(x)}{R}+\frac{v(x+1)v(x)}{R}=0.$$
Multiplying through by $`R`$ and solving for $`v(x)`$ gives
$$v(x)=\frac{v(x+1)+v(x1)}{2}$$
for $`x=1,2,\mathrm{},N1`$. Therefore, $`v(x)`$ also satisfies property (c).
We have seen that $`p(x)`$ and $`v(x)`$ both satisfy properties (a), (b), and (c) of Section 1.1.3. This raises the question: are $`p(x)`$ and $`v(x)`$ equal? For this simple example, we can easily find $`v(x)`$ using Ohm’s Law, find $`p(x)`$ using elementary probability, and see that they are the same. However, we want to illustrate a principle that will work for very general circuits. So instead we shall prove that these two functions are the same by showing that there is only one function that satisfies these properties, and we shall prove this by a method that will apply to more general situations than points connected together in a straight line.
###### Exercise 1.1.1
Referring to the random walk along Madison Avenue, let $`X=p(1)`$, $`Y=p(2)`$, $`Z=p(3)`$, and $`W=p(4)`$. Show that properties (a), (b), and (c) of Section 1.1.3 determine a set of four linear equations with variables $`X`$, $`Y`$, $`Z`$ and $`W`$. Show that these equations have a unique solution. What does this say about $`p(x)`$ and $`v(x)`$ for this special case?
###### Exercise 1.1.2
Assume that our walker has a tendency to drift in one direction: more specifically, assume that each step is to the right with probability $`p`$ or to the left with probability $`q=1p`$. Show that properties (a), (b), and (c) of Section 1.1.3 should be replaced by
(a) $`p(0)=0`$.
(b) $`p(N)=1`$.
(c) $`p(x)=qp(x1)+pp(x+1)`$.
###### Exercise 1.1.3
In our electric network problem, assume that the resistors are not necessarily equal. Let $`R_x`$ be the resistance between $`x`$ and $`x+1`$. Show that
$$v(x)=\frac{\frac{1}{R_{x1}}}{\frac{1}{R_{x1}}+\frac{1}{R_x}}v(x1)+\frac{\frac{1}{R_x}}{\frac{1}{R_{x1}}+\frac{1}{R_x}}v(x+1).$$
How should the resistors be chosen to correspond to the random walk of Exercise 1.1.2?
#### 1.1.5 Harmonic functions in one dimension; the Uniqueness Principle
Let $`S`$ be the set of points $`S=\{0,1,2,\mathrm{},N\}`$. We call the points of the set $`D=\{1,2,\mathrm{},N1\}`$ the *interior points* of $`S`$ and those of $`B=\{0,N\}`$ the *boundary points* of $`S`$. A function $`f(x)`$ defined on $`S`$ is *harmonic* if, at points of $`D`$, it satisfies the averaging property
$$f(x)=\frac{f(x1)+f(x+1)}{2}.$$
As we have seen, $`p(x)`$ and $`v(x)`$ are harmonic functions on $`S`$ having the same values on the boundary: $`p(0)=v(0)=0`$; $`p(N)=v(N)=1`$. Thus both $`p(x)`$ and $`v(x)`$ solve the problem of finding a harmonic function having these boundary values. Now the problem of finding a harmonic function given its boundary values is called the *Dirichlet problem*, and the *Uniqueness Principle* for the Dirichlet problem asserts that there cannot be two different harmonic functions having the same boundary values. In particular, it follows that $`p(x)`$ and $`v(x)`$ are really the same function, and this is what we have been hoping to show. Thus the fact that $`p(x)=v(x)`$ is an aspect of a general fact about harmonic functions.
We will approach the Uniqueness Principle by way of the *Maximum Principle* for harmonic functions, which bears the same relation to the Uniqueness Principle as Rolle’s Theorem does to the Mean Value Theorem of Calculus.
Maximum Principle . A harmonic function $`f(x)`$ defined on $`S`$ takes on its maximum value $`M`$ and its minimum value $`m`$ on the boundary.
Proof. Let $`M`$ be the largest value of $`f`$. Then if $`f(x)=M`$ for $`x`$ in $`D`$, the same must be true for $`f(x1)`$ and $`f(x+1)`$ since $`f(x)`$ is the average of these two values. If $`x1`$ is still an interior point, the same argument implies that $`f(x2)=M`$; continuing in this way, we eventually conclude that $`f(0)=M`$. That same argument works for the minimum value $`m`$. $`\mathrm{}`$
Uniqueness Principle. If $`f(x)`$ and $`g(x)`$ are harmonic functions on $`S`$ such that $`f(x)=g(x)`$ on $`B`$, then $`f(x)=g(x)`$ for all $`x`$.
Proof. Let $`h(x)=f(x)g(x)`$. Then if $`x`$ is any interior point,
$$\frac{h(x1)+h(x+1)}{2}=\frac{f(x1)+f(x+1)}{2}\frac{g(x1)+g(x+1)}{2},$$
and $`h`$ is harmonic. But $`h(x)=0`$ for $`x`$ in $`B`$, and hence, by the Maximum Principle, the maximum and mininium values of $`h`$ are 0. Thus $`h(x)=0`$ for all $`x`$, and $`f(x)=g(x)`$ for all $`x`$. $`\mathrm{}`$
Thus we finally prove that $`p(x)=v(x)`$; but what does $`v(x)`$ equal? The Uniqueness Principle shows us a way to find a concrete answer: just guess. For if we can find any harmonic function $`f(x)`$ having the right boundary values, the Uniqueness Principle guarantees that
$$p(x)=v(x)=f(x).$$
The simplest function to try for $`f(x)`$ would be a linear function; this leads to the solution $`f(x)=x/N`$. Note that $`f(0)=0`$ and $`f(N)=1`$ and
$$\frac{f(x1)+f(x+1)}{2}=\frac{x1+x+1}{2N}=\frac{x}{N}=f(x).$$
Therefore $`f(x)=p(x)=v(x)=x/N`$.
As another application of the Uniqueness Principle, we prove that our walker will eventually reach 0 or $`N`$. Choose a starting point $`x`$ with $`0<x<N`$. Let $`h(x)`$ be the probability that the walker never reaches the boundary $`B=\{0,N\}`$. Then
$$h(x)=\frac{1}{2}h(x+1)+\frac{1}{2}h(x1)$$
and h is harmonic. Also $`h(0)=h(N)=0`$; thus, by the Maximum Principle, $`h(x)=0`$ for all $`x`$.
###### Exercise 1.1.4
Show that you can choose $`A`$ and $`B`$ so that the function $`f(x)=A(q/p)^x+B`$ satisfies the modified properties (a), (b) and (c) of Exercise 1.1.2. Does this show that $`f(x)=p(x)`$?
###### Exercise 1.1.5
Let $`m(x)`$ be the expected number of steps, starting at $`x`$, required to reach 0 or $`N`$ for the first time. It can be proven that $`m(x)`$ is finite. Show that $`m(x)`$ satisfies the conditions
(a) $`m(0)=0`$.
(b) $`m(N)=0`$.
(c) $`m(x)=\frac{1}{2}m(x+1)+\frac{1}{2}m(x1)+1`$.
###### Exercise 1.1.6
Show that the conditions in Exercise 1.1.5 have a unique solution. Hint: show that if $`m`$ and $`\overline{m}`$ are two solutions, then $`f=m\overline{m}`$ is harmonic with $`f(0)=f(N)=0`$ and hence $`f(x)=0`$ for all $`x`$.
###### Exercise 1.1.7
Show that you can choose $`A`$, $`B`$, and $`C`$ such that $`f(x)=A+Bx+Cx^2`$ satisfies all the conditions of Exercise 1.1.5. Does this show that $`f(x)=m(x)`$ for this choice of $`A`$, $`B`$, and $`C`$?
###### Exercise 1.1.8
Find the expected duration of the walk down Madison Avenue as a function of the walker’s starting point (1, 2, 3, or 4).
#### 1.1.6 The solution as a fair game (martingale)
Let us return to our interpretation of a random walk as Peter’s fortune in a game of penny matching with Paul. On each match, Peter wins one penny with probability 1/2 and loses one penny with probability 1/2. Thus, when Peter has $`k`$ pennies his expected fortune after the next play is
$$\frac{1}{2}(k1)+\frac{1}{2}(k+1)=k,$$
so his expected fortune after the next play is equal to his present fortune. This says that he is playing a *fair game*; a chance process that can be interpreted as a player’s fortune in a fair game is called a *martingale*.
Now assume that Peter and Paul have a total of $`N`$ pennies. Let $`p(x)`$ be the probability that, when Peter has $`x`$ pennies, he will end up with all $`N`$ pennies. Then Peter’s expected final fortune in this game is
$$(1p(x))0+p(x)N=p(x)N.$$
If we could be sure that a fair game remains fair to the end of the game, then we could conclude that Peter’s expected final fortune is equal to his starting fortune $`x`$, i.e., $`x=p(x)N`$. This would give $`p(x)=x/N`$ and we would have found the probability that Peter wins using the fact that a fair game remains fair to the end. Note that the time the game ends is a random time, namely, the time that the walk first reaches 0 or $`N`$ for the first time. Thus the question is, is the fairness of a game preserved when we stop at a random time?
Unfortunately, this is not always the case. To begin with, if Peter somehow has knowledge of what the future holds in store for him, he can decide to quit when he gets to the end of a winning streak. But even if we restrict ourselves to stopping rules where the decision to stop or continue is independent of future events, fairness may not be preserved. For example, assume that Peter is allowed to go into debt and can play as long as he wants to. He starts with 0 pennies and decides to play until his fortune is 1 and then quit. We shall see that a random walk on the set of all integers, starting at 0, will reach the point 1 if we wait long enough. Hence, Peter will end up one penny ahead by this system of stopping.
However, there are certain conditions under which we can guarantee that a fair game remains fair when stopped at a random time. For our purposes, the following standard result of martingale theory will do:
Martingale Stopping Theorem. A fair game that is stopped at a random time will remain fair to the end of the game if it is assumed that there is a finite amount of money in the world and a player must stop if he wins all this money or goes into debt by this amount.
This theorem would justify the above argument to obtain $`p(x)=x/N`$.
Let’s step back and see how this martingale argument worked. We began with a harmonic function, the function $`f(x)=x`$, and interpreted it as the player’s fortune in a fair game. We then considered the player’s expected final fortune in this game. This was another harmonic function having the same boundary values and we appealed to the Martingale Stopping Theorem to argue that this function must be the same as the original function. This allowed us to write down an expression for the probability of winning, which was what we were looking for.
Lurking behind this argument is a general principle: If we are given boundary values of a function, we can come up with a harmonic function having these boundary values by assigning to each point the player’s expected final fortune in a game where the player starts from the given point and carries out a random walk until he reaches a boundary point, where he receives the specified payoff. Furthermore, the Martingale Stopping Theorern allows us to conclude that there can be no other harmonic function with these boundary values. Thus martingale theory allows us to establish existence and uniqueness of solutions to a Dirichlet problem. All this isn’t very exciting for the cases we’ve been considering, but the nice thing is that the same arguments carry through to the more general situations that we will be considering later on.
The study of martingales was originated by Levy and Ville . Kakutani showed the connection between random walks and harmonic functions. Doob developed martingale stopping theorems and showed how to exploit the preservation of fairness to solve a wide variety of problems in probability theory. An informal discussion of martingales may be found in Snell .
###### Exercise 1.1.9
Consider a random walk with a drift; that is, there is a probability $`p\frac{1}{2}`$ of going one step to the right and a probability $`q=1p`$ of going one step to the left. (See Exercise 1.1.2.) Let $`w(x)=(q/p)^x`$; show that, if you interpret $`w(x)`$ as your fortune when you are at $`x`$, the resulting game is fair. Then use the Martingale Stopping Theorem to argue that
$$w(x)=p(x)w(N)+(1p(x))w(0).$$
Solve for $`p(x)`$ to obtain
$$p(x)=\frac{\left(\frac{q}{p}\right)^x1}{\left(\frac{q}{p}\right)^N1}.$$
###### Exercise 1.1.10
You are gambling against a professional gambler; you start with $`A`$ dollars and the gambler with $`B`$ dollars; you play a game in which you win one dollar with probability $`p<\frac{1}{2}`$ and lose one dollar with probability $`q=1p`$; play continues until you or the gambler runs out of money. Let $`R_A`$ be the probability that you are ruined. Use the result of Exercise 1.1.9 to show that
$$R_A=\frac{1\left(\frac{p}{q}\right)^B}{1\left(\frac{p}{q}\right)^N}$$
with $`N=A+B`$. If you start with 20 dollars and the gambler with 50 dollars and $`p=.45`$, find the probability of being ruined.
###### Exercise 1.1.11
The gambler realizes that the probability of ruining you is at least $`1(p/q)^B`$ (Why?). The gambler wants to make the probability at least .999. For this, $`(p/q)^B`$ should be at most .001. If the gambler offers you a game with $`p=.499`$, how large a stake should she have?
### 1.2 Random walks in two dimensions
#### 1.2.1 An example
We turn now to the more complicated problem of a random walk on a two-dimensional array. In Figure 3 we illustrate such a walk.
The large dots represent boundary points; those marked $`E`$ indicate escape routes and those marked $`P`$ are police. We wish to find the probability $`p(x)`$ that our walker, starting at an interior point $`x`$, will reach an escape route before he reaches a policeman. The walker moves from $`x=(a,b)`$ to each of the four neighboring points $`(a+1,b)`$, $`(a1,b)`$, $`(a,b+1)`$, $`(a,b1)`$ with probability $`\frac{1}{4}`$. If he reaches a boundary point, he remains at this point.
The corresponding voltage problem is shown in Figure 4.
The boundary points $`P`$ are grounded and points $`E`$ are connected and fixed at one volt by a one-volt battery. We ask for the voltage $`v(x)`$ at the interior points.
#### 1.2.2 Harmonic functions in two dimensions
We now define harmonic functions for sets of lattice points in the plane (a lattice point is a point with integer coordinates). Let $`S=DB`$ be a finite set of lattice points such that (a) $`D`$ and $`B`$ have no points in common, (b) every point of $`D`$ has its four neighboring points in $`S`$, and (c) every point of $`B`$ has at least one of its four neighboring points in $`D`$. We assume further that $`S`$ hangs together in a nice way, namely, that for any two points $`P`$ and $`Q`$ in $`S`$, there is a sequence of points $`P_j`$ in $`D`$ such that $`P,P_1,P_2,\mathrm{},P_n,Q`$ forms a path from $`P`$ to $`A`$. We call the points of $`D`$ the *interior points* of $`S`$ and the points of $`B`$ the *boundary points* of $`S`$.
A function $`f`$ defined on $`S`$ is *harmonic* if, for points $`(a,b)`$ in $`D`$, it has the averaging property
$$f(a,b)=\frac{f(a+1,b)+f(a1,b)+f(a,b+1)+f(a,b1)}{4}.$$
Note that there is no restriction on the values of $`f`$ at the boundary points.
We would like to prove that $`p(x)=v(x)`$ as we did in the one-dimensional case. That $`p(x)`$ is harmonic follows again by considering all four possible first steps; that $`v(x)`$ is harmonic follows again by Kirchhoff’s Laws since the current coming into $`x=(a,b)`$ is
$$\frac{v(a+1,b)v(a,b)}{R}+\frac{v(a1,b)v(a,b)}{R}+\frac{v(a,b+1)v(a,b)}{R}+\frac{v(a,b1)v(a,b)}{R}=0.$$
Multiplying through by $`R`$ and solving for $`v(a,b)`$ gives
$$v(a,b)=\frac{v(a+1,b)+v(a1,b)+v(a,b+1)+v(a,b1)}{4}.$$
Thus $`p(x)`$ and $`v(x)`$ are harmonic functions with the same boundary values. To show from this that they are the same, we must extend the Uniqueness Principle to two dimensions.
We first prove the Maximum Principle. If $`M`$ is the maximum value of $`f`$ and if $`f(P)=M`$ for $`P`$ an interior point, then since $`f(P)`$ is the average of the values of $`f`$ at its neighbors, these values must all equal $`M`$ also. By working our way due south, say, repeating this argument at every step, we eventually reach a boundary point $`Q`$ for which we can conclude that $`f(Q)=M`$. Thus a harmonic function always attains its maximum (or minimum) on the boundary; this is the Maximum Principle. The proof of the Uniqueness Principle goes through as before since again the difference of two harmonic functions is harmonic.
The fair game argument, using the Martingale Stopping Theorem, holds equally well and again gives an alternative proof of the existence and uniqueness to the solution of the Dirichlet problem.
###### Exercise 1.2.1
Show that if $`f`$ and $`g`$ are harmonic functions so is $`h=af+bg`$ for constants $`a`$ and $`b`$. This is called the *superposition principle*.
###### Exercise 1.2.2
Let $`B_1,B_2,\mathrm{},B_n`$ be the boundary points for a region $`S`$. Let $`e_j(a,b)`$ be a function that is harmonic in $`S`$ and has boundary value 1 at $`B_j`$ and 0 at the other boundary points. Show that if arbitrary boundary values $`v_1,v_2,\mathrm{},v_n`$ are assigned, we can find the harmonic function $`v`$ with these values from the solutions $`e_1,e_2,\mathrm{},e_n`$.
#### 1.2.3 The Monte Carlo solution
Finding the exact solution to a Dirichlet problem in two dimensions is not always a simple matter, so before taking on this problem, we will consider two methods for generating approximate solutions. In this section we will present a method using random walks. This method is known as a *Monte Carlo method*, since random walks are random, and gambling involves randomness, and there is a famous gambling casino in Monte Carlo. In Section 1.2.4, we will describe a much more effective method for finding approximate solutions, called the *method of relaxations*.
We have seen that the solution to the Dirichlet problem can be found by finding the value of a player’s final winning in the following game: Starting at $`x`$ the player carries out a random walk until reaching a boundary point. He is then paid an amount $`f(y)`$ if $`y`$ is the boundary point first reached. Thus to find $`f(x)`$, we can start many random walks at $`x`$ and find the average final winnings for these walks. By the law of averages (the law of large numbers in probability theory), the estimate that we obtain this way will approach the true expected final winning $`f(x)`$.
Here are some estimates obtained this way by starting 10,000 random walks from each of the interior points and, for each $`x`$, estimating $`f(x)`$ by the average winning of the random walkers who started at this point.
$$\begin{array}{cccc}& & 1& 1\\ & 1.824& .785& 1\\ 1& .876& .503& .317& 0\\ & 1& 0& 0\end{array}$$
This method is a colorful way to solve the problem, but quite inefficient. We can use probability theory to estimate how inefficient it is. We consider the case with boundary values I or 0 as in our example. In this case, the expected final winning is just the probability that the walk ends up at a boundary point with value 1. For each point $`x`$, assume that we carry out $`n`$ random walks; we regard each random walk to be an experiment and interpret the outcome of the $`i`$th experiment to be a “success” if the walker ends at a boundary point with a 1 and a “failure” otherwise. Let $`p=p(x)`$ be the unknown probability for success for a walker starting at $`x`$ and $`q=1p`$. How many walks should we carry out to get a reasonable estimate for $`p`$? We estimate $`p`$ to be the fraction $`\overline{p}`$ of the walkers that end at a 1.
We are in the position of a pollster who wishes to estimate the proportion $`p`$ of people in the country who favor candidate $`A`$ over $`B`$. The pollster chooses a random sample of $`n`$ people and estimates $`p`$ as the proportion $`\overline{p}`$ of voters in his sample who favor $`A`$. (This is a gross oversimplification of what a pollster does, of course.) To estimate the number $`n`$ required, we can use the central limit theorem. This theorem states that, if $`S_n`$, is the number of successes in $`n`$ independent experiments, each having probability $`p`$ for success, then for any $`k>0`$
$$𝐏\left(k<\frac{S_nnp}{\sqrt{npq}}<k\right)A(k),$$
where $`A(k)`$ is the area under the normal curve between $`k`$ and $`k`$. For $`k=2`$ this area is approximately .95; what does this say about $`\overline{p}=S_n/n`$? Doing a little rearranging, we see that
$$𝐏\left(2<\frac{\overline{p}p}{\sqrt{\frac{pq}{n}}}<2\right).95$$
or
$$𝐏\left(2\frac{\sqrt{pq}}{n}<\overline{p}p<2\frac{\sqrt{pq}}{n}\right).95.$$
Since $`\sqrt{pq}\frac{1}{2}`$,
$$𝐏\left(\frac{1}{\sqrt{n}}<\overline{p}p<\frac{1}{\sqrt{n}}\right)\stackrel{>}{}.95.$$
Thus, if we choose $`\frac{1}{\sqrt{n}}=.01`$, or $`n=10,000`$, there is a 95 percent chance that our estimate $`\overline{p}=S_n/n`$ will not be off by more than .01. This is a large number for rather modest accuracy; in our example we carried out 10,000 walks from each point and this required about 5 seconds on the Dartmouth computer. We shall see later, when we obtain an exact solution, that we did obtain the accuracy predicted.
###### Exercise 1.2.3
You play a game in which you start a random walk at the center in the grid shown in Figure 5.
When the walk reaches the boundary, you receive a payment of $`+1`$ or $`1`$ as indicated at the boundary points. You wish to simulate this game to see if it is a favorable game to play; how many simulations would you need to be reasonably certain of the value of this game to an accuracy of .01? Carry out such a simulation and see if you feel that it is a favorable game.
#### 1.2.4 The original Dirichlet problem; the method of relaxations
The Dirichlet problem we have been studying is not the original Dirichlet problem, but a discrete version of it. The original Dirichlet problem concerns the distribution of temperature, say, in a continuous medium; the following is a representative example.
Suppose we have a thin sheet of metal gotten by cutting out a small square from the center of a large square. The inner boundary is kept at temperature 0 and the outer boundary is kept at temperature 1 as indicated in Figure 6.
The problem is to find the temperature at points in the rectangle’s interior. If $`u(x,y)`$ is the temperature at $`(x,y)`$, then $`u`$ satisfies Laplace’s differential equation
$$u_{xx}+u_{yy}=0.$$
A function that satisfies this differential equation is called *harmonic*. It has the property that the value $`u(x,y)`$ is equal to the average of the values over any circle with center $`(x,y)`$ lying inside the region. Thus to determine the temperature $`u(x,y)`$, we must find a harmonic function defined in the rectangle that takes on the prescribed boundary values. We have a problem entirely analogous to our discrete Dirichlet problem, but with continuous domain.
The *method of relaxations* was introduced as a way to get approximate solutions to the original Dirichlet problem. This method is actually more closely connected to the discrete Dirichlet problem than to the continuous problem. Why? Because, faced with the continuous problem just described, no physicist will hesitate to replace it with an analogous discrete problem, approximating the continuous medium by an array of lattice points such as that depicted in Figure 7, and searching for a function that is harmonic in our discrete sense and that takes on the appropriate boundary values. It is this approximating discrete problem to which the method of relaxations applies.
Here’s how the method goes. Recall that we are looking for a function that has specified boundary values, for which the value at any interior point is the average of the values at its neighbors. Begin with any function having the specified boundary values, pick an interior point, and see what is happening there. In general, the value of the function at the point we are looking at will not be equal to the average of the values at its neighbors. So adjust the value of the function to be equal to the average of the values at its neighbors. Now run through the rest of the interior points, repeating this process. When you have adjusted the values at all of the interior points, the function that results will not be harmonic, because most of the time after adjusting the value at a point to be the average value at its neighbors, we afterwards came along and adjusted the values at one or more of those neighbors, thus destroying the harmony. However, the function that results after running through all the interior points, if not harmonic, is more nearly harmonic than the function we started with; if we keep repeating this averaging process, running through all of the interior points again and again, the function will approximate more and more closely the solution to our Dirichlet problem.
We do not yet have the tools to prove that this method works for a general initial guess; this will have to wait until later (see Exercise 1.3.12). We will start with a special choice of initial values for which we can prove that the method works (see Exercise 1.2.5).
We start with all interior points 0 and keep the boundary points fixed.
$$\begin{array}{cccc}& & 1& 1\\ & 1& 0& 0& 1\\ 1& 0& 0& 0& 0\\ & 1& 0& 0\end{array}$$
After one iteration we have:
$$\begin{array}{cccc}& & 1& 1\\ & 1& .547& .648& 1\\ 1& .75& .188& .047& 0\\ & 1& 0& 0\end{array}$$
Note that we go from left to right moving up each column replacing each value by the average of the four neighboring values. The computations for this first iteration are
$$.75=(1/4)(1+1+1+0)$$
$$.1875=(1/4)(.75+0+0+0)$$
$$.5469=(1/4)(.1875+1+1+0)$$
$$.0469=(1/4)(.1875+0+0+0)$$
$$.64844=(1/4)(.0469+.5769+1+1)$$
We have printed the results to three decimal places. We continue the iteration until we obtain the same results to three decimal places. This occurs at iterations 8 and 9. Here’s what we get:
$$\begin{array}{cccc}& & 1& 1\\ & 1& .823& .787& 1\\ 1& .876& .506& .323& 0\\ & 1& 0& 0\end{array}$$
We see that we obtain the same result to three places after only nine iterations and this took only a fraction of a second of computing time. We shall see that these results are correct to three place accuracy. Our Monte Carlo method took several seconds of computing time and did not even give three place accuracy.
The classical reference for the method of relaxations as a means of finding approximate solutions to continuous problems is Courant, Friedrichs, and Lewy . For more information on the relationship between the original Dirichlet problem and the discrete analog, see Hersh and Griego .
###### Exercise 1.2.4
Apply the method of relaxations to the discrete problem illustrated in Figure 7.
###### Exercise 1.2.5
Consider the method of relaxations started with an initial guess with the property that the value at each point is $``$ the average of the values at the neighbors of this point. Show that the successive values at a point $`u`$ are monotone increasing with a limit $`f(u)`$ and that these limits provide a solution to the Dirichlet problem.
#### 1.2.5 Solution by solving linear equations
In this section we will show how to find an exact solution to a two-dimensional Dirichlet problem by solving a system of linear equations. As usual, we will illustrate the method in the case of the example introduced in Section 1.2.1. This example is shown again in Figure 8; the interior points have been labelled $`a`$, $`b`$, $`c`$, $`d`$, and $`e`$.
By our averaging property, we have
$$x_a=\frac{x_b+x_d+2}{4}$$
$$x_b=\frac{x_a+x_c+2}{4}$$
$$x_c=\frac{x_d+3}{4}$$
$$x_d=\frac{x_a+x_c+x_e}{4}$$
$$x_e=\frac{x_b+x_d}{4}.$$
We can rewrite these equations in matrix form as
$$\left(\begin{array}{ccccc}1& 1/4& 0& 1/4& 0\\ 1/4& 1& 0& 0& 1/4\\ 0& 0& 1& 1/4& 0\\ 1/4& 0& 1/4& 1& 1/4\\ 0& 1/4& 0& 1/4& 1\end{array}\right)\left(\begin{array}{c}x_a\\ x_b\\ x_c\\ x_d\\ x_e\end{array}\right)=\left(\begin{array}{c}1/2\\ 1/2\\ 3/4\\ 0\\ 0\end{array}\right).$$
We can write this in symbols as
$$\mathrm{𝐀𝐱}=𝐮.$$
Since we know there is a unique solution, $`𝐀`$ must have an inverse and
$$𝐱=𝐀^1𝐮.$$
Carrying out this calculation we find
$$\text{Calculated }𝐱=\left(\begin{array}{c}.823\\ .787\\ .876\\ .506\\ .323\end{array}\right).$$
Here, for comparison, are the approximate solutions found earlier:
$$\text{Monte Carlo }𝐱=\left(\begin{array}{c}.824\\ .785\\ .876\\ .503\\ .317\end{array}\right).$$
$$\text{Relaxed }𝐱=\left(\begin{array}{c}.823\\ .787\\ .876\\ .506\\ .323\end{array}\right).$$
We see that our Monte Carlo approximations were fairly good in that no error of the simulation is greater than .01, and our relaxed approximations were very good indeed, in that the error does not show up at all.
###### Exercise 1.2.6
Consider a random walker on the graph of Figure 9.
Find the probability of reaching the point with a 1 before any of the points with 0’s for each starting point $`a,b,c,d`$.
###### Exercise 1.2.7
Solve the discrete Dirichlet problem for the graph shown in Figure 10.
The interior points are $`a,b,c,d`$. (Hint: See Exercise 1.2.2.)
###### Exercise 1.2.8
Find the exact value, for each possible starting point, for the game described in Exercise 1.2.3. Is the game favorable starting in the center?
#### 1.2.6 Solution by the method of Markov chains
In this section, we describe how the Dirichlet problem can be solved by the method of Markov chains. This method may be viewed as a more sophisticated version of the method of linear equations.
A *finite Markov chain* is a special type of chance process that may be described informally as follows: we have a set $`S=\{s_1,s_2,\mathrm{},s_r\}`$ of *states* and a chance process that moves around through these states. When the process is in state $`s_i`$, it moves with probability $`P_{ij}`$ to the state $`s_j`$. The transition probabilities $`P_{ij}`$ are represented by an $`r`$-by-$`r`$ matrix $`𝐏`$ called the *transition matrix*. To specify the chance process completely we must give, in addition to the transition matrix, a method for starting the process. We do this by specifying a specific state in which the process starts.
According to Kemeny, Snell, and Thompson , in the Land of Oz, there are three kinds of weather: rain, nice, and snow. There are never two nice days in a row. When it rains or snows, half the time it is the same the next day. If the weather changes, the chances are equal for a change to each of the other two types of weather. We regard the weather in the Land of Oz as a Markov chain with transition matrix:
𝐏=(
RNSR/12/14/14N/120/12S/14/14/12
).𝐏
RNSR/12/14/14N/120/12S/14/14/12
{\mathbf{P}}=\hbox{}\;\vbox{\kern 49.83331pt\hbox{$\kern 90.83315pt\kern-8.75pt\left(\kern-90.83315pt\vbox{\kern-49.83331pt\vbox{\halign{$#$\hfil\kern 2\p@\kern\@tempdima&\thinspace\hfil$#$\hfil&&\quad\hfil$#$\hfil\cr\hfil\crcr\kern-12.0pt\cr$\hfil\kern 2.0pt\kern 8.75pt&\mbox{R}&\mbox{N}&\mbox{S}\crcr\kern 2.0pt\cr\mbox{R}$\hfil\kern 2.0pt\kern 8.75pt&1/2&1/4&1/4\cr\mbox{N}$\hfil\kern 2.0pt\kern 8.75pt&1/2&0&1/2\cr\mbox{S}$\hfil\kern 2.0pt\kern 8.75pt&1/4&1/4&1/2\crcr\cr}}\kern-12.0pt}\,\right)$}}.
When we start in a particular state, it is natural to ask for the probability that the process is in each of the possible states after a specific number of steps. In the study of Markov chains, it is shown that this information is provided by the powers of the transition matrix. Specifically, if $`𝐏^n`$ is the matrix $`𝐏`$ raised to the $`n`$th power, the entries $`P_{ij}^n`$ represent the probability that the chain, started in state $`s_i`$, will, after $`n`$ steps, be in state $`s_j`$. For example, the fourth power of the transition matrix $`𝐏`$ for the weather in the Land of Oz is
$$𝐏^4=\left(\text{RNSR.402.199.398N.398.203.398S.398.199.402}\right).$$
Thus, if it is raining today in the Land of Oz, the probability that the weather will be nice four days from now is .199. Note that the probability of a particular type of weather four days from today is essentially independent of the type of weather today. This Markov chain is an example of a type of chain called a regular chain. A Markov chain is a *regular* chain if some power of the transition matrix has no zeros. In the study of regular Markov chains, it is shown that the probability of being in a state after a large number of steps is independent of the starting state.
As a second example, we consider a random walk in one dimension. Let us assume that the walk is stopped when it reaches either state 0 or 4. (We could use 5 instead of 4, as before, but we want to keep the matrices small.) We can regard this random walk as a Markov chain with states 0, 1, 2, 3, 4 and transition matrix given by
𝐏=(
012340100001/120/120020/120/120300/120/12400001
).𝐏
012340100001/120/120020/120/120300/120/12400001
{\mathbf{P}}=\hbox{}\;\vbox{\kern 62.33333pt\hbox{$\kern 147.22194pt\kern-8.75pt\left(\kern-147.22194pt\vbox{\kern-62.33333pt\vbox{\halign{$#$\hfil\kern 2\p@\kern\@tempdima&\thinspace\hfil$#$\hfil&&\quad\hfil$#$\hfil\cr\hfil\crcr\kern-12.0pt\cr$\hfil\kern 2.0pt\kern 8.75pt&0&1&2&3&4\crcr\kern 2.0pt\cr 0$\hfil\kern 2.0pt\kern 8.75pt&1&0&0&0&0\cr 1$\hfil\kern 2.0pt\kern 8.75pt&1/2&0&1/2&0&0\cr 2$\hfil\kern 2.0pt\kern 8.75pt&0&1/2&0&1/2&0\cr 3$\hfil\kern 2.0pt\kern 8.75pt&0&0&1/2&0&1/2\cr 4$\hfil\kern 2.0pt\kern 8.75pt&0&0&0&0&1\crcr\cr}}\kern-12.0pt}\,\right)$}}.
The states 0 and 4 are *traps* or *absorbing states*. These are states that, once entered, cannot be left. A Markov chain is called *absorbing* if it has at least one absorbing state and if, from any state, it is possible (not necessarily in one step) to reach at least one absorbing state. Our Markov chain has this property and so is an absorbing Markov chain. The states of an absorbing chain that are not traps are called *non-absorbing*.
When an absorbing Markov chain is started in a non-absorbing state, it will eventually end up in an absorbing state. For non-absorbing state $`s_i`$ and absorbing state $`s_j`$, we denote by $`B_{ij}`$ the probability that the chain starting in $`s_i`$ will end up in state $`s_j`$. We denote by $`𝐁`$ the matrix with entries $`B_{ij}`$. This matrix will have as many rows as non-absorbing states and as many columns as there are absorbing states. For our random walk example, the entries $`B_{x,4}`$ will give the probability that our random walker, starting at $`x`$, will reach 4 before reaching 0. Thus, if we can find the matrix $`𝐁`$ by Markov chain techniques, we will have a way to solve the Dirichlet problem.
We shall show, in fact, that the Dirichlet problem has a natural generalization in the context of absorbing Markov chains and can be solved by Markov chain methods.
Assume now that $`𝐏`$ is an absorbing Markov chain and that there are $`u`$ absorbing states and $`v`$ non-absorbing states. We reorder the states so that the absorbing states come first and the non-absorbing states come last. Then our transition matrix has the canonical form:
$$𝐏=\left(\begin{array}{cc}𝐈& \mathrm{𝟎}\\ 𝐑& 𝐐\end{array}\right).$$
Here $`𝐈`$ is a $`u`$-by-$`u`$ identity matrix; $`\mathrm{𝟎}`$ is a matrix of dimension $`u`$-by-$`v`$ with all entries 0.
For our random walk example this canonical form is:
(
041230100004010001/1200/120200/120/1230/120/120
).
041230100004010001/1200/120200/120/1230/120/120
\hbox{}\;\vbox{\kern 62.33333pt\hbox{$\kern 147.22194pt\kern-8.75pt\left(\kern-147.22194pt\vbox{\kern-62.33333pt\vbox{\halign{$#$\hfil\kern 2\p@\kern\@tempdima&\thinspace\hfil$#$\hfil&&\quad\hfil$#$\hfil\cr\hfil\crcr\kern-12.0pt\cr$\hfil\kern 2.0pt\kern 8.75pt&0&4&1&2&3\crcr\kern 2.0pt\cr 0$\hfil\kern 2.0pt\kern 8.75pt&1&0&0&0&0\cr 4$\hfil\kern 2.0pt\kern 8.75pt&0&1&0&0&0\cr 1$\hfil\kern 2.0pt\kern 8.75pt&1/2&0&0&1/2&0\cr 2$\hfil\kern 2.0pt\kern 8.75pt&0&0&1/2&0&1/2\cr 3$\hfil\kern 2.0pt\kern 8.75pt&0&1/2&0&1/2&0\crcr\cr}}\kern-12.0pt}\,\right)$}}.
The matrix $`𝐍=(𝐈𝐐)^1`$ is called the *fundamental matrix* for the absorbing chain $`𝐏`$. The entries $`N_{ij}`$ of this matrix have the following probabilistic interpretation: $`N_{ij}`$ is the expected number of times that the chain will be in state $`s_j`$ before absorption when it is started in $`s_i`$. (To see why this is true, think of how $`(𝐈𝐐)^1`$ would look if it were written as a geometric series.) Let $`\mathrm{𝟏}`$ be a column vector of all 1’s. Then the vector $`𝐭=\mathrm{𝐍𝐈}`$ gives the expected number of steps before absorption for each starting state.
The absorption probabilities $`𝐁`$ are obtained from $`𝐍`$ by the matrix formula
$$𝐁=(𝐈𝐐)^1𝐑.$$
This simply says that to get the probability of ending up at a given absorbing state, we add up the probabilities of going there from all the non-absorbing states, weighted by the number of times we expect to be in those (non-absorbing) states.
For our random walk example
$$𝐐=\left(\begin{array}{ccc}0& \frac{1}{2}& 0\\ \frac{1}{2}& 0& \frac{1}{2}\\ 0& \frac{1}{2}& 0\end{array}\right)$$
$$𝐈𝐐=\left(\begin{array}{ccc}1& \frac{1}{2}& 0\\ \frac{1}{2}& 1& \frac{1}{2}\\ 0& \frac{1}{2}& 1\end{array}\right)$$
𝐍=(𝐈𝐐)1=(
1231321122121312132
)𝐍superscript𝐈𝐐1
1231321122121312132
{\mathbf{N}}=({\mathbf{I}}-{\mathbf{Q}})^{-1}=\hbox{}\;\vbox{\kern 43.93332pt\hbox{$\kern 45.00005pt\kern-8.75pt\left(\kern-45.00005pt\vbox{\kern-43.93332pt\vbox{\halign{$#$\hfil\kern 2\p@\kern\@tempdima&\thinspace\hfil$#$\hfil&&\quad\hfil$#$\hfil\cr\hfil\crcr\kern-12.0pt\cr$\hfil\kern 2.0pt\kern 8.75pt&1&2&3\crcr\kern 2.0pt\cr 1$\hfil\kern 2.0pt\kern 8.75pt&\frac{3}{2}&1&{\frac{1}{2}}\cr 2$\hfil\kern 2.0pt\kern 8.75pt&1&2&1\cr 3$\hfil\kern 2.0pt\kern 8.75pt&{\frac{1}{2}}&1&\frac{3}{2}\crcr\cr}}\kern-12.0pt}\,\right)$}}
$$𝐭=\mathrm{𝐍𝟏}=\left(\begin{array}{ccc}\frac{3}{2}& 1& \frac{1}{2}\\ 1& 2& 1\\ \frac{1}{2}& 1& \frac{3}{2}\end{array}\right)\left(\begin{array}{c}1\\ 1\\ 1\end{array}\right)=\left(\begin{array}{c}3\\ 4\\ 3\end{array}\right)$$
𝐁=𝐍𝐑=(3211212112132)(12000012)=(
04134142121231434
).𝐁𝐍𝐑matrix3211212112132matrix12000012
04134142121231434
{\mathbf{B}}={\mathbf{N}}{\mathbf{R}}=\pmatrix{\frac{3}{2}&1&{\frac{1}{2}}\cr 1&2&1\cr{\frac{1}{2}}&1&\frac{3}{2}}\pmatrix{{\frac{1}{2}}&0\cr 0&0\cr 0&{\frac{1}{2}}}=\hbox{}\;\vbox{\kern 46.5111pt\hbox{$\kern 30.00002pt\kern-8.75pt\left(\kern-30.00002pt\vbox{\kern-46.5111pt\vbox{\halign{$#$\hfil\kern 2\p@\kern\@tempdima&\thinspace\hfil$#$\hfil&&\quad\hfil$#$\hfil\cr\hfil\crcr\kern-12.0pt\cr$\hfil\kern 2.0pt\kern 8.75pt&0&4\crcr\kern 2.0pt\cr 1$\hfil\kern 2.0pt\kern 8.75pt&\frac{3}{4}&\frac{1}{4}\cr 2$\hfil\kern 2.0pt\kern 8.75pt&{\frac{1}{2}}&{\frac{1}{2}}\cr 3$\hfil\kern 2.0pt\kern 8.75pt&\frac{1}{4}&\frac{3}{4}\crcr\cr}}\kern-12.0pt}\,\right)$}}.
Thus, starting in state 3, the probability is $`3/4`$ of reaching 4 before 0; this is in agreement with our previous results. From $`𝐭`$ we see that the expected duration of the game, when we start in state 2, is 4.
For an absorbing chain $`𝐏`$, the $`n`$th power $`𝐏^n`$ of the transition probabilities will approach a matrix $`𝐏^{\mathrm{}}`$ of the form
$$𝐏^{\mathrm{}}=\left(\begin{array}{cc}𝐈& \mathrm{𝟎}\\ 𝐁& 𝐐\end{array}\right).$$
We now give our Markov chain version of the Dirichlet problem. We interpret the absorbing states as boundary states and the non-absorbing states as interior states. Let $`B`$ be the set of boundary states and $`D`$ the set of interior states. Let $`f`$ be a function with domain the state space of a Markov chain $`𝐏`$ such that for $`i`$ in $`D`$
$$f(i)=\underset{j}{}P_{ij}f(j).$$
Then $`f`$ is a *harmonic function* for $`𝐏`$. Now $`f`$ again has an averaging property and extends our previous definition. If we represent $`f`$ as a column vector $`𝐟`$, $`f`$ is harmonic if and only if
$$\mathrm{𝐏𝐟}=𝐟.$$
This implies that
$$𝐏^2𝐟=𝐏\mathrm{𝐏𝐟}=\mathrm{𝐏𝐟}=𝐟$$
and in general
$$𝐏^n𝐟=𝐟.$$
Let us write the vector $`𝐟`$ as
$$𝐟=\left(\begin{array}{c}𝐟_B\\ 𝐟_D\end{array}\right)$$
where $`𝐟_B`$ represents the values of $`f`$ on the boundary and $`𝐟_D`$ values on the interior. Then we have
$$\left(\begin{array}{c}𝐟_B\\ 𝐟_D\end{array}\right)=\left(\begin{array}{cc}𝐈& \mathrm{𝟎}\\ 𝐁& 𝐐\end{array}\right)\left(\begin{array}{c}𝐟_B\\ 𝐟_D\end{array}\right)$$
and
$$𝐟_D=\mathrm{𝐁𝐟}_B.$$
We again see that the values of a harmonic function are determined by the values of the function at the boundary points.
Since the entries $`B_{ij}`$ of $`𝐁`$ represent the probability, starting in $`i`$, that the process ends at $`j`$, our last equation states that if you play a game in which your fortune is $`f_j`$ when you are in state $`j`$, then your expected final fortune is equal to your initial fortune; that is, fairness is preserved. As remarked above, from Markov chain theory $`𝐁=\mathrm{𝐍𝐑}`$ where $`𝐍=(𝐈𝐐)^1`$. Thus
$$𝐟_D=(𝐈𝐐)^1\mathrm{𝐑𝐟}_B.$$
(To make the correspondence between this solution and the solution of Section 1.2.5, put $`𝐀=𝐈𝐐`$ and $`𝐮=\mathrm{𝐑𝐟}_B`$.)
A general discussion of absorbing Markov chains may be found in Kemeny, Snell, and Thompson .
###### Exercise 1.2.9
Consider the game played on the grid in Figure 11.
You start at an interior point and move randomly until a boundary point is reached and obtain the payment indicated at this point. Using Markov chain methods find, for each starting state, the expected value of the game. Find also the expected duration of the game.
### 1.3 Random walks on more general networks
#### 1.3.1 General resistor networks and reversible Markov chains
Our networks so far have been very special networks with unit resistors. We will now introduce general resistor networks, and consider what it means to carry out a random walk on such a network.
A *graph* is a finite collection of *points* (also called *vertices* or *nodes*) with certain pairs of points connected by *edges* (also called *branches*). The graph is *connected* if it is possible to go between any two points by moving along the edges. (See Figure 12.)
We assume that $`G`$ is a connected graph and assign to each edge $`xy`$ a resistance $`R_{xy}`$; an example is shown in Figure 13.
The *conductance* of an edge $`xy`$ is $`C_{xy}=1/R_{xy}`$; conductances for our example are shown in Figure 14.
We define a *random walk* on $`G`$ to be a Markov chain with transition matrix $`𝐏`$ given by
$$P_{xy}=\frac{C_{xy}}{C_x}$$
with $`C_x=_yC_{xy}`$. For our example, $`C_a=2`$, $`C_b=3`$, $`C_c=4`$, and $`C_d=5`$, and the transition matrix $`𝐏`$ for the associated random walk is
(
abcda001212b001323c1414012d1525250
)
abcda001212b001323c1414012d1525250
\kern 60.77667pt\kern-8.75pt\left(\kern-60.77667pt\vbox{\kern-56.38887pt\vbox{\halign{$#$\hfil\kern 2\p@\kern\@tempdima&\thinspace\hfil$#$\hfil&&\quad\hfil$#$\hfil\cr\hfil\crcr\kern-12.0pt\cr$\hfil\kern 2.0pt\kern 8.75pt&a&b&c&d\crcr\kern 2.0pt\cr a$\hfil\kern 2.0pt\kern 8.75pt&0&0&{\frac{1}{2}}&{\frac{1}{2}}\cr b$\hfil\kern 2.0pt\kern 8.75pt&0&0&\frac{1}{3}&\frac{2}{3}\cr c$\hfil\kern 2.0pt\kern 8.75pt&\frac{1}{4}&\frac{1}{4}&0&{\frac{1}{2}}\cr d$\hfil\kern 2.0pt\kern 8.75pt&\frac{1}{5}&\frac{2}{5}&\frac{2}{5}&0\crcr\cr}}\kern-12.0pt}\,\right)
Its graphical representation is shown in Figure 15.
Since the graph is connected, it is possible for the walker to go between any two states. A Markov chain with this property is called an *ergodic* Markov chain. Regular chains, which were introduced in Section 1.2.6, are always ergodic, but ergodic chains are not always regular (see Exercise 1.3.1).
For an ergodic chain, there is a unique probability vector $`𝐰`$ that is a fixed vector for $`𝐏`$, i.e., $`\mathrm{𝐰𝐏}=𝐰`$. The component $`w_j`$ of $`𝐰`$ represents the proportion of times, in the long run, that the walker will be in state $`j`$. For random walks determined by electric networks, the fixed vector is given by $`w_j=C_j/C`$, where $`C=_xC_x`$. (You are asked to prove this in Exercise 1.3.2.) For our example $`C_a=2`$, $`C_b=3`$, $`C_c=4`$, $`C_d=5`$, and $`C=14`$. Thus $`𝐰=(2/14,3/14,4/14,5/14)`$. We can check that $`𝐰`$ is a fixed vector by noting that
$$\left(\begin{array}{cccc}\frac{2}{14}& \frac{3}{14}& \frac{4}{14}& \frac{5}{14}\end{array}\right)\left(\begin{array}{cccc}0& 0& \frac{1}{2}& \frac{1}{2}\\ 0& 0& \frac{1}{3}& \frac{2}{3}\\ \frac{1}{4}& \frac{1}{4}& 0& \frac{1}{2}\\ \frac{1}{5}& \frac{2}{5}& \frac{2}{5}& 0\end{array}\right)=\left(\begin{array}{cccc}\frac{2}{14}& \frac{3}{14}& \frac{4}{14}& \frac{5}{14}\end{array}\right).$$
In addition to being ergodic, Markov chains associated with networks have another property called *reversibility*. An ergodic chain is said to be *reversible* if $`w_xP_{xy}=w_yP_{yx}`$ for all $`x,y`$. That this is true for our network chains follows from the fact that
$$C_xP_{xy}=C_x\frac{C_{xy}}{C_x}=C_{xy}=C_{yx}=C_y\frac{C_{yx}}{C_y}=C_yP_{yx}.$$
Thus, dividing the first and last term by $`C`$, we have $`w_xP_{xy}=w_yP_{yx}`$.
To see the meaning of reversibility, we start our Markov chain with initial probabilities $`𝐰`$ (in equilibrium) and observe a few states, for example
$$acbd.$$
The probability that this sequence occurs is
$$w_aP_{ac}P_{cb}P_{bd}=\frac{2}{14}\frac{1}{2}\frac{1}{4}\frac{2}{3}=\frac{1}{84}.$$
The probability that the reversed sequence
$$dbca$$
occurs is
$$w_dP_{db}P_{bc}P_{ca}=\frac{5}{14}\frac{2}{5}\frac{1}{3}\frac{1}{4}=\frac{1}{84}.$$
Thus the two sequences have the same probability of occurring.
In general, when a reversible Markov chain is started in equilibrium, probabilities for sequences in the correct order of time are the same as those with time reversed. Thus, from data, we would never be able to tell the direction of time.
If $`𝐏`$ is any reversible ergodic chain, then $`𝐏`$ is the transition matrix for a random walk on an electric network; we have only to define $`C_{xy}=w_xP_{xy}`$. Note, however, if $`P_{xx}0`$ the resulting network will need a conductance from $`x`$ to $`x`$ (see Exercise 1.3.4). Thus reversibility characterizes those ergodic chains that arise from electrical networks. This has to do with the fact that the physical laws that govern the behavior of steady electric currents are invariant under time-reversal (see Onsager ).
When all the conductances of a network are equal, the associated random walk on the graph $`G`$ of the network has the property that, from each point, there is an equal probability of moving to each of the points connected to this point by an edge. We shall refer to this random walk as *simple random walk* on $`G`$. Most of the examples we have considered so far are simple random walks. Our first example of a random walk on Madison Avenue corresponds to simple random walk on the graph with points $`0,1,2,\mathrm{},N`$ and edges the streets connecting these points. Our walks on two dimensional graphs were also simple random walks.
###### Exercise 1.3.1
Give an example of an ergodic Markov chain that is not regular. (Hint: a chain with two states will do.)
###### Exercise 1.3.2
Show that, if $`𝐏`$ is the transition matrix for a random walk determined by an electric network, then the fixed vector $`𝐰`$ is given by $`w_x=\frac{C_x}{C}`$ where $`C_x=_yC_{xy}`$ and $`C=_xC_x`$.
###### Exercise 1.3.3
Show that, if $`𝐏`$ is a reversible Markov chain and $`a,b,c`$ are any three states, then the probability, starting at $`a`$, of the cycle $`abca`$ is the same as the probability of the reversed cycle $`acba`$. That is $`P_{ab}P_{bc}P_{ca}=P_{ac}P_{cb}P_{ba}`$. Show, more generally, that the probability of going around any cycle in the two different directions is the same. (Conversely, if this cyclic condition is satisfied, the process is reversible. For a proof, see Kelly .)
###### Exercise 1.3.4
Assume that $`𝐏`$ is a reversible Markov chain with $`P_{xx}=0`$ for all $`x`$. Define an electric network by $`C_{xy}=w_xP_{xy}`$. Show that the Markov chain associated with this circuit is $`𝐏`$. Show that we can allow $`P_{xx}>0`$ by allowing a conductance from $`x`$ to $`x`$.
###### Exercise 1.3.5
For the *Ehrenfest urn model*, there are two urns that together contain $`N`$ balls. Each second, one of the $`N`$ balls is chosen at random and moved to the other urn. We form a Markov chain with states the number of balls in one of the urns. For $`N=4`$, the resulting transition matrix is
𝐏=(
01234001000114034002012012030034014400010
).𝐏
01234001000114034002012012030034014400010
{\mathbf{P}}=\hbox{}\;\vbox{\kern 59.39998pt\hbox{$\kern 75.0001pt\kern-8.75pt\left(\kern-75.0001pt\vbox{\kern-59.39998pt\vbox{\halign{$#$\hfil\kern 2\p@\kern\@tempdima&\thinspace\hfil$#$\hfil&&\quad\hfil$#$\hfil\cr\hfil\crcr\kern-12.0pt\cr$\hfil\kern 2.0pt\kern 8.75pt&0&1&2&3&4\crcr\kern 2.0pt\cr 0$\hfil\kern 2.0pt\kern 8.75pt&0&1&0&0&0\cr 1$\hfil\kern 2.0pt\kern 8.75pt&\frac{1}{4}&0&\frac{3}{4}&0&0\cr 2$\hfil\kern 2.0pt\kern 8.75pt&0&{\frac{1}{2}}&0&{\frac{1}{2}}&0\cr 3$\hfil\kern 2.0pt\kern 8.75pt&0&0&\frac{3}{4}&0&\frac{1}{4}\cr 4$\hfil\kern 2.0pt\kern 8.75pt&0&0&0&1&0\crcr\cr}}\kern-12.0pt}\,\right)$}}.
Show that the fixed vector $`𝐰`$ is the binomial distribution $`𝐰=(\frac{1}{16},\frac{4}{16},\frac{6}{16},\frac{4}{16},\frac{1}{16})`$. Determine the electric network associated with this chain.
#### 1.3.2 Voltages for general networks; probabilistic interpretation
We assume that we have a network of resistors assigned to the edges of a connected graph. We choose two points $`a`$ and $`b`$ and put a one-volt battery across these points establishing a voltage $`v_a=1`$ and $`v_b=0`$, as illustrated in Figure 16.
We are interested in finding the voltages $`v_x`$ and the currents $`i_{xy}`$ in the circuit and in giving a probabilistic interpretation to these quantities.
We begin with the probabilistic interpretation of voltage. It will come as no surprise that we will interpret the voltage as a hitting probability, observing that both functions are harmonic and that they have the same boundary values.
By Ohm’s Law, the currents through the resistors are determined by the voltages by
$$i_{xy}=\frac{v_xv_y}{R_{xy}}=(V_xv_y)C_{xy}.$$
Note that $`i_{xy}=i_{yx}`$. Kirchhoff’s Current Law requires that the total current flowing into any point other than $`a`$ or $`b`$ is 0. That is, for $`xa,b`$
$$\underset{y}{}i_{xy}=0.$$
This will be true if
$$\underset{y}{}(v_xv_y)C_{xy}=0$$
or
$$v_x\underset{y}{}C_{xy}=\underset{y}{}C_{xy}v_y.$$
Thus Kirchhoff’s Current Law requires that our voltages have the property that
$$v_x=\underset{y}{}\frac{C_{xy}}{C_x}v_y=\underset{y}{}P_{xy}v_y$$
for $`xa,b`$. This means that the voltage $`v_x`$ is harmonic at all points $`xa,b`$.
Let $`h_x`$ be the probability, starting at $`x`$, that state $`a`$ is reached before $`b`$. Then $`h_x`$ is also harmonic at all points $`xa,b`$. Furthermore
$$v_a=h_a=1$$
and
$$v_b=h_b=0.$$
Thus if we modify $`𝐏`$ by making $`a`$ and $`b`$ absorbing states, we obtain an absorbing Markov chain $`\overline{𝐏}`$ and $`v`$ and $`h`$ are both solutions to the Dirichlet problem for the Markov chain with the same boundary values. Hence $`v=h`$.
For our example, the transition probabilities $`\overline{P}_{xy}`$ are shown in Figure 17.
The function $`v_x`$ is harmonic for $`\overline{P}`$ with boundary values $`v_a=1,v_b=0`$.
To sum up, we have the following:
Intrepretation of Voltage. When a unit voltage is applied between $`a`$ and $`b`$, making $`v_a=1`$ and $`v_b=0`$, the voltage $`v_x`$ at any point $`x`$ represents the probability that a walker starting from $`x`$ will return to $`a`$ before reaching $`b`$.
In this probabilistic interpretation of voltage, we have assumed a unit voltage, but we could have assumed an arbitrary voltage $`v_a`$ between $`a`$ and $`b`$. Then the hitting probability $`h_x`$ would be replaced by an expected value in a game where the player starts at $`x`$ and is paid $`v_a`$ if $`a`$ is reached before $`b`$ and 0 otherwise.
Let’s use this interpretation of voltage to find the voltages for our example. Referring back to Figure 17, we see that
$$v_a=1$$
$$v_b=0$$
$$v_c=\frac{1}{4}+\frac{1}{2}v_d$$
$$v_d=\frac{1}{5}+\frac{2}{5}v_c.$$
Solving these equations yields $`v_c=\frac{7}{16}`$ and $`v_d=\frac{3}{8}`$. From these voltages we obtain the current $`i_{xy}`$. For example $`i_{cd}=(\frac{7}{16}\frac{3}{8})2=\frac{1}{8}`$. The resulting voltages and currents are shown in Figure 18.
The voltage at $`c`$ is $`\frac{7}{16}`$ and so this is also the probability, starting at $`c`$, of reaching $`a`$ before $`b`$.
#### 1.3.3 Probabilistic interpretation of current
We turn now to the probabilistic interpretation of current. This interpretation is found by taking a naive view of the process of electrical conduction: We imagine that positively charged particles enter the network at point $`a`$ and wander around from point to point until they finally arrive at point $`b`$, where they leave the network. (It would be more realistic to imagine negatively charged particles entering at $`b`$ and leaving at $`a`$, but realism is not what we’re after.) To determine the current $`i_{xy}`$ along the branch from $`x`$ to $`y`$, we consider that in the course of its peregrinations the point may pass once or several times along the branch from $`x`$ to $`y`$, and in the opposite direction from $`y`$ to $`x`$. We may now hypothesize that the current $`i_{xy}`$ is proportional to the expected net number of movements along the edge from $`x`$ to $`y`$, where movements from $`y`$ back to $`x`$ are counted as negative. This hypothesis is correct, as we will now show.
The walker begins at $`a`$ and walks until he reaches $`b`$; note that if he returns to $`a`$ before reaching $`b`$, he keeps on going. Let $`u_x`$ be the expected number of visits to state $`x`$ before reaching $`b`$. Then $`u_b=0`$ and, for $`xa,b`$,
$$u_x=\underset{y}{}u_yP_{yx}.$$
This last equation is true because, for $`xa,b`$, every entrance to $`x`$ must come from some $`y`$.
We have seen that $`C_xP_{xy}=C_yP_{yx}`$; thus
$$u_x=\underset{y}{}u_y\frac{P_{xy}C_x}{C_y}$$
or
$$\frac{u_x}{C_x}=\underset{y}{}P_{xy}\frac{u_y}{C_y}.$$
This means that $`v_x=u_x/C_x`$ is harmonic for $`xa,b`$. We have also $`v_b=0`$ and $`v_a=u_a/C_a`$. This implies that $`v_x`$ is the voltage at $`x`$ when we put a battery from $`a`$ to $`b`$ that establishes a voltage $`u_a/C_a`$ at $`a`$ and voltage 0 at $`b`$. (We remark that the expression $`v_x=u_x/C_x`$ may be understood physically by viewing $`u_x`$ as charge and $`C_x`$ as capacitance; see Kelly for more about this.)
We are interested in the current that flows from $`x`$ to $`y`$. This is
$$i_{xy}=(v_xv_y)C_{xy}=\left(\frac{u_x}{C_x}\frac{u_y}{C_y}\right)C_{xy}=\frac{u_xC_{xy}}{C_x}\frac{u_yC_{yx}}{C_y}=u_xP_{xy}u_yP_{yx}.$$
Now $`u_xP_{xy}`$ is the expected number of times our walker will go from $`x`$ to $`y`$ and $`u_yP_{yx}`$ is the expected number of times he will go from $`y`$ to $`x`$. Thus the current $`i_{xy}`$ is the expected value for the net number of times the walker passes along the edge from $`x`$ to $`y`$. Note that for any particular walk this net value will be an integer, but the expected value will not.
As we have already noted, the currents $`i_{xy}`$ here are not those of our original electrical problem, where we apply a 1-volt battery, but they are proportional to those original currents. To determine the constant of proportionality, we note the following characteristic property of the new currents $`i_{xy}`$: The total current flowing into the network at $`a`$ (and out at $`b`$) is 1. In symbols,
$$\underset{y}{}i_{ay}=1.$$
Indeed, from our probabilistic interpretation of $`i_{xy}`$ this sum represents the expected value of the difference between the number of times our walker leaves $`a`$ and enters $`a`$. This number is necessarily one and so the current flowing into $`a`$ is 1.
This unit current flow from $`a`$ to $`b`$ can be obtained from the currents in the original circuit, corresponding to a 1-volt battery, by dividing through by the total amount of current $`_yi_{ay}`$ flowing into $`a`$; doing this to the currents in our example yields the unit current flow shown in Figure 19.
This shows that the constant of proportionality we were seeking to determine is the reciprocal of the amount of current that flows through the circuit when a I-volt battery is applied between $`a`$ and $`b`$. This quantity, called the effective resistance between $`a`$ and $`b`$, is discussed in detail in Section 1.3.4.
To sum up, we have the following:
Interpretation of Current. When a unit current flows into $`a`$ and out of $`b`$, the current $`i_{xy}`$ flowing through the branch connecting $`x`$ to $`y`$ is equal to the expected net number of times that a walker, starting at $`a`$ and walking until he reaches $`b`$, will move along the branch from $`x`$ to $`y`$. These currents are proportional to the currents that arise when a unit voltage is applied between $`a`$ and $`b`$, the constant of proportionality being the effective resistance of the network.
We have seen that we can estimate the voltages by simulation. We can now do the same for the currents. We have to estimate the expected value for the net number of crossings of $`xy`$. To do this, we start a large number of walks at $`a`$ and, for each one, record the net number of crossings of each edge and then average these as an estimate for the expected value. Carrying out 10,000 such walks yielded the results shown in Figure 20.
The results of simulation are in good agreement with the theoretical values of current. As was the case for estimates of the voltages by simulation, we have statistical errors. Our estimates have the property that the total current flowing into $`a`$ is 1, out of $`b`$ is 1, and into any other point it is 0. This is no accident; the explanation is that the history of each walk would have these properties, and these properties are not destroyed by averaging.
###### Exercise 1.3.6
Kingman introduced a different model for current flow. Kelly gave a new interpretation of this model. Both authors use continuous time. A discrete time version of Kelly’s interpretation would be the following: At each point of the graph there is a black or a white button. Each second an edge is chosen; edge $`xy`$ is chosen with probability $`C_{xy}/C`$ where $`C`$ is the sum of the conductances. The buttons on the edge chosen are then interchanged. When a button reaches $`a`$ it is painted black, and when a button reaches $`b`$ it is painted white. Show that there is a limiting probability $`p_x`$ that site $`x`$ has a black button and that $`p_x`$ is the voltage $`v_x`$ at $`x`$ when a unit voltage is imposed between $`a`$ and $`b`$. Show that the current $`i_{xy}`$ is proportional to the net flow of black buttons along the edge $`xy`$. Does this suggest a hypothesis about the behavior of conduction electrons in metals?
#### 1.3.4 Effective resistance and the escape probability
When we impose a voltage $`v`$ between points $`a`$ and $`b`$, a voltage $`v_a=v`$ is established at $`a`$ and $`v_b=0`$, and a current $`i_a=_xi_{ax}`$ will flow into the circuit from the outside source. The amount of current that flows depends upon the overall resistance in the circuit. We define the *effective resistance* $`R_{\text{eff}}`$ between $`a`$ and $`b`$ by $`R_{\text{eff}}=v_a/i_a`$. The reciprocal quantity $`C_{\text{eff}}=1/R_{\text{eff}}=i_a/v_a`$ is the *effective conductance*. If the voltage between $`a`$ and $`b`$ is multiplied by a constant, then the currents are multiplied by the same constant, so $`R_{\text{eff}}`$ depends only on the ratio of $`v_a`$ to $`i_a`$.
Let us calculate $`R_{\text{eff}}`$ for our example. When a unit voltage was imposed, we obtained the currents shown in Figure 18. The total current flowing into the circuit is $`i_a=9/16+10/16=19/16`$. Thus the effective resistance is
$$R_{\text{eff}}=\frac{v_a}{i_a}=\frac{1}{\frac{19}{16}}=\frac{16}{19}.$$
We can interpret the effective conductance probabilistically as an escape probability. When $`v_a=1`$, the effective conductance equals the total current $`i_a`$ flowing into $`a`$. This current is
$$i_a=\underset{y}{}(v_av_y)C_{ay}=\underset{y}{}(v_av_y)\frac{C_{ay}}{C_a}C_a=C_a(1\underset{y}{}P_{ay}v_y)=C_ap_{\text{esc}}$$
where $`p_{\text{esc}}`$ is the probability, starting at $`a`$, that the walk reaches $`b`$ before returning to $`a`$. Thus
$$C_{\text{eff}}=C_ap_{\text{esc}}$$
and
$$p_{\text{esc}}=\frac{C_{\text{eff}}}{C_a}.$$
In our example $`C_a=2`$ and we found that $`i_a=19/16`$. Thus
$$p_{\text{esc}}=\frac{19}{32}.$$
In calculating effective resistances, we shall use two important facts about electric networks. First, if two resistors are connected in series, they may be replaced by a single resistor whose resistance is the sum of the two resistances. (See Figure 21.)
Secondly, two resistors in parallel may be replaced by a single resistor with resistance $`R`$ such that
$$\frac{1}{R}=\frac{1}{R_1}+\frac{1}{R_2}=\frac{R_1R_2}{R_1+R_2}.$$
(See Figure 22.)
The second rule can be stated more simply in terms of conductances: If two resistors are connected in parallel, they may be replaced by a single resistor whose conductance is the sum of the two conductances.
We illustrate the use of these ideas to compute the effective resistance between two adjacent points of a unit cube of unit resistors, as shown in Figure 23.
We put a unit battery between $`a`$ and $`b`$. Then, by symmetry, the voltages at $`c`$ and $`d`$ will be the same as will those at $`e`$ and $`f`$. Thus our circuit is equivalent to the circuit shown in Figure 24.
Using the laws for the effective resistance of resistors in series and parallel, this network can be successively reduced to a single resistor of resistance $`7/12`$ ohms, as shown in Figure 25.
Thus the effective resistance is $`7/12`$. The current flowing into $`a`$ from the battery will be $`i_a=\frac{1}{R_{\text{eff}}}=12/7`$. The probability that a walk starting at $`a`$ will reach $`b`$ before returning to $`a`$ is
$$p_{\text{esc}}=\frac{i_a}{C_a}=\frac{\frac{12}{7}}{3}=\frac{4}{7}.$$
This example and many other interesting connections between electric networks and graph theory may be found in Bollobas .
###### Exercise 1.3.7
A bug walks randomly on the unit cube (see Figure 26).
If the bug starts at $`a`$, what is the probability that it reaches food at $`b`$ before returning to $`a`$?
###### Exercise 1.3.8
Consider the Ehrenfest urn model with $`N=4`$ (see Exercise 1.3.5). Find the probability, starting at 0, that state 4 is reached before returning to 0.
###### Exercise 1.3.9
Consider the ladder network shown in Figure 27.
Show that if $`R_n`$ is the effective resistance of a ladder with $`n`$ rungs then $`R_1=2`$ and
$$R_{n+1}=\frac{2+2R_n}{2+R_n}.$$
Use this to show that $`lim_n\mathrm{}R_n=\sqrt{2}`$.
###### Exercise 1.3.10
A drunken tourist starts at her hotel and walks at random through the streets of the idealized Paris shown in Figure 28.
Find the probability that she reaches the Arc de Triomphe before she reaches the outskirts of town.
#### 1.3.5 Currents minimize energy dissipation
We have seen that when we impose a voltage between $`a`$ and $`b`$ voltages $`v_x`$ are established at the points and currents $`i_{xy}`$ flow through the resistors. In this section we shall give a characterization of the currents in terms of a quantity called *energy dissipation*. When a current $`i_{xy}`$ flows through a resistor, the energy dissipated is
$$i_{xy}^2R_{xy};$$
this is the product of the current $`i_{xy}`$ and the voltage $`v_{xy}=i_{xy}R_{xy}`$. The *total energy dissipation* in the circuit is
$$E=\frac{1}{2}\underset{x,y}{}i_{xy}^2R_{xy}.$$
Since $`i_{xy}R_{xy}=v_xv_y`$, we can also write the energy dissipation as
$$E=\frac{1}{2}\underset{x,y}{}i_{xy}(v_xv_y).$$
The factor 1/2 is necessary in this formulation since each edge is counted twice in this sum. For our example, we see from Figure 18 that
$$E=\left(\frac{9}{16}\right)^21+\left(\frac{10}{16}\right)^21+\left(\frac{7}{16}\right)^21+\left(\frac{2}{16}\right)^2\frac{1}{2}+\left(\frac{12}{16}\right)^2\frac{1}{2}=\frac{19}{16}.$$
If a source (battery) establishes voltages $`v_a`$ and $`v_b`$ at $`a`$ and $`b`$, then the energy supplied is $`(v_av_b)i_a`$ where $`i_a=_xi_{ax}`$. By conservation of energy, we would expect this to be equal to the energy dissipated. In our example $`v_av_b=1`$ and $`i_a=\frac{19}{16}`$, so this is the case. We shall show that this is true in a somewhat more general setting.
Define a *flow* $`𝐣`$ from $`a`$ to $`b`$ to be an assignment of numbers $`j_{xy}`$ to pairs $`xy`$ such that
(a) $`j_{xy}=j_{yx}`$
(b) $`_yj_{xy}=0`$ if $`xa,b`$
(c) $`j_{xy}=0`$ if $`x`$ and $`y`$ are not adjacent.
We denote by $`j_x=_yj_{xy}`$ the flow into $`x`$ from the outside. By (b) $`j_x=0`$ for $`xa,b`$. Of course $`j_b=j_a`$. To verify this, note that
$$j_a+j_b=\underset{x}{}j_x=\underset{x}{}\underset{y}{}j_{xy}=\frac{1}{2}\underset{x,y}{}(j_{xy}+j_{yx})=0,$$
since $`j_{xy}=j_{yx}`$.
With this terminology, we can now formulate the following version of the principle of conservation of energy:
Conservation of Energy. Let $`w`$ be any function defined on the points of the graph and $`𝐣`$ a flow from $`a`$ to $`b`$. Then
$$(w_aw_b)j_a=\frac{1}{2}\underset{x,y}{}(w_xw_y)j_{xy}.$$
Proof.
$`{\displaystyle \underset{x,y}{}}(w_xw_y)j_{xy}`$ $`=`$ $`{\displaystyle \underset{x}{}}(w_x{\displaystyle \underset{y}{}}j_{xy}){\displaystyle \underset{y}{}}(w_y{\displaystyle \underset{x}{}}j_{xy})`$
$`=`$ $`w_a{\displaystyle \underset{y}{}}j_{ay}+w_b{\displaystyle \underset{y}{}}j_{by}w_a{\displaystyle \underset{x}{}}j_{xa}w_b{\displaystyle \underset{x}{}}j_{xb}`$
$`=`$ $`w_aj_a+w_bj_bw_a(j_a)w_b(j_b)`$
$`=`$ $`2(w_aw_b)j_a.`$
Thus
$$(w_aw_b)j_a=\frac{1}{2}\underset{x,y}{}(w_xw_y)j_{xy}$$
as was to be proven. $`\mathrm{}`$
If we now impose a voltage $`v_a`$ between $`a`$ and $`b`$ with $`v_b=0`$, we obtain voltages $`v_x`$ and currents $`i_{xy}`$. The currents $`𝐢`$ give a flow from $`a`$ to $`b`$ and so by the previous result, we conclude that
$$v_ai_a=\frac{1}{2}\underset{x,y}{}(v_xv_y)i_{xy}=\frac{1}{2}\underset{x,y}{}i_{xy}^2R_{xy}.$$
Recall that $`R_{\text{eff}}=v_a/i_a`$. Thus in terms of resistances we can write this as
$$i_{xy}^2R_{\text{eff}}=\frac{1}{2}\underset{x,y}{}i_{xy}^2R_{xy}.$$
If we adjust $`v_a`$ so that $`i_a=1`$, we call the resulting flow *the unit current flow* from $`a`$ to $`b`$. The unit current flow from $`a`$ to $`b`$ is a particular example of a *unit flow* from $`a`$ to $`b`$, which we define to be any flow $`i_{xy}`$ from $`a`$ to $`b`$ for which $`i_a=i_b=1`$. The formula above shows that the energy dissipated by the unit current flow is just $`R_{\text{eff}}`$. According to a basic result called Thomson’s Principle, this value is smaller than the energy dissipated by any other unit flow from $`a`$ to $`b`$. Before proving this principle, let us watch it in action in the example worked out above.
Recall that, for this example, we found the true values and some approximate values for the unit current flow; these were shown in Figure 20. The energy dissipation for the true currents is
$$E=R_{\text{eff}}=\frac{16}{19}=.8421053.$$
Our approximate currents also form a unit flow and, for these, the energy dissipation is
$$\overline{E}=(.4754)^21+(.5246)^21+(.3672)^21+(.1082)^2\frac{1}{2}+(.6328)^2\frac{1}{2}=.8421177.$$
We note that $`\overline{E}`$ is greater than $`E`$, though just barely.
Thomson’s Principle. (Thomson ). If $`𝐢`$ is the unit flow from $`a`$ to $`b`$ determined by Kirchhoff’s Laws, then the energy dissipation $`\frac{1}{2}_{x,y}i_{xy}^2R_{xy}`$ minimizes the energy dissipation $`\frac{1}{2}_{x,y}j_{xy}^2R_{xy}`$ among all unit flows $`𝐣`$ from $`a`$ to $`b`$.
Proof. Let $`𝐣`$ be any unit flow from $`a`$ to $`b`$ and let $`d_{xy}=j_{xy}i_{xy}`$. Then $`𝐝`$ is a flow from $`a`$ to $`b`$ with $`d_a=_xd_{ax}=11=0`$.
$`{\displaystyle \underset{x,y}{}}j_{xy}^2R_{xy}`$ $`=`$ $`{\displaystyle \underset{x,y}{}}(i_{xy}+d_{xy})^2R_{xy}`$
$`=`$ $`{\displaystyle \underset{x,y}{}}i_{xy}^2R_{xy}+2{\displaystyle \underset{x,y}{}}i_{xy}R_{xy}d_{xy}+{\displaystyle \underset{x,y}{}}d_{xy}^2R_{xy}`$
$`=`$ $`{\displaystyle \underset{x,y}{}}i_{xy}^2R_{xy}+2{\displaystyle \underset{x,y}{}}(v_xv_y)d_{xy}+{\displaystyle \underset{x,y}{}}d_{xy}^2R_{xy}.`$
Setting $`𝐰=𝐯`$ and $`𝐣=𝐝`$ in the conservation of energy result above shows that the middle term is $`4(v_av_b)d_a=0`$. Thus
$$\underset{x,y}{}j_{xy}^2R_{xy}=\underset{x,y}{}i_{xy}^2R_{xy}+\underset{x,y}{}d_{xy}^2R_{xy}\underset{x,y}{}i_{xy}^2R_{xy}.$$
This completes the proof. $`\mathrm{}`$
###### Exercise 1.3.11
The following is the so-called “dual form” of Thomson’s Principle. Let $`u`$ be any function on the points of the graph $`G`$ of a circuit such that $`u_a=1`$ and $`u_b=0`$. Then the energy dissipation
$$\frac{1}{2}\underset{x,y}{}(u_xu_y)^2C_{xy}$$
is minimized by the voltages $`v_x`$ that result when a unit voltage is established between $`a`$ and $`b`$, i.e., $`v_a=1`$, $`v_b=0`$, and the other voltages are determined by Kirchhoff’s Laws. Prove this dual principle. This second principle is known nowadays as the *Dirichlet Principle*, though it too was discovered by Thomson.
###### Exercise 1.3.12
In Section 1.2.4 we stated that, to solve the Dirichlet problem by the method of relaxations, we could start with an arbitrary initial guess. Show that when we replace the value at a point by the average of the neighboring points the energy dissipation, as expressed in Exercise 1.3.11, can only decrease. Use this to prove that the relaxation method converges to a solution of the Dirichlet problem for an arbitrary initial guess.
### 1.4 Rayleigh’s Monotonicity Law
#### 1.4.1 Rayleigh’s Monotonicity Law
Next we will study Rayleigh’s Monotonicity Law. This law from electric network theory will be an important tool in our future study of random walks. In this section we will give an example of the use of this law.
Consider a random walk on streets of a city as in Figure 29.
Let $`p_{\text{esc}}`$ be the probability that a walker starting from $`a`$ reaches $`b`$ before returning to $`a`$. Assign to each edge a unit resistance and maintain a voltage of one volt between $`a`$ and $`b`$; then a current $`i_a`$ will flow into the circuit and we showed in Section 1.3.4 that
$$p_{\text{esc}}=\frac{i_a}{C_a}=\frac{i_a}{4}.$$
Now suppose that one of the streets (not connected to $`a`$) becomes blocked. Our walker must choose from the remaining streets if he reaches a corner of this street. The modified graph will be as in Figure 30.
We want to show that the probability of escaping to $`b`$ from $`a`$ is decreased.
Consider this problem in terms of our network. Blocking a street corresponds to replacing a unit resistor by an infinite resistor. This should have the effect of increasing the effective resistance $`R_{\text{eff}}`$ of the circuit between $`a`$ and $`b`$. If so, when we put a unit voltage between $`a`$ and $`b`$ less current will flow into the circuit and
$$p_{\text{esc}}=\frac{i_a}{4}=\frac{1}{4R_{\text{eff}}}$$
will decrease.
Thus we need only show that when we increase the resistance in one part of a circuit, the effective resistance increases. This fact, known as Rayleigh’s Monotonicity Law, is almost self-evident. Indeed, the father of electromagnetic theory, James Clerk Maxwell, regarded this to be the case. In his *Treatise on Electricity and Magnetism* (, p. 427), he wrote
> If the specific resistance of any portion of the conductor be changed, that of the remainder being unchanged, the resistance of the whole conductor will be increased if that of the portion is increased, and diminished if that of the portion is diminished. This principle may be regarded as self-evident ….
Rayleigh’s Monotonicity Law. If the resistances of a circuit are increased, the effective resistance $`R_{\text{eff}}`$ between any two points can only increase. If they are decreased, it can only decrease.
Proof. Let $`𝐢`$ be the unit current flow from $`a`$ to $`b`$ with the resistors $`R_{xy}`$. Let $`𝐣`$ be the unit current flow from $`a`$ to $`b`$ with the resistors $`\overline{R}_{xy}`$ with $`\overline{R}_{xy}R_{xy}`$. Then
$$\overline{R}_{\text{eff}}=\frac{1}{2}\underset{x,y}{}j_{xy}^2\overline{R}_{xy}\frac{1}{2}\underset{x,y}{}j_{xy}^2R_{xy}.$$
But since $`𝐣`$ is a unit flow from $`a`$ to $`b`$, Thomson’s Principle tells us that the energy dissipation, calculated with resistors $`R_{xy}`$, is bigger than that for the true currents determined by these resistors: that is
$$\frac{1}{2}\underset{x,y}{}j_{xy}^2R_{xy}\frac{1}{2}\underset{x,y}{}i_{xy}^2R_{xy}=R_{\text{eff}}.$$
Thus $`\overline{R}_{\text{eff}}R_{\text{eff}}`$. The proof for the case of decreasing resistances is the same.
###### Exercise 1.4.1
Consider a graph $`G`$ and let $`R_{xy}`$ and $`\overline{R}_{xy}`$ be two different assignments of resistances to the edges of $`G`$. Let $`\widehat{R}_{xy}=\overline{R}_{xy}+R_{xy}`$. Let $`R_{\text{eff}}`$, $`\overline{R}_{\text{eff}}`$, and $`\widehat{R}_{\text{eff}}`$ be the effective resistances when $`R`$, $`\overline{R}`$, and $`\widehat{R}`$, respectively, are used. Prove that
$$\widehat{R}_{\text{eff}}\overline{R}_{\text{eff}}+R_{\text{eff}}.$$
Conclude that the effective resistance of a network is a concave function of any of its component resistances (Shannon and Hagelbarger .)
###### Exercise 1.4.2
Show that the effective resistance of the circuit in Figure 31 is greater than or equal to the effective resistance of the circuit in Figure 32.
Use this to show the following inequality for $`R_{ij}0`$:
$$\frac{1}{\frac{1}{R_{11}+R_{12}}+\frac{1}{R_{21}+R_{22}}}\frac{1}{\frac{1}{R_{11}}+\frac{1}{R_{21}}}+\frac{1}{\frac{1}{R_{12}}+\frac{1}{R_{22}}}.$$
See the note by Lehman for a proof of the general Minkowski inequality by this method.
###### Exercise 1.4.3
Let $`𝐏`$ be the transition matrix associated with an electric network and let $`a,b,r,s`$ be four points on the network. Let $`\overline{𝐏}`$ be a transition matrix defined on the state-space $`S=\{a,b,r,s\}`$. Let $`\overline{P}_{ii}=0`$ and for $`ij`$ let $`\overline{P}_{ij}`$ be the probability that, if the chain $`𝐏`$ is started in state $`i`$, then the next time it is in the set $`S\{i\}`$ it is in the state $`j`$. Show that $`\overline{𝐏}`$ is a reversible Markov chain and hence corresponds to an electric network of the form of a Wheatstone Bridge, shown in Figure 33.
Explain how this proves that, in order to prove the Monotonicity Law, it is sufficient to prove that $`R_{\text{eff}}`$ is a monotone function of the component resistances for a Wheatstone Bridge. Give a direct proof of the Monotonicity Law for this special case.
#### 1.4.2 A probabilistic explanation of the Monotonicity Law
We have quoted Maxwell’s assertion that Rayleigh’s Monotonicity Law may be regarded as self-evident, but one might feel that any argument in terms of electricity is only self-evident if we know what electricity is. In Cambridge, they tell the following story about Maxwell: Maxwell was lecturing and, seeing a student dozing off, awakened him, asking, “Young man, what is electricity?” “I’m terribly sorry, sir,” the student replied, ‘I knew the answer but I have forgotten it.” Maxwell’s response to the class was, “Gentlemen, you have just witnessed the greatest tragedy in the history of science. The one person who knew what electricity is has forgotten it.”
To say that our intuition about the Monotonicity Law is only as solid as our understanding of electricity is not really a valid argument, of course, because in saying that this law is self-evident we are secretly depending on the analogy between electricity and the flow of water (see Feynman , Vol. 2, Chapter 12). We just can’t believe that if a water main gets clogged the total rate of flow out of the local reservoir is going to increase. But as soon as we admit this, some pedant will ask if we’re talking about flows with low Reynolds number, or what, and we’ll have to admit that we don’t understand water any better than we understand electricity.
Whatever our feelings about electricity or the flow of water, it seems desirable to have an explanation of the Monotonicity Law in terms of our random walker. We now give such an explanation.
We begin by collecting some notation and results from previous sections. As usual, we have a network of conductances (streets) and a walker who moves from point $`x`$ to point $`y`$ with probability
$$P_{xy}=\frac{C_{xy}}{C_x}$$
where $`C_{xy}`$ is the conductance from $`x`$ to $`y`$ and $`C_x=_yC_{xy}`$. We choose two preferred points $`a`$ and $`b`$. The walker starts at $`a`$ and walks until he reaches $`b`$ or returns to $`a`$. We denote by $`v_x`$ the probability that the walker, starting at $`a`$, reaches $`a`$ before $`b`$. Then $`v_a=1`$, $`v_b=0`$, and the function $`v_x`$ is harmonic at all points $`xa,b`$. We denote by $`p_{\text{esc}}`$ the probability that the walker, starting at $`a`$, reaches $`b`$ before returning to $`a`$. Then
$$p_{\text{esc}}=1\underset{x}{}p_{ax}v_x.$$
Now we have seen that the effective conductance between $`a`$ and $`b`$ is
$$C_ap_{\text{esc}}.$$
We wish to show that this increases whenever one of the conductances $`C_{rs}`$ is increased. If $`a`$ is different from $`r`$ or $`s`$, we need only show that $`p_{\text{esc}}`$ increases. The case where $`r`$ or $`s`$ coincides with $`a`$ is easily disposed of (see Exercise 1.4.4). The case where $`r`$ or $`s`$ coincides with $`b`$ is also easy (see Exercise 1.4.5). Hence from now on we will assume that $`r,sa`$ and $`r,sb`$.
Instead of increasing $`C_{rs}`$, we can think of adding a new edge of conductance $`ϵ`$ between $`r`$ and $`s`$. (See Figure 34.)
We will call this new edge a “bridge” to distinguish it from the other edges. Note that the graph with the bridge added will have more than one edge between $`r`$ and $`s`$ (unless there was no edge between $`r`$ and $`s`$ in the original graph), and this will complicate any expression that involves summing over edges. Everything we have said or will say holds for graphs with “multiple edges” as well as for graphs without them. So far, we have chosen to keep our expressions simple by assuming that an edge is determined by its endpoints. The trade-off is that in the manipulations below, whenever we write a sum over edges we will have to add an extra term to account for the bridge.
Why should adding the bridge increase the escape probability? The first thing you think is, “Of course, it opens up new possibilities of escaping!” The next instant you think, “Wait a minute, it also opens up new possibilities of returning to the starting point. What ensures that the first effect will outweigh the second?” As we shall see, the proper reply is, “Because the walker will cross the bridge more often in the good direction than in the bad direction.” To turn this blithe reply into a real explanation will require a little work, however.
To begin to make sense of the notion that the bridge gets used more often in the good direction than the bad, we will make a preliminary argument that applies to any edge of any graph. Let $`G`$ be any graph, and let $`rs`$ be any edge with endpoints not $`a`$ or $`b`$. $`v_r>v_s`$. Since the walker has a better chance to escape from $`s`$ than from $`r`$, this means that to cross this edge in the good direction is to go from $`r`$ to $`s`$ We shall show that the walker will cross the edge from $`r`$ to $`s`$ more often on the average than from $`s`$ to $`r`$.
Let $`u_x`$ be the expected number of times the walker is at $`x`$ and $`u_{xy}`$ the expected number of times he crosses the edge $`xy`$ from $`x`$ to $`y`$ before he reaches $`b`$ or returns to $`a`$. The calculation carried out in Section 1.3.3 shows that $`u_x/C_x`$ is harmonic for $`xa,b`$ with $`u_a/C_a=1/C_a`$ and $`b_b/C_b=0`$. But the function $`v_x/C_a`$ also has these properties, so by the Uniqueness Principle
$$\frac{u_x}{C_x}=\frac{v_x}{C_a}.$$
Now
$$u_{rs}=u_rP_{rs}=u_r\frac{C_{rs}}{C_r}=v_r\frac{C_{rs}}{C_a}$$
and
$$u_{sr}=u_sP_{sr}=u_s\frac{C_{sr}}{C_s}=v_s\frac{C_{sr}}{C_a}.$$
Since $`C_{rs}=C_{sr}`$, and since by assumption $`v_rv_s`$, this means that $`u_{rs}u_{sr}`$. Therefore, we see that any edge leads the walker more often to the more favorable of the points of the edge.
Now let’s go back and think about the graph with the bridge. The above discussion shows that the bridge helps in the sense that, on the average, the bridge is crossed more often in the direction that improves the chance of escaping. While this fact is suggestive, it doesn’t tell us that we are more likely to escape than if the bridge weren’t there; it only tells us what goes on once the bridge is in place. What we need is to make a “before and after” comparison.
Recall that we are denoting the conductance of the bridge by $`ϵ`$. To distinguish the quantities pertaining to the walks with and without the bridge, we will put $`(ϵ)`$ superscripts on quantities that refer to the walk with the bridge, so that, e.g., $`p_{\text{esc}}^{(ϵ)}`$ denotes the escape probability with the bridge.
Now let $`d^{(ϵ)}`$ denote the expected net number of times the walker crosses from $`r`$ to $`s`$. As above, we have
$$d^{(ϵ)}=u_r^{(ϵ)}\frac{ϵ}{C_r+ϵ}u_s^{(ϵ)}\frac{ϵ}{C_s+ϵ}=\left(\frac{u_r^{(ϵ)}}{C_r+ϵ}\frac{u_s^{(ϵ)}}{C_s+ϵ}\right)ϵ.$$
Claim.
$$p_{\text{esc}}^{(ϵ)}=p_{\text{esc}}+(v_rv_s)d^{(ϵ)}.$$
Why. Every time you use the bridge to go from $`r`$ to $`s`$, you improve your chances of escaping by
$$(1v_s)(1v_r)=v_rv_s$$
assuming that you would continue your walk without using the bridge. To get the probability of escaping with the bridge, you take the probability of escaping without the bridge, and correct it by adding in the change attributable to the bridge, which is the difference in the original escape probabilities at the ends of the bridge, multiplied by the net number of times you expect to cross the bridge. $`\mathrm{}`$
Proof. Suppose you’re playing a game where you walk around the graph with the bridge, and your fortune when you’re at $`x`$ is $`v_x`$, which is the probability that you would return to $`a`$ before reaching $`b`$ if the bridge weren’t there. You start at $`a`$, and walk until you reach $`b`$ or return to $`a`$.
This is not a fair game. Your initial fortune is 1 since you start at $`a`$ and $`v_a=1`$. Your expected final fortune is
$$1(1p_{\text{esc}}^{(ϵ)})+0p_{\text{esc}}^{(ϵ)}=1p_{\text{esc}}^{(ϵ)}.$$
The amount you expect to lose by participating in the game is
$$p_{\text{esc}}^{(ϵ)}.$$
(Note that escaping has suddenly become a bad thing!)
Let’s analyze where it is that you expect to lose money. First of all, you lose money when you take the first step away from $`a`$. The amount you expect to lose is
$$1\underset{x}{}P_{ax}^{(ϵ)}v_x=p_{\text{esc}}.$$
Now if your fortune were given by $`v_x^{(ϵ)}`$ instead of $`v_x`$, the game would be fair after this first step. However, the function $`v_x`$ is not harmonic for the walk with the bridge; it fails to be harmonic at $`r`$ and $`s`$. Every time you step away from $`r`$, you expect to lose an amount
$$v_r\left(\underset{x}{}\frac{C_{rs}}{C_r+ϵ}v_x+\frac{ϵ}{C_r+ϵ}v_s\right)=(v_rv_s)\frac{ϵ}{C_r+ϵ}.$$
Similarly, every time you step away from $`s`$ you expect to lose an amount
$$(v_sv_r)\frac{ϵ}{C_s+ϵ}.$$
The total amount you expect to lose by participating in the game is:
$`\text{expected loss at first step}+`$
$`(\text{expected loss at }r)(\text{expected number of times at }r)+`$
$`(\text{expected loss at }s)(\text{expected number of times at }s)`$
$`=`$ $`p_{\text{esc}}+`$
$`(v_rv_s){\displaystyle \frac{ϵ}{C_r+ϵ}}u_r^{(ϵ)}+`$
$`(v_sv_r){\displaystyle \frac{ϵ}{C_s+ϵ}}u_s^{(ϵ)}`$
$`=`$ $`p_{\text{esc}}+(v_rv_s)d^{(ϵ)}.`$
Equating this with our first expression for the expected cost of playing the game yields the formula we were trying to prove.
According to the formula just established,
$$p_{\text{esc}}^{(ϵ)}p_{\text{esc}}=(v_rv_s)\left(\frac{u_r^{(ϵ)}}{C_r+ϵ}\frac{u_s^{(ϵ)}}{C_s+ϵ}\right)ϵ.$$
For small $`ϵ`$, we will have
$$\frac{u_r^{(ϵ)}}{C_r+ϵ}\frac{u_s^{(ϵ)}}{C_s+ϵ}\frac{u_r}{C_r}\frac{u_s}{C_s}=\frac{v_r}{C_a}\frac{v_s}{C_a},$$
so for small $`ϵ`$
$$p_{\text{esc}}^{(ϵ)}p_{\text{esc}}(v_rv_s)^2\frac{ϵ}{C_a}.$$
This approximation allows us to conclude that
$$p_{\text{esc}}^{(ϵ)}p_{\text{esc}}0$$
for small $`ϵ`$. But this is enough to establish the monotonicity law, since any finite change in $`ϵ`$ can be realized by making an infinite chain of graphs each of which is obtained from the one before by adding a bridge of infinitesimal conductance.
To recapitulate, the difference in the escape probabilities with and without the bridge is obtained by taking the difference between the original escape probabilities at the ends of the bridge, and multiplying by the expected net number of crossings of the bridge. This quantity is positive because the walker tends to cross the bridge more often in the good direction than in the bad direction.
###### Exercise 1.4.4
Give a probabilistic argument to show that $`C_ap_{\text{esc}}`$ increases with $`C_{ar}`$ for any $`r`$. Give an example to show that $`p_{\text{esc}}`$ by itself may actually decrease.
###### Exercise 1.4.5
Give a probabilistic argument to show that $`C_ap_{\text{esc}}`$ increases with $`C_{rb}`$ for any $`r`$.
###### Exercise 1.4.6
Show that when $`v_r=v_s`$, changing the value of $`C_{rs}`$ does not change $`p_{\text{esc}}`$.
###### Exercise 1.4.7
Show that
$$\frac{}{R_{rs}}R_{\text{eff}}=i_{rs}^2.$$
###### Exercise 1.4.8
In this exercise we ask you to derive an exact formula for the change in escape probability
$$p_{\text{esc}}^{(ϵ)}p_{\text{esc}},$$
in terms of quantities that refer only to the walk without the bridge.
(a) Let $`N_{xy}`$ denote the expected number of times in state $`y`$ for a walker who starts at $`x`$ and walks around the graph without the bridge until he reaches $`a`$ or $`b`$. It is a fact that
$$u_r^{(ϵ)}=u_r+u_r^{(ϵ)}\frac{ϵ}{C_r+ϵ}(N_{sr}+1N_{rr})+u_s^{(ϵ)}\frac{ϵ}{C_s+ϵ}(N_{rr}N_{sr}).$$
Explain in words why this formula is true.
(b) This equation for $`u_r^{(ϵ)}`$ can be rewritten as follows:
$$\frac{C_r}{C_r+ϵ}u_r^{(ϵ)}=u_r+d^{(ϵ)}(N_{sr}N_{rr}).$$
Prove this formula. (Hint: Consider a game where your fortune at $`x`$ is $`N_{xr}`$, and where you start from $`a`$ and walk on the graph with the bridge until you reach $`b`$ or return to $`a`$.)
(c) Write down the corresponding formula for $`u_s^{(ϵ)}`$, and use this formula to get an expression for $`d^{(ϵ)}`$ in terms of quantities that refer to the walk without the bridge.
(d) Use the expression for $`d^{(ϵ)}`$ to express $`p_{\text{esc}}^{(ϵ)}p_{\text{esc}}`$ in terms of quantities that refer to the walk without the bridge, and verify that the value of your expression is $`0`$ for $`ϵ0`$.
###### Exercise 1.4.9
Give a probabilistic interpretation of the energy dissipation rate.
#### 1.4.3 A Markov chain proof of the Monotonicity Law
Let $`𝐏`$ be the ergodic Markov chain associated with an electric network. When we add an $`ϵ`$ bridge from $`r`$ to $`s`$, we obtain a new transition matrix $`𝐏^{(ϵ)}`$ that differs from $`𝐏`$ only for transitions from $`r`$ and $`s`$. We can minimize the differences between $`𝐏`$ and $`𝐏^{(ϵ)}`$ by replacing $`𝐏`$ by the matrix $`\widehat{𝐏}`$ corresponding to the circuit without the bridge but with an $`ϵ`$ conductance added from $`r`$ to $`r`$ and from $`s`$ to $`s`$. This allows the chain to stay in states $`r`$ and $`s`$ but does not change the escape probability from $`a`$ to $`b`$. Thus, we can compare the escape probabilities for the two matrices $`\widehat{𝐏}`$ and $`𝐏^{(ϵ)}`$, which differ only by
$$\begin{array}{cc}\widehat{P}_{rr}=\frac{ϵ}{C_r+ϵ}\hfill & P_{rr}^{(ϵ)}=0\hfill \\ \widehat{P}_{rs}=\frac{C_{rs}}{C_r+ϵ}\hfill & P_{rs}^{(ϵ)}=\frac{C_{rs}+ϵ}{C_r+ϵ}\hfill \\ \widehat{P}_{ss}=\frac{ϵ}{C_s+ϵ}\hfill & P_{ss}^{(ϵ)}=0\hfill \\ \widehat{P}_{sr}=\frac{C_{sr}}{C_S+ϵ}\hfill & P_{sr}^{(ϵ)}=\frac{C_{sr}+ϵ}{C_s+ϵ}\hfill \end{array}.$$
We make states $`a`$ and $`b`$ into absorbing states. Let $`\widehat{𝐍}`$ and $`𝐍^{(ϵ)}`$ be the fundamental matrices for the absorbing chains obtained from $`\widehat{𝐏}`$ and $`𝐏^{(ϵ)}`$ respectively. Then $`\widehat{𝐍}=(𝐈\widehat{𝐐})^1`$ and $`𝐍^{(ϵ)}=(𝐈𝐐^{(ϵ)})^1`$ where $`\widehat{𝐐}`$ and $`𝐐^{(ϵ)}`$ differ only for the four components involving only $`r`$ and $`s`$. That is,
𝐐(ϵ)=𝐐^+(
rsr-ϵ+Crϵϵ+Crϵsϵ+Csϵ-ϵ+Csϵ
)=𝐐^+𝐡𝐤superscript𝐐italic-ϵ^𝐐
rsr-ϵ+Crϵϵ+Crϵsϵ+Csϵ-ϵ+Csϵ
^𝐐𝐡𝐤{\mathbf{Q^{(\epsilon)}}}={\mathbf{\hat{Q}}}+\hbox{}\;\vbox{\kern 34.6222pt\hbox{$\kern 57.82578pt\kern-8.75pt\left(\kern-57.82578pt\vbox{\kern-34.6222pt\vbox{\halign{$#$\hfil\kern 2\p@\kern\@tempdima&\thinspace\hfil$#$\hfil&&\quad\hfil$#$\hfil\cr\hfil\crcr\kern-12.0pt\cr$\hfil\kern 2.0pt\kern 8.75pt&r&s\crcr\kern 2.0pt\cr r$\hfil\kern 2.0pt\kern 8.75pt&-\frac{\epsilon}{C_{r}+\epsilon}&\frac{\epsilon}{C_{r}+\epsilon}\cr s$\hfil\kern 2.0pt\kern 8.75pt&\frac{\epsilon}{C_{s}+\epsilon}&-\frac{\epsilon}{C_{s}+\epsilon}\crcr\cr}}\kern-12.0pt}\,\right)$}}={\mathbf{\hat{Q}}}+{\mathbf{h}}{\mathbf{k}}
where $`𝐡`$ is the column vector with only components $`r`$ and $`s`$ non-zero
𝐡=(
rϵ+Crϵs-ϵ+Csϵ
)𝐡
rϵ+Crϵs-ϵ+Csϵ
{\mathbf{h}}=\hbox{}\;\vbox{\kern 30.31665pt\hbox{$\kern 28.78717pt\kern-8.75pt\left(\kern-28.78717pt\vbox{\kern-30.31665pt\vbox{\halign{$#$\hfil\kern 2\p@\kern\@tempdima&\thinspace\hfil$#$\hfil&&\quad\hfil$#$\hfil\cr\hfil\crcr\kern-12.0pt\cr$\hfil\kern 2.0pt\kern 8.75pt&\ \crcr\kern 2.0pt\cr r$\hfil\kern 2.0pt\kern 8.75pt&\frac{\epsilon}{C_{r}+\epsilon}\cr s$\hfil\kern 2.0pt\kern 8.75pt&-\frac{\epsilon}{C_{s}+\epsilon}\crcr\cr}}\kern-12.0pt}\,\right)$}}
and $`𝐤`$ is a row vector with only components $`r`$ and $`s`$ non-zero
𝐤=(
rs−11
).𝐤
rs−11
{\mathbf{k}}=\hbox{}\;\vbox{\kern 23.74998pt\hbox{$\kern 30.5555pt\kern-8.75pt\left(\kern-30.5555pt\vbox{\kern-23.74998pt\vbox{\halign{$#$\hfil\kern 2\p@\kern\@tempdima&\thinspace\hfil$#$\hfil&&\quad\hfil$#$\hfil\cr\hfil\crcr\kern-12.0pt\cr$\hfil\kern 2.0pt\kern 8.75pt&r&s\crcr\kern 2.0pt\cr$\hfil\kern 2.0pt\kern 8.75pt&-1&1\crcr\cr}}\kern-12.0pt}\,\right)$}}.
J. G. Kemeny has pointed out to us that if $`𝐀`$ is any matrix with inverse $`𝐍`$ and we add to $`𝐀`$ a matrix of the form $`\mathrm{𝐡𝐤}`$, then $`𝐀\mathrm{𝐡𝐤}`$ has an inverse if and only if $`\mathrm{𝐤𝐍𝐡}1`$ and, if so, $`\overline{𝐍}=(𝐀\mathrm{𝐡𝐤})^1`$ is given by
$$\overline{𝐍}=𝐍+c(\mathrm{𝐍𝐡})(\mathrm{𝐤𝐍})$$
where $`c=1/(1\mathrm{𝐤𝐍𝐡})`$. You are asked to prove this in Exercise 1.4.10. Adding $`\mathrm{𝐡𝐤}`$ to $`𝐀=𝐈\widehat{𝐐}`$ and using this result, we obtain
$$𝐍^{(ϵ)}=\widehat{𝐍}+c(\widehat{𝐍}𝐡)(𝐤\widehat{𝐍}).$$
Using the simple nature of $`𝐡`$ and $`𝐤`$ we obtain
$$N_{ij}^{(ϵ)}=\widehat{N}_{ij}+\left(\frac{\widehat{N}_{ir}ϵ}{C_r+ϵ}\frac{\widehat{N}_{is}ϵ}{C_s+ϵ}\right)(\widehat{N}_{sj}\widehat{N}_{rj})$$
and
$$c=\frac{1}{1+\frac{\widehat{N}_{rr}ϵ}{C_r+ϵ}\frac{\widehat{N}_{sr}ϵ}{C_r+ϵ}+\frac{\widehat{N}_{ss}ϵ}{C_s+ϵ}\frac{\widehat{N}_{rs}ϵ}{C_s+ϵ}}.$$
Since $`\widehat{N}_{rr}`$ is the expected number of times in $`r`$ starting in $`r`$ and $`\widehat{N}_{sr}`$ is the expected number of times in $`r`$ starting in $`s`$, $`\widehat{N}_{rr}\widehat{N}_{sr}`$. Similarly $`\widehat{N}_{ss}\widehat{N}_{rs}`$ and so the denominator of $`c`$ is $`1`$. In particular, it is positive.
Recall that the absorption probabilities for state $`b`$ are given by
$$B_{xb}=\underset{y}{}N_{xy}P_{yb}.$$
Since $`P_{xb}^{(ϵ)}=\widehat{P}_{xb}`$,
$$B_{xb}^{(ϵ)}=\widehat{B}_{xb}+c\left(\frac{\widehat{N}_{xr}ϵ}{C_r+ϵ}\frac{\widehat{N}_{xs}ϵ}{C_s+ϵ}\right)(\widehat{B}_{sb}\widehat{B}_{rb}).$$
Since $`P_{ax}^{(ϵ)}=\widehat{P}_{ax}`$,
$$p_{\text{esc}}^{(ϵ)}=\widehat{p}_{\text{esc}}+c\left(\frac{\widehat{u}_rϵ}{C_r+ϵ}\frac{\widehat{u}_sϵ}{C_s+ϵ}\right)(\widehat{B}_{sb}\widehat{B}_{rb})$$
where $`\widehat{u}_x`$ is the expected number of times that the ergodic chain $`\widehat{P}`$, started at state $`a`$, is in state $`x`$ before returning to $`a`$ reaching $`b`$ for the first time. The absorption probability $`B_{xa}`$ is the quantity $`v_x`$ introduced in the previous section. As shown there, reversibility allows us to conclude that
$$\frac{\widehat{u}_x}{\widehat{C}_x}=\frac{\widehat{B}_{xa}}{\widehat{C}_a}=\frac{\widehat{B}_{xa}}{C_a}$$
so that
$$p_{\text{esc}}^{(ϵ)}=p_{\text{esc}}+\frac{ϵc}{C_a}(\widehat{B}_{sb}\widehat{B}_{rb})^2$$
and this shows that the Monotonicity Law is true.
The change from $`𝐏`$ to $`\widehat{𝐏}`$ was merely to make the calculations easier. As we have remarked, the escape probabilities are the same for the two chains as are the absorption probabilities $`B_{ib}`$. Thus we can remove the hats and write the same formula.
$$p_{\text{esc}}^{(ϵ)}=p_{\text{esc}}+\frac{ϵc}{C_a}(B_{sb}B_{rb})^2.$$
The only quantity in this final expression that seems to depend upon quantities from $`\widehat{𝐏}`$ is $`c`$. In Exercise 1.4.11 you are asked to show that $`c`$ can also be expressed in terms of the fundamental matrix $`𝐍`$ obtained from the original $`𝐏`$.
###### Exercise 1.4.10
Let $`𝐀`$ be a matrix with inverse $`𝐍=𝐀^1`$ Let $`𝐡`$ be a column vector and $`𝐤`$ a row vector. Show that
$$\overline{𝐍}=(𝐀\mathrm{𝐡𝐤})^1$$
exists if and only if $`\mathrm{𝐤𝐍𝐡}1`$ and, if so,
$$\overline{𝐍}=𝐍+\frac{(\mathrm{𝐍𝐡})(\mathrm{𝐤𝐍})}{1\mathrm{𝐤𝐍𝐡}}.$$
###### Exercise 1.4.11
Show that $`c`$ can be expressed in terms of the fundamental matrix $`𝐍`$ of the original Markov chain $`𝐏`$ by
$$c=\frac{1}{1+\frac{N_{rr}ϵ}{C_r}\frac{N_{sr}ϵ}{C_r}+\frac{N_{ss}ϵ}{C_s}\frac{N_{rs}ϵ}{C_s}}.$$
## 2 Random walks on infinite networks
### 2.1 Polya’s recurrence problem
#### 2.1.1 Random walks on lattices
In 1921 George Polya investigated random walks on certain infinite graphs, or as he called them, “street networks”. The graphs he considered, which we will refer to as lattices, are illustrated in Figure 35.
To construct a $`d`$-dimensional lattice, we take as vertices those points $`(x_1,\mathrm{},x_d)`$ of $`𝐑^d`$ all of whose coordinates are integers, and we join each vertex by an undirected line segment to each of its $`2d`$ nearest neighbors. These connecting segments, which represent the edges of our graph, each have unit length and run parallel to one of the coordinate axes of $`𝐑^d`$. We will denote this $`d`$-dimensional lattice by $`𝐙^d`$. We will denote the origin $`(0,0,\mathrm{},0)`$ by $`\mathrm{𝟎}`$.
Now let a point walk around at random on this lattice. As usual, by walking at random we mean that, upon reaching any vertex of the graph, the probability of choosing any one of the $`2d`$ edges leading out of that vertex is $`\frac{1}{2d}`$. We will call this random walk *simple random walk* in $`d`$ dimensions.
When $`d=1`$, our lattice is just an infinite line divided into segments of length one. We may think of the vertices of this graph as representing the fortune of a gambler betting on heads or tails in a fair coin tossing game. Simple random walk in one dimension then represents the vicissitudes of his or her fortune, either increasing or decreasing by one unit after each round of the game.
When $`d=2`$, our lattice looks like an infinite network of streets and avenues, which is why we describe the random motion of the wandering point as a “walk”.
When $`d=3`$, the lattice looks like an infinite “jungle gym”, so perhaps in this case we ought to talk about a “random climb”, but we will not do so. It is worth noting that when $`d=3`$, the wanderings of our point could be regarded as an approximate representation of the random path of a molecule diffusing in an infinite cubical crystal. Figure 36 shows a simulation of a simple random walk in three dimensions.
#### 2.1.2 The question of recurrence
The question that Polya posed amounts to this: “Is the wandering point certain to return to its starting point during the course of its wanderings?” If so, we say that the walk is *recurrent*. If not, that is, if there is a positive probability that the point will never return to its starting point, then we say that the walk is *transient*.
If we denote the probability that the point never returns to its starting point by $`p_{\text{esc}}`$, then the chain is recurrent if $`p_{\text{esc}}=0`$, and transient if $`p_{\text{esc}}>0`$.
We will call the problem of determining recurrence or transience of a random walk the *type problem*.
#### 2.1.3 Polya’s original question
The definition of recurrence that we have given differs from Polya’s original definition. Polya defined a walk to be recurrent if, with probability one, it will pass through every single point of the lattice in the course of its wanderings. In our definition, we require only that the point return to its starting point. So we have to ask ourselves, “Can the random walk be recurrent in our sense and fail to be recurrent in Polya’s sense?”
The answer to this question is, “No, the two definitions of recurrence are equivalent.” Why? Because if the point must return once to its starting point, then it must return there again and again, and each time it starts away from the origin, it has a certain non-zero probability of hitting a specified target vertex before returning to the origin. And anyone can get a bull’s-eye if he or she is allowed an infinite number of darts, so eventually the point will hit the target vertex.
###### Exercise 2.1.1
Write out a rigorous version of the argument just given.
#### 2.1.4 Polya’s Theorem: recurrence in the plane, transience in space
In , Polya proved the following theorem:
Polya’s Theorem. Simple random walk on a $`d`$-dimensional lattice is recurrent for $`d=1,2`$ and transient for $`d>2`$.
The rest of this work will be devoted to trying to understand this theorem. Our approach will be to exploit once more the connections between questions about a random walk on a graph and questions about electric currents in a corresponding network of resistors. We hope that this approach, by calling on methods that appeal to our physical intuition, will leave us feeling that we know “why” the theorem is true.
###### Exercise 2.1.2
Show that Polya’s theorem implies that if two random walkers start at $`\mathrm{𝟎}`$ and wander independently, then in one and two dimensions they will eventually meet again, but in three dimensions there is positive probability that they won’t.
###### Exercise 2.1.3
Show that Polya’s theorem implies that a random walker in three dimensions will eventually hit the line defined by $`x=2,z=0`$.
#### 2.1.5 The escape probability as a limit of escape probabilities for finite graphs
We can determine the type of an infinite lattice from properties of bigger and bigger finite graphs that sit inside it. The simplest way to go about this is to look at the lattice analog of balls (solid spheres) in space. These are defined as follows: Let $`r`$ be an integer—this will be the radius of the ball. Let $`G^{(r)}`$ be the graph gotten from $`𝐙^d`$ by throwing out vertices whose distance from the origin is $`>r`$. By “distance from the origin” we mean here not the usual Euclidean distance, but the distance “in the lattice”; that is, the length of the shortest path along the edges of the lattice between the two points. Let $`G^{(r)}`$ be the “sphere” of radius $`r`$ about the origin, i.e., those points that are exactly $`r`$ units from the origin. In two dimensions, $`G^{(r)}`$ looks like a square. (See Figure 37.)
In three dimensions, it looks like an octahedron.
We define a random walk on $`G^{(r)}`$ as follows: The walk starts at $`\mathrm{𝟎}`$ and continues as it would on $`𝐙^d`$ until it reaches a point on $`G^{(r)}`$ and then it stays at this point. Thus the walk on $`G^{(r)}`$ is an absorbing Markov chain with every point of $`G^{(r)}`$ an absorbing state.
Let $`p_{\text{esc}}^{(r)}`$ be the probability that a random walk on $`G^{(r)}`$ starting at $`\mathrm{𝟎}`$, reaches $`G^{(r)}`$ before returning to $`\mathrm{𝟎}`$. Then $`p_{\text{esc}}^{(r)}`$ decreases as $`r`$ increases and $`p_{\text{esc}}=lim_r\mathrm{}p_{\text{esc}}^{(r)}`$ is the *escape probability* for the infinite graph. If this limit is 0, the infinite walk is recurrent. If it is greater than 0, the walk is transient.
#### 2.1.6 Electrical formulation of the type problem
Now that we have expressed things in terms of finite graphs, we can make use of electrical methods. To determine $`p_{\text{esc}}`$ electrically, we simply ground all the points of $`G^{(r)}`$, maintain $`\mathrm{𝟎}`$ at one volt, and measure the current $`i^{(r)}`$ flowing into the circuit. (See Figure 38.)
From Section 1.3.4, we have
$$p_{\text{esc}}^{(r)}=\frac{i^{(r)}}{2d},$$
where $`d`$ is the dimension of the lattice. (Remember that we have to divide by the number of branches coming out of the starting point.) Since the voltage being applied is 1, $`i^{(r)}`$ is just the effective conductance between $`\mathrm{𝟎}`$ and $`G^{(r)}`$, i.e.,
$$i^{(r)}=\frac{1}{R_{\text{eff}}^{(r)}}.$$
where $`R_{\text{eff}}^{(r)}`$ is the effective resistance from $`\mathrm{𝟎}`$ to $`G^{(r)}`$. Thus
$$p_{\text{esc}}^{(r)}=\frac{1}{2dR_{\text{eff}}^{(r)}}.$$
Define $`R_{\text{eff}}`$, the *effective resistance from the origin to infinity*, to be
$$R_{\text{eff}}=\underset{r\mathrm{}}{lim}R_{\text{eff}}^{(r)}.$$
This limit exists since $`R_{\text{eff}}^{(r)}`$ is an increasing function of $`r`$. Then
$$p_{\text{esc}}=\frac{1}{2dR_{\text{eff}}}.$$
Of course $`R_{\text{eff}}`$ may be infinite; in fact, this will be the case if and only if $`p_{\text{esc}}=0`$. Thus the walk is recurrent if and only if the resistance to infinity is infinite, which makes sense.
The success of this electrical formulation of the type problem will come from the fact that the resistance to infinity can be estimated using classical methods of electrical theory.
#### 2.1.7 One Dimension is easy, but what about higher dimensions?
We now know that simple random walk on a graph is recurrent if and only if a corresponding network of 1-ohm resistors has infinite resistance “out to infinity”. Since an infinite line of resistors obviously has infinite resistance, it follows that simple random walk on the 1-dimensional lattice is recurrent, as stated by Polya’s theorem.
What happens in higher dimensions? We are asked to decide whether a $`d`$-dimensional lattice has infinite resistance to infinity. The difficulty is that the $`d`$-dimensional lattice $`𝐙^d`$ lacks the rotational symmetry of the Euclidean space $`𝐑^d`$ in which it sits.
To see how this lack of symmetry complicates electrical problems, we determine, by solving the appropriate discrete Dirichlet problem, the voltages for a one-volt battery attached between $`\mathrm{𝟎}`$ and the points of $`G^{(3)}`$ in $`𝐙^2`$. The resulting voltages are:
$$\begin{array}{cccc}& & & 0\\ & & 0& .091& 0\\ & 0& .182& .364& .182& 0\\ 0& .091& .364& 1& .364& .091& 0\\ & 0& .182& .364& .182& 0\\ & & 0& .091& 0\\ & & & 0\end{array}.$$
The voltages at points of $`G^{(1)}`$ are equal, but the voltages at points of $`G^{(2)}`$ are not. This means that the resistance from $`\mathrm{𝟎}`$ to $`G^{(3)}`$ cannot be written simply as the sum of the resistances from $`\mathrm{𝟎}`$ to $`G^{(1)}`$, $`G^{(1)}`$ to $`G^{(2)}`$, and $`G^{(2)}`$ to $`G^{(3)}`$. This is in marked contrast to the case of a continuous resistive medium to be discussed in Section 2.1.8.
###### Exercise 2.1.4
Using the voltages given for $`G^{(3)}`$, find $`R_{\text{eff}}^{(3)}`$ and $`p_{\text{esc}}^{(3)}`$.
###### Exercise 2.1.5
Consider a one-dimensional infinite network with resistors $`R_n=1/2^n`$ from $`n`$ to $`n+1`$ for $`n=\mathrm{},2,1,0,1,2,\mathrm{}`$. Describe the associated random walk and determine whether the walk is recurrent or transient.
###### Exercise 2.1.6
A random walk moves on the non-negative integers; when it is in state $`n`$, $`n>0`$, it moves with probability $`p_n`$ to $`n+1`$ and with probability $`1p_n`$, to $`n1`$. When at $`\mathrm{𝟎}`$, it moves to 1. Determine a network that gives this random walk and give a criterion in terms of the $`p_n`$ for recurrence of the random walk.
#### 2.1.8 Getting around the lack of rotational symmetry of the lattice
Suppose we replace our $`d`$-dimensional resistor lattice by a (homogeneous, isotropic) resistive medium filling all of $`𝐑^d`$ and ask for the effective resistance to infinity. Naturally we expect that the rotational symmetry will make this continuous problem easier to solve than the original discrete problem. If we took this problem to a physicist, he or she would probably produce something like the scribblings illustrated in Figure 41, and conclude that the effective resistance is infinite for $`d=1,2`$ and finite for $`d>2`$. The analogy to Polya’s theorem is obvious, but is it possible to translate these calculations for continuous media into information about what happens in the lattice?
This can indeed be done, and this would certainly be the most natural approach to take. We will come back to this approach at the end of the work. For now, we will take a different approach to getting around the asymmetry of the lattice. Our method will be to modify the lattice in such a way as to obtain a graph that is symmetrical enough so that we can calculate its resistance out to infinity. Of course, we will have to think carefully about what happens to that resistance when we make these modifications.
#### 2.1.9 Rayleigh: shorting shows recurrence in the plane, cutting shows transience in space
Here is a sketch of the method we will use to prove Polya’s theorem.
To take care of the case $`d=2`$, we will modify the 2-dimensional resistor network by shorting certain sets of nodes together so as to get a new network whose resistance is readily seen to be infinite. As shorting can only decrease the effective resistance of the network, the resistance of the original network must also be infinite. Thus the walk is recurrent when $`d=2`$.
To take care of the case $`d=3`$, we will modify the 3-dimensional network by cutting out certain of the resistors so as to get a new network whose resistance is readily seen to be finite. As cutting can only increase the resistance of the network, the resistance of the original network must also be finite. Thus the walk is transient when $`d=3`$.
The method of applying shorting and cutting to get lower and upper bounds for the resistance of a resistive medium was introduced by Lord Rayleigh. (Rayleigh ; see also Maxwell , Jeans , PoIya and Szego ). We will refer to Rayleigh’s techniques collectively as *Rayleigh’s short-cut method*. This does not do Rayleigh justice, for Rayleigh’s method is a whole bag of tricks that goes beyond mere shorting and cutting—but who can resist a pun?
Rayleigh’s method was apparently first applied to random walks by C. St. J. A. Nash-Williams , who used the shorting method to establish recurrence for random walks on the 2-dimensional lattice.
### 2.2 Rayleigh’s short-cut method
#### 2.2.1 Shorting and cutting
In its simplest form, Rayleigh’s method involves modifying the network whose resistance we are interested in so as to get a simpler network. We consider two kinds of modifications, shorting and cutting. Cutting involves nothing more than clipping some of the branches of the network, or what is the same, simply deleting them from the network. Shorting involves connecting a given set of nodes together with perfectly conducting wires, so that current can pass freely between them. In the resulting network, the nodes that were shorted together behave as if they were a single node.
#### 2.2.2 The Shorting Law and the Cutting Law; Rayleigh’s idea
The usefulness of these two procedures (shorting and cutting) stems from the following observations:
Shorting Law. Shorting certain sets of nodes together can only decrease the effective resistance of the network between two given nodes.
Cutting Law. Cutting certain branches can only increase the effective resistance between two given nodes.
These laws are both equivalent to Rayleigh’s Monotonicity Law, which was introduced in Section 1.4.1 (see Exercise 2.2.1):
Monotonicity Law. The effective resistance between two given nodes is monotonic in the branch resistances.
Rayleigh’s idea was to use the Shorting Law and the Cutting Law above to get lower and upper bounds for the resistance of a network. In Section 2.2.3 we apply this method to solve the recurrence problem for simple random walk in dimensions 2 and 3.
###### Exercise 2.2.1
Show that the Shorting Law and the Cutting Law are both equivalent to the Monotonicity Law.
#### 2.2.3 The plane is easy
When $`d=2`$, we apply the Shorting Law as follows: Short together nodes on squares about the origin, as shown in Figure 43.
The network we obtain is equivalent to the network shown in Figure 44.
Now as $`n`$ 1-ohm resistors in parallel are equivalent to a single resistor of resistance $`\frac{1}{n}`$ ohms, the modified network is equivalent to the network shown in Figure 45.
The resistance of this network out to infinity is
$$\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{8n+4}=\mathrm{}.$$
As the resistance of the old network can only be bigger, we conclude that it too must be infinite, so that the walk is recurrent when $`d=2`$.
###### Exercise 2.2.2
Using the shorting technique, give an upper bound for $`p_{\text{esc}}^{(3)}`$, and compare this with the exact value obtained in Exercise 2.1.4.
#### 2.2.4 Space: searching for a residual network
When $`d=3`$, what we want to do is delete certain of the branches of the network so as to leave behind a residual network having manifestly finite resistance. The problem is to reconcile the “manifestly” with the “finite”. We want to cut out enough edges so that the effective resistance of what is left is easy to calculate, while leaving behind enough edges so that the result of the calculation is finite.
#### 2.2.5 Trees are easy to analyze
Trees—that is, graphs without circuits—are undoubtedly the easiest to work with. For instance, consider the full binary tree, shown in Figure 46.
Notice that sitting inside this tree just above the root are two copies of the tree itself. This self-similarity property can be used to compute the effective resistance $`R_{\mathrm{}}`$ from the root out to infinity. (See Exercise 2.2.3.) It turns out that $`R_{\mathrm{}}=1`$. We will demonstrate this below by a more direct method.
To begin with, let us determine the effective resistance $`R_n`$ between the root and the set of $`n`$th generation branch points. To do this, we should ground the set of branch points, hook the root up to a 1-volt battery, and compute
$$R_n=\frac{1}{\text{current through battery}}.$$
For $`n=3`$, the circuit that we would obtain is shown in Figure 48.
In the resulting circuit, all branch points of the same generation are at the same voltage (by symmetry). Nothing happens when you short together nodes that are at the same potential. Thus shorting together branch points of the same generation will not affect the distribution of currents in the branches. In particular, this modification will not affect the current through the battery, and we conclude that
$$R_n=\frac{1}{\text{current in original circuit}}=\frac{1}{\text{current in modified circuit}}.$$
For $`n=3`$, the modified circuit is shown in Figure 49.
This picture shows that
$$R_3=\frac{1}{2}+\frac{1}{4}+\frac{1}{8}=1\frac{1}{2^3}.$$
More generally,
$$R_n=\frac{1}{2}+\frac{1}{4}+\mathrm{}+\frac{1}{2^n}=1\frac{1}{2^n}.$$
Letting $`n\mathrm{}`$, we get
$$R_{\mathrm{}}=\underset{n\mathrm{}}{lim}R_n=\underset{n\mathrm{}}{lim}1\frac{1}{2^n}=1.$$
Figure 50 shows another closely related tree, the *tree homogeneous of degree three*: Note that all nodes of this tree are similar—there is no intrinsic way to distinguish one from another.
This tree is obviously a close relative of the full binary tree. Its resistance to infinity is $`2/3`$.
###### Exercise 2.2.3
(a) Show, using the self-similarity of the full binary tree, that the resistance $`R_{\mathrm{}}`$ to infinity satisfies the equation
$$R_{\mathrm{}}=\frac{R_{\mathrm{}}+1}{2}$$
and conclude that either $`R_{\mathrm{}}=1`$ or $`R_{\mathrm{}}=\mathrm{}`$.
(b) Again using the self-similarity of the tree, show that
$$R_{n+1}=\frac{R_n+1}{2}$$
where $`R_n`$ denotes the resistance out to the set of the $`n`$th-generation branch points. Conclude that
$$R_{\mathrm{}}=\underset{n\mathrm{}}{lim}R_n=1.$$
#### 2.2.6 The full binary tree is too big
Nothing could be nicer than the two trees we have just described. They are the prototypes of networks having manifestly finite resistance to infinity. Unfortunately, we can’t even come close to finding either of these trees as a subgraph of the three-dimensional lattice. For in these trees, the number of nodes in a “ball” of radius $`r`$ grows exponentially with $`r`$, whereas in a $`d`$-dimensional lattice, it grows like $`r`$, i.e., much slower. (See Figure 51.)
There is simply no room for these trees in any finite-dimensional lattice.
#### 2.2.7 $`\text{NT}_3`$: a “three-dimensional” tree
These observations suggest that we would do well to look for a nice tree $`\text{NT}_3`$ where the number of nodes within a radius $`r`$ of the root is on the order of $`r^3`$. For we might hope to find something resembling $`\text{NT}_3`$ in the 3-dimensional lattice, and if there is any justice in the world, this tree would have finite resistance to infinity, and we would be done.
Before introducing $`\text{NT}_3`$, let’s have a look at $`\text{NT}_2`$, our choice for the tree most likely to succeed in the 2-dimensional lattice (see Figure 52).
The idea behind $`\text{NT}_2`$ is that, since a ball of radius $`r`$ in the graph ought to contain something like $`r^2`$ points, a sphere of radius $`r`$ ought to contain something like $`r`$ points, so the number of points in a sphere should roughly double when the radius of the sphere is doubled. For this reason, we make the branches of our tree split in two every time the distance from the origin is (roughly) doubled.
Similarly, in a 3-dimensional tree, when we double the radius, the size of a sphere should roughly quadruple. Thus in $`\text{NT}_3`$, we make the branches of our tree split in four where the branches of $`\text{NT}_2`$ would have split in two. $`\text{NT}_3`$ is shown in Figure 53.
Obviously, $`\text{NT}_3`$ is none too happy about being drawn in the plane.
#### 2.2.8 $`\text{NT}_3`$ has finite resistance
To see if we’re on the right track, let’s work out the resistance of our new trees. These calculations are shown in Figures 54 and 55.
As we would hope, the resistance of $`\text{NT}_2`$ is infinite, but the resistance of $`\text{NT}_3`$ is not.
###### Exercise 2.2.4
Use self-similarity arguments to compute the resistance of $`\text{NT}_2`$ and $`\text{NT}_3`$.
#### 2.2.9 But does $`\text{NT}_3`$ fit in the three-dimensional lattice?
We would like to embed $`\text{NT}_3`$ in $`𝐙^3`$. We start by trying to embed $`\text{NT}_2`$ in $`𝐙^2`$. The result is shown in Figure 56.
To construct this picture, we start from the origin and draw 2 rays, one going north, one going east. Whenever a ray intersects the line $`x+y=2^n1`$ for some $`n`$, it splits into 2 rays, one going north, and one going east. The sequence of pictures in Figure 57 shows successively larger portions of the graph, along with the corresponding portions of $`\text{NT}_2`$.
Of course this isn’t really an embedding, since certain pairs of points that were distinct in $`\text{NT}_2`$ get identified, that is, they are made to correspond to a single point in the picture. In terms of our description, sometimes a ray going north and a ray going east pass through each other. This could have been avoided by allowing the rays to “bounce” instead of passing through each other, at the expense of embedding not $`\text{NT}_2`$ but a close relative—see Exercise 2.2.7. However, because the points of each identified pair are at the same distance from the root of $`\text{NT}_2`$, when we put a battery between the root and the $`n`$th level they will be at the same potential. Hence, the current flow is not affected by these identifications, so the identifications have no effect on $`R_{\text{eff}}`$. For our purposes, then, we have done just as well as if we had actually embedded $`\text{NT}_2`$.
To construct the analogous picture in three dimensions, we start three rays off from the origin going north, east, and up. Whenever a ray intersects the plane $`x+y+z=2^n1`$ for some $`n`$, it splits into three rays, going north, east, and up. This process is illustrated in Figure 58.
Surprisingly, the subgraph of the 3-dimensional lattice obtained in this way is not $`\text{NT}_3`$! Rather, it represents an attempt to embed the tree shown in Figure 59.
We call this tree $`\text{NT}_{2.5849\mathrm{}}`$ because it is $`2.5849\mathrm{}`$-dimensional in the sense that when you double the radius of a ball, the number of points in the ball gets multiplied roughly by 6 and
$$6=2^{\mathrm{log}_26}=2^{2.5849\mathrm{}}.$$
Again, certain pairs of points of $`\text{NT}_{2.5849\mathrm{}}`$ have been allowed to correspond to the same point in the lattice, but once again the intersections have no effect on $`R_{\text{eff}}`$.
So we haven’t come up with our embedded $`\text{NT}_3`$ yet. But why bother? The resistance of $`\text{NT}_{2.5849\mathrm{}}`$ out to infinity is
$$\frac{1}{3}+\frac{2}{9}+\frac{4}{27}+\mathrm{}=\frac{1}{3}\left(1+\frac{2}{3}+(\frac{2}{3})^2+\mathrm{}\right)=\frac{1}{3}\frac{1}{1\frac{2}{3}}=1.$$
Thus we have found an infinite subgraph of the 3-dimensional lattice having finite resistance out to infinity, and we are done.
###### Exercise 2.2.5
This exercise deals with the escape probability $`p_{\text{esc}}`$ for simple random walk in 3 dimensions. The idea is to turn upper and lower bounds for the resistance of the lattice into bounds for $`p_{\text{esc}}`$. Bounds are the best we can ask for using our method. The determination of the exact value will be discussed in Section 2.3.5. It is roughly .66.
(a) Use a shorting argument to get an upper bound for $`p_{\text{esc}}`$.
(b) We have seen that the resistance of the 3-dimensional lattice is at most one ohm. Show that the corresponding lower bound for $`p_{\text{esc}}`$ is $`1/6`$ Show that this lower bound can be increased to $`1/3`$ with no extra effort.
###### Exercise 2.2.6
Prove that simple random walk in any dimension $`d>3`$ is transient.
###### Exercise 2.2.7
Show how the not-quite embeddings of $`\text{NT}_2`$ and $`\text{NT}_{2.5849\mathrm{}}`$ can be altered to yield honest-to-goodness embeddings of “stretched-out” versions of these trees, obtained by replacing each edge of the tree by three edges in series. (Hint: “bounce”.)
#### 2.2.10 What we have done; what we will do
We have finally finished our electrical proof of Polya’s theorem. The proof of recurrence for $`d=1,2`$ was straight-forward, but this could hardly be said of the proof for $`d=3`$. After all, we started out trying to embed $`\text{NT}_3`$ and ended up by not quite embedding something that was not quite $`\text{NT}_3`$!
This is not bad in itself, for one frequently sets out to do something and in the process of trying to do it gets a better idea. The real problem is that this explicit construction is just too clever, too artificial. We seem to be saying that a simple random walk in 3 dimensions is transient because it happens to contain a beautifully symmetrical subgraph that is in some sense $`2.5849\mathrm{}`$-dimensional! Fine, but what if we hadn’t stumbled upon this subgraph? Isn’t there some other, more natural way?
We will see that indeed there are more natural approaches to showing transience for $`d=3`$. One such approach uses the same idea of embedding trees, but depends on the observation that one doesn’t need to be too careful about sending edges to edges. Another approach, based on relating the lattice not to a tree but to Euclidean space, was already hinted at in Section 2.1.8. The main goal for the rest of this work will be to explore these more natural electrical approaches to Polya’s theorem.
Before jumping into this, however, we are going to go back and take a look at a classical—i.e., probabilistic—approach to Polya’s theorem. This will give us something to compare our electrical proofs with.
### 2.3 The classical proofs of Polya’s Theorem
#### 2.3.1 Recurrence is equivalent to an infinite expected number of returns
For the time being, all of our random walks will be simple. Let $`u`$ be the probability that a random walker, starting at the origin, will return to the origin. The probability that the walker will be there exactly $`k`$ times (counting the initial time) is $`u^k(1u)`$. Thus, if $`m`$ is the expected number of times at the origin,
$$m=\underset{k=1}{}\mathrm{}ku^{k1}(1u)=\frac{1}{1u}.$$
If $`m=\mathrm{}`$ then $`u=1`$, and hence the walk is recurrent. If $`m<\mathrm{}`$ then $`u<1`$, so the walk is transient. Thus $`m`$ determines the type of the walk.
We shall use an alternate expression for $`m`$. Let $`u_n`$ be the probability that the walk, starting at $`\mathrm{𝟎}`$, is at $`\mathrm{𝟎}`$ on the nth step. Since the walker starts at $`\mathrm{𝟎}`$, $`u_0=1`$. Let $`e_n`$ be a random variable that takes on the value 1 if, at time $`n`$ the walker is at $`\mathrm{𝟎}`$ and 0 otherwise. Then
$$T=\underset{n=0}{\overset{\mathrm{}}{}}e_n$$
is the total number of times at $`\mathrm{𝟎}`$ and
$$m=𝐄(T)=\underset{n=0}{\overset{\mathrm{}}{}}𝐄(e_n).$$
But $`𝐄(e_n)=1u_n+0(1u_n)=u_n`$. Thus
$$m=\underset{n=0}{\overset{\mathrm{}}{}}u_n.$$
Therefore, the walk will be recurrent if the series $`_{n=0}^{\mathrm{}}`$ diverges and transient if it converges.
###### Exercise 2.3.1
Let $`N_{\mathrm{𝐱𝐲}}`$ be the expected number of visits to $`𝐲`$ for a random walker starting in $`𝐱`$. Show that $`N_{\mathrm{𝐱𝐲}}`$is finite if and only if the walk is transient.
#### 2.3.2 Simple random walk in one dimension
Consider a random walker in one dimension, started at $`\mathrm{𝟎}`$. To return to $`\mathrm{𝟎}`$, the walker must take the same number of steps to the right as to the left; hence, only even times are possible. Let us compute $`u_{2n}`$. Any path that returns in $`2n`$ steps has probability $`1/2^n`$. The number of possible paths equals the number of ways that we can choose the $`n`$ times to go right from the $`2n`$ possible times. Thus
$$u_{2n}=\left(\genfrac{}{}{0pt}{}{2n}{n}\right)\frac{1}{2^{2n}}.$$
We shall show that $`_nu_{2n}=\mathrm{}`$ by using Stirling’s approximation:
$$n!\sqrt{2\pi n}e^nn^n.$$
Thus
$$u_{2n}=\frac{(2n!)}{n!n!}\frac{1}{2^{2n}}\frac{\sqrt{2\pi 2n}e^{2n}(2n)^{2n}}{(\sqrt{2\pi n}e^nn^n)^22^{2n}}=\frac{1}{\sqrt{\pi n}}.$$
Therefore,
$$\underset{n}{}u_{2n}\underset{n}{}\frac{1}{\sqrt{\pi n}}=\mathrm{}$$
and a simple random walk in one dimension is recurrent.
Recall that this case was trivial by the resistor argument.
###### Exercise 2.3.2
We showed in Section 1.1.5 that a random walker starting at $`x`$ with $`0<x<N`$ has probability $`x/N`$ of reaching $`N`$ before 0. Use this to show that a simple random walk in one dimension is recurrent.
###### Exercise 2.3.3
Consider a random walk in one dimension that moves from $`n`$ to $`n+1`$ with probability $`p`$ and to $`n1`$ with probability $`q=1p`$. Assume that $`p>1/2`$. Let $`h_x`$ be the probability, starting at $`x`$, that the walker ever reaches 0. Use Exercise 1.1.9 to show that $`h_x=(q/p)^x`$ for $`x0`$ and $`h_x=1`$ for $`x<0`$. Show that this walk is transient.
###### Exercise 2.3.4
For a simple random walk in one dimension, it follows from Exercise 1.1.7 that the expected time, for a walker starting at $`x`$ with $`0<x<N`$, to reach 0 or $`n`$ is $`x(Nx)`$. Prove that for the infinite walk, the expected time to return to 0 is infinite.
###### Exercise 2.3.5
Let us regard a simple random walk in one dimension as the fortune of a player in a penny matching game where the players have unlimited credit. Show that the result is a martingale (see Section 1.1.6). Show that you can describe a stopping system that guarantees that you make money.
#### 2.3.3 Simple random walk in two dimensions
For a random walker in two dimensions to return to the origin, the walker must have gone the same number of times north and south and the same number of times east and west. Hence, again, only even times for return are possible. Every path that returns in $`2n`$ steps has probability $`1/4^{2n}`$. The number of paths that do this by taking $`k`$ steps to the north, $`k`$ south, $`nk`$ east and $`nk`$ west is
$$\left(\genfrac{}{}{0pt}{}{2n}{k,k,nk,nk}\right)=\frac{2n!}{k!k!(nk)!(nk)!}.$$
Thus
$`u_{2n}`$ $`=`$ $`{\displaystyle \frac{1}{4^{2n}}}{\displaystyle \underset{k=0}{\overset{n}{}}}{\displaystyle \frac{(2n)!}{k!k!(nk)!(nk)!}}`$
$`=`$ $`{\displaystyle \frac{1}{4^{2n}}}{\displaystyle \underset{k=0}{\overset{n}{}}}{\displaystyle \frac{(2n)!}{n!n!}}{\displaystyle \frac{n!n!}{k!k!(nk)!(nk)!}}`$
$`=`$ $`{\displaystyle \frac{1}{4^{2n}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2n}{n}}\right){\displaystyle \underset{k=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{k}}\right)^2.`$
But $`_{k=0}^n\left(\genfrac{}{}{0pt}{}{n}{k}\right)^2=\left(\genfrac{}{}{0pt}{}{2n}{n}\right)`$ (see Exercise 2.3.6). Hence
$$u_{2n}=\left(\frac{1}{2^{2n}}\left(\genfrac{}{}{0pt}{}{2n}{n}\right)\right)^2.$$
This is just the square of the one dimension result (not by accident, either—see Section 2.3.6). Thus we have
$$m=\underset{n}{}u_{2n}\underset{n}{}\frac{1}{\pi n}=\mathrm{}.$$
Recall that the resistor argument was very simple in this case also.
###### Exercise 2.3.6
Show that $`_{k=0}^n\left(\genfrac{}{}{0pt}{}{n}{k}\right)^2=\left(\genfrac{}{}{0pt}{}{2n}{n}\right)`$ (Hint: Think of choosing $`n`$ balls from a box that has $`2n`$ balls, $`n`$ black and $`n`$ white.)
#### 2.3.4 Simple random walk in three dimensions
For a walker in three dimensions to return, the walker must take an equal number of steps back and forth in each of the three different directions. Thus we have
$$u_{2n}=\frac{1}{6^{2n}}\underset{j,k}{}\frac{(2n)!}{j!j!k!k!(njk)!(njk)!}$$
where the sum is taken over all $`j,k`$ with $`j+kn`$. Following Feller , we rewrite this sum as
$$u_{2n}=\frac{1}{2^{2n}}\left(\genfrac{}{}{0pt}{}{2n}{n}\right)\underset{j,k}{}\left(\frac{1}{3^n}\frac{n!}{j!k!(njk)!}\right)^2.$$
Now consider placing $`n`$ balls randomly into three boxes, $`A,B,C`$. The probability that $`j`$ balls fall in $`A`$, $`k`$ in $`B`$, and $`njk`$ in $`C`$ is
$$\frac{1}{3^n}\left(\genfrac{}{}{0pt}{}{n}{j,k,njk}\right)=\frac{1}{3^n}\frac{n!}{j!k!(njk)!}.$$
Intuitively, this probability is largest when $`j`$, $`k`$, and $`njk`$ are as near as possible to $`n/3`$, and this can be proved (see Exercise 2.3.7). Hence, replacing one of the factors in the square by this larger value, we have:
$$u_{2n}\frac{1}{2^{2n}}\left(\genfrac{}{}{0pt}{}{2n}{n}\right)\left(\frac{1}{3^n}\frac{n!}{\frac{n}{3}!\frac{n}{3}!\frac{n}{3}!}\right)\left(\underset{j,k}{}\frac{1}{3^n}\frac{n!}{j!k!(njk)!}\right),$$
where $`n/3`$ denotes the greatest integer $`n/3`$. The last sum is 1 since it is the sum of all probabilities for the outcomes of putting $`n`$ balls into three boxes. Thus
$$u_{2n}\frac{1}{2^{2n}}\left(\genfrac{}{}{0pt}{}{2n}{n}\right)\left(\frac{1}{3^n}\frac{n!}{\frac{n}{3}!^3}\right).$$
Applying Stirling’s approximation yields
$$u_{2n}\frac{K}{n^{3/2}}$$
for suitable constant $`K`$. Therefore
$$m=\underset{n}{}u_{2n}K\underset{n}{}\frac{1}{n^{3/2}}<\mathrm{},$$
and a simple random walk in three dimensions is recurrent.
While this is a complex calculation, the resistor argument was also complicated in this case. We will try to make amends for this presently.
###### Exercise 2.3.7
Prove that $`\left(\genfrac{}{}{0pt}{}{n}{j,k,njk}\right)`$ is largest when $`j`$, $`k`$, and $`njk`$ are as close as possible to $`n/3`$.
###### Exercise 2.3.8
Find an appropriate value for the ”suitable constant” $`K`$ that was mentioned above, and derive an upper bound for $`m`$. Use this to get a lower bound for the probability of escape for simple random walk in three dimensions.
#### 2.3.5 The probability of return in three dimensions: exact calculations
Since the probability of return in three dimensions is less than one, it is natural to ask, “What is this probability?” For this we need an exact calculation. The first such calculation was carried out in a pioneering paper of McCrea and Whipple . The solution outlined here follows Feller , Exercise 28, Chapter 9, and Montroll and West .
Let $`p(a,b,c;n)`$ be the probability that a random walker, starting at $`\mathrm{𝟎}`$, is at $`(a,b,c)`$ after $`n`$ steps. Then $`p(a,b,c;n)`$ is completely determined by the fact that
$$p(0,0,0;0)=1$$
and
$`p(a,b,c;n)`$ $`=`$ $`{\displaystyle \frac{1}{6}}p(a1,b,c;n1)+{\displaystyle \frac{1}{6}}p(a+1,b,c;n1)+`$
$`{\displaystyle \frac{1}{6}}p(a,b1,c;n1)+{\displaystyle \frac{1}{6}}p(a,b+1,c;n1)+`$
$`{\displaystyle \frac{1}{6}}p(a,b,c1;n1)+{\displaystyle \frac{1}{6}}p(a,b,c+1;n1).`$
Using the technique of generating functions, it is possible to derive a solution of these equations as
$`p(a,b,c;n)`$
$`=`$ $`{\displaystyle \frac{1}{(2\pi )^3}}`$
$`{\displaystyle _\pi ^\pi }{\displaystyle _\pi ^\pi }{\displaystyle _\pi ^\pi }\left({\displaystyle \frac{\mathrm{cos}x+\mathrm{cos}y+\mathrm{cos}z}{3}}\right)^n\mathrm{cos}(xa)\mathrm{cos}(yb)\mathrm{cos}(zc)𝑑x𝑑y𝑑z.`$
Of course, we can just verify that this formula satisfies our equations once someone has suggested it. Having this formula, we put $`a=b=c=0`$ and sum over $`n`$ to obtain the expected number of returns as
$$m=\frac{3}{(2\pi )^3}_\pi ^\pi _\pi ^\pi _\pi ^\pi \frac{1}{3(\mathrm{cos}x+\mathrm{cos}y+\mathrm{cos}z)}𝑑x𝑑y𝑑z.$$
This integral was first evaluated by Watson in terms of elliptic integrals, which are tabulated. A simpler result was obtained by Glasser and Zucker who evaluated this integral in terms of gamma functions. Using this result, we get
$$m=\frac{\sqrt{6}}{32\pi ^3}\mathrm{\Gamma }\left(\frac{1}{24}\right)\mathrm{\Gamma }\left(\frac{5}{24}\right)\mathrm{\Gamma }\left(\frac{7}{24}\right)\mathrm{\Gamma }\left(\frac{11}{24}\right)=1.516386059137\mathrm{},$$
where
$$\mathrm{\Gamma }(x)=_0^{\mathrm{}}e^tt^{x1}𝑑t$$
is Euler’s gamma function. (Incidentally, the value given by Glasser and Zucker for the integral above needs to be corrected by multiplying by a factor of $`1/(384\pi )`$.)
Recall that $`m=1/(1u)`$ so that $`u=(m1)/m`$. This gives
$$u=.340537329544\mathrm{}.$$
#### 2.3.6 Simple random walk in two dimensions is the same as two independent one-dimensional random walks
We observed that the probability of return at time $`2n`$ in two dimensions is the square of the corresponding probability in one dimension. Thus it is the same as the probability that two independent walkers, one walking in the $`x`$ direction and the other in the $`y`$ direction, will, at time $`2n`$, both be at $`0`$. Can we see that this should be the case? The answer is yes. Just change our axes by 45 degrees to new axes $`\overline{x}`$ and $`\overline{y}`$ as in Figure 60.
Look at the possible outcomes for the first step using the $`x,y`$ coordinates and the $`\overline{x},\overline{y}`$ coordinates. We have
$$\begin{array}{cc}x,y\text{ coordinates}\hfill & \overline{x},\overline{y}\text{ coordinates}\hfill \\ (0,1)\hfill & (1/\sqrt{2},1/\sqrt{2})\hfill \\ (0,1)\hfill & (1/\sqrt{2},1/\sqrt{2})\hfill \\ (1,0)\hfill & (1/\sqrt{2},1/\sqrt{2})\hfill \\ (1,0)\hfill & (1/\sqrt{2},1/\sqrt{2})\hfill \end{array}.$$
Assume that we have two independent walkers, one moving with step size $`\frac{1}{\sqrt{2}}`$ randomly along the $`\overline{x}`$ axis and the other moving with the same step size along the $`\overline{y}`$ axis. Then, if we plot their positions using the $`x,y`$ axes, the four possible outcomes for the first step would agree with those given in the second column of the table above. The probabilities for each of the four pairs of outcomes would also be $`(1/2)(1/2)=1/4`$. Thus, we cannot distinguish a simple random walk in two dimensions from two independent walkers along the $`\overline{x}`$ and $`\overline{y}`$ axes making steps of magnitude $`1/\sqrt{2}`$. Since the probability of return does not depend upon the magnitude of the steps, the probability that our two independent walkers are at $`(0,0)`$ at time $`2n`$ is equal to the product of the probabilities that each separately is at 0 at time $`2n`$, namely $`(1/2^{2n})\left(\genfrac{}{}{0pt}{}{2n}{n}\right)`$. Therefore, the probability that the standard walk will be at $`(0,0)`$ at time $`2n`$ is $`((1/2^{2n})\left(\genfrac{}{}{0pt}{}{2n}{n}\right))^2`$ as observed earlier.
#### 2.3.7 Simple random walk in three dimensions is not the same as three independent random walks
In three dimensions, the probability that three independent walkers are each back to 0 after time $`2n`$ is
$$u_{2n}=\left(\left(\genfrac{}{}{0pt}{}{2n}{n}\right)\frac{1}{2^{2n}}\right)^3.$$
This does not agree with our result for a simple random walk in three dimensions. Hence, the same trick cannot work. However, it is interesting to consider a random walk which is the result of three independent walkers. Let $`(i,j,k)`$ be the position of three independent random walkers. The next position is one of the eight possibilities $`(i\pm 1,j\pm 1,k\pm 1)`$ Thus we may regard their progress as a random walk on the lattice points in three dimensions. If we center a cube of side 2 at $`(i,j,k)`$, then the walk moves with equal probabilities to each of the eight corners of the cube. It is easier to show that this random walk is transient (using classical methods) than it is for simple random walk. This is because we can again use the one-dimension calculation. The probability $`u_{2n}`$ for return at time $`2n`$ is
$$u_{2n}=\left(\left(\genfrac{}{}{0pt}{}{2n}{n}\right)\frac{1}{2^{2n}}\right)^3\left(\frac{1}{\sqrt{\pi n}}\right)^3.$$
Thus
$$m=\underset{n}{}u_{2n}\underset{n}{}\left(\frac{1}{\pi n}\right)^{3/2}<\mathrm{}$$
and the walk is transient.
The fact that this three independent walkers model and simple random walk are of the same type suggests that when two random walks are “really about the same”, they should either both be transient or both be recurrent. As we will soon see, this is indeed the case. Thus we may infer the transience of simple random walk in 3 dimensions from the transience of the three independent walkers model without going through the involved calculation of Section 2.3.4.
### 2.4 Random walks on more general infinite networks
#### 2.4.1 Random walks on infinite networks
From now on we assume that $`G`$ is an infinite connected graph. We assume that it is of *bounded degree*, by which we mean that there is some integer $`E`$ such that the number of edges from any point is at most $`E`$. We assign to each edge $`xy`$ of $`G`$ a conductance $`C_{xy}>0`$ with $`R_{xy}=\frac{1}{C_{xy}}`$. The graph $`G`$ together with the conductances $`𝐂=(C_{xy})`$ is called a *network* and denoted by $`(G,𝐂)`$. Given a network $`(G,𝐂)`$, we define a random walk by
$$P_{xy}=\frac{C_{xy}}{C_x}$$
where $`C_x=_yC_{xy}`$. When all the conductances are equal, we obtain a random walk that moves along each edge with the same probability: In agreement with our previous terminology, we call this walk *simple random walk* on $`G`$.
We have now a quite general class of infinite-state Markov chains. As in the case of finite networks, the chains are reversible Markov chains: That is, there is a positive vector $`𝐰`$ such that $`w_xP_{xy}=w_yP_{yx}`$. As in the finite case, we can take $`w_x=C_x`$, since $`C_xP_{xy}=C_{xy}=C_{yx}=C_yP_{yx}`$.
#### 2.4.2 The type problem
Let $`(G,𝐂)`$ be an infinite network with random walk $`𝐏`$. Let $`\mathrm{𝟎}`$ be a reference point. Let $`p_{\text{esc}}`$ be the probability that a walk starting at $`\mathrm{𝟎}`$ will never return to $`\mathrm{𝟎}`$. If $`p_{\text{esc}}=0`$ we say that $`𝐏`$ is *recurrent*, and if $`p_{\text{esc}}>0`$ we say that it is *transient*. You are asked to show in Exercise 2.4.1 that the question of recurrence or transience of $`𝐏`$ does not depend upon the choice of the reference point. The *type problem* is the problem of determining if a random walk (network) is recurrent or transient.
In Section 2.1.5 we showed how to rephrase the type problem for a lattice in terms of finite graphs sitting inside it. In Section 2.1.6 we showed that the type problem is equivalent to an electrical network problem by showing that simple random walk on a lattice is recurrent if and only if the lattice has infinite resistance to infinity. The same arguments apply with only minor modifications to the more general infinite networks as well. This means that we can use Rayleigh’s short-cut method to determine the type of these more general networks.
###### Exercise 2.4.1
Show that the question of recurrence or transience of $`𝐏`$ does not depend upon the choice of the reference point.
#### 2.4.3 Comparing two networks
Given two sets of conductances $`𝐂`$ and $`\overline{𝐂}`$ on $`G`$, we say that $`(G,\overline{𝐂})<(G,𝐂)`$ if $`\overline{C}_{xy}<C_{xy}`$ for all $`xy`$, or equivalently, if $`\overline{R}_{xy}>R_{xy}`$ for all $`xy`$. Assume that $`(G,\overline{𝐂})<(G,𝐂)`$. Then by the Monotonicity Law, $`\overline{R}_{\text{eff}}R_{\text{eff}}`$. Thus if random walk on $`(G,\overline{𝐂})`$ is transient, i.e., if $`\overline{R}_{\text{eff}}<\mathrm{}`$, then random walk on $`(G,𝐂)`$ is also transient. If random walk on $`(G,𝐂)`$ is recurrent, i.e., if $`R_{\text{eff}}=\mathrm{}`$, then random walk on $`(G,\overline{𝐂})`$ is also recurrent.
Theorem. If $`(G,𝐂)`$ and $`(G,\overline{𝐂})`$ are networks, and if there exist constants $`u,v`$ with $`0<uv<\mathrm{}`$ such that
$$uC_{xy}\overline{C}_{xy}vC_{xy}$$
for all $`x`$ and $`y`$, then random walk on $`(G,\overline{𝐂})`$ is of the same type as random walk on $`(G,𝐂)`$.
Proof. Let $`U_{xy}=uC_{xy}`$ and $`V_{xy}=vC_{xy}`$. Then $`(G,𝐔)(G,\overline{𝐂})(G,𝐕)`$. But the random walks for $`(G,𝐔)`$ and $`(G,𝐕)`$ are the same as random walk on $`(G,𝐂)`$. Thus random walk for $`(G,\overline{𝐂})`$ is of the same type as random walk on $`(G,𝐂)`$. $`\mathrm{}`$
Corollary. Let $`(G,𝐂)`$ be a network. If for every edge $`xy`$ of $`G`$ we have $`0<u<C_{xy}<v<\mathrm{}`$ for some constants $`u`$ and $`v`$, then the random walk on $`(G,𝐂)`$ has the same type as simple random walk on $`G`$.
###### Exercise 2.4.2
Consider the two-dimensional lattice. For each edge, we toss a coin to decide what kind of resistor to put across this edge. If heads turns up, we put a two-ohm resistor across this edge; if tails turns up, we put a one-ohm resistor across the edge. Show that the random walk on the resulting network is recurrent.
###### Exercise 2.4.3
Consider the analogous problem to Exercise 2.4.2 in 3 dimensions.
#### 2.4.4 The $`k`$-fuzz of a graph
For any integer $`k`$, the $`k`$-*fuzz* of a graph $`G`$ is the graph $`G_k`$ obtained from $`G`$ by adding an edge $`xy`$ if it is possible to go from $`x`$ to $`y`$ in at most $`k`$ steps. For example, the 2-fuzz of the two-dimensional lattice is shown in Figure 61; please note that horizontal and vertical edges of length 2, such as those joining $`(0,0)`$ to $`(0,2)`$, have not been indicated.
Theorem. Simple random walk on $`G`$ and on the $`k`$-fuzz $`G_k`$ of $`G`$ have the same type.
Proof. Let $`𝐏`$ be simple random walk on $`G`$. Define $`\overline{𝐏}=(𝐏+𝐏^2+\mathrm{}+𝐏^k)/k`$. Then $`\overline{𝐏}`$ may be considered to be $`𝐏`$, watched at one of the first $`k`$ steps chosen at random, then at a time chosen at random from the next $`k`$ steps after this time, etc. Thinking of $`\overline{𝐏}`$ in this way, we see that $`𝐏`$ is in state $`\mathrm{𝟎}`$ at least once for every time $`\overline{𝐏}`$ in state $`\mathrm{𝟎}`$. Hence, if $`\overline{𝐏}`$ is recurrent so is $`𝐏`$. Assume now that $`\overline{𝐏}`$ is transient. Choose a finite set $`S`$ so that $`\mathrm{𝟎}`$ cannot be reached in $`k`$ steps from a point outside of $`S`$. Then, since the walk $`\overline{𝐏}`$ will be outside $`S`$ from some time on, the walk $`𝐏`$ cannot be at $`\mathrm{𝟎}`$ after this time, and $`𝐏`$ is also transient. Therefore, $`𝐏`$ and $`\overline{𝐏}`$ are of the same type.
Finally, we show that $`\overline{𝐏}`$ has the same type as simple random walk on $`G_k`$. Here it is important to remember our restriction that $`G`$ is of bounded degree, so that for some $`E`$ no vertex has degree $`>E`$. We know that $`𝐏`$ is reversible with $`\mathrm{𝐰𝐏}=𝐰`$, where $`w_x`$ is the number of edges coming out of $`x`$. From its construction, $`\overline{𝐏}`$ is also reversible and $`𝐰\overline{𝐏}=𝐰`$. $`\overline{𝐏}`$ is the random walk on a network $`(G_k,\overline{𝐂})`$ with $`\overline{C}_{xy}=w_x\overline{P}_{xy}`$. If $`\overline{P}_{xy}>0`$, there is a path $`x,x_1,x_2,\mathrm{},x_{m1},y`$ in $`G`$ from $`x`$ to $`y`$ of length $`mk`$. Then
$$\overline{P}_{xy}\frac{1}{k}(\frac{1}{E})^m\frac{1}{k}(\frac{1}{E})^k.$$
Thus
$$0<\frac{1}{k}(\frac{1}{E})^k\overline{P}_{xy}1$$
and
$$0<\frac{1}{k}(\frac{1}{E})^k\overline{C}_{xy}E.$$
Therefore, by the theorem on the irrelevance of bounded twiddling proven in Section 2.4.3, $`\overline{𝐏}`$ and simple random walk on $`G_k`$ are of the same type. So $`G`$ and $`G_k`$ are of the same type.
NOTE: This is the only place where we use probabilistic methods of proof. For the purist who wishes to avoid probabilistic methods, Exercise 2.4.9 indicates an alternative electrical proof.
We show how this theorem can be used. We say that a graph $`G`$ can be *embedded* in a graph $`\overline{G}`$ if the points $`x`$ of $`G`$ can be made to correspond in a one-to-one fashion to points $`\overline{x}`$ of $`\overline{G}`$ in such a way that if $`xy`$ is an edge in $`G`$, then $`\overline{x}\overline{y}`$ is an edge in $`\overline{G}`$.
Theorem. If simple random walk on $`G`$ is transient, and if $`\overline{G}`$ can be embedded in a $`k`$-fuzz $`\overline{G}_k`$ of $`\overline{G}`$ then simple random walk on $`\overline{G}`$ is also transient. Simple random walk on $`G`$ and $`\overline{G}`$ are of the same type if each graph can be embedded in a $`k`$-fuzz of the other graph.
Proof. Assume that simple random walk on $`G`$ is transient and that $`G`$ can be embedded in a $`k`$-fuzz $`\overline{G}_k`$ of $`\overline{G}`$. Since $`R_{\text{eff}}`$ for $`G`$ is finite and $`G`$ can be embedded in $`\overline{G}_k`$, $`R_{\text{eff}}`$ for $`\overline{G}_k`$ is finite. By our previous theorem, the same is true for $`\overline{G}`$ and simple random walk on $`\overline{G}`$ is transient.
If we can embed $`G`$ in $`\overline{G}_k`$ and $`\overline{G}`$ in $`G_k`$, then the random walk on $`G`$ is transient if and only if the random walk on $`\overline{G}`$ is. $`\mathrm{}`$
###### Exercise 2.4.4
We have assumed that there is a bound $`E`$ for the number of edges coming out of any point. Show that if we do not assume this, it is not necessarily true that $`G`$ and $`G_k`$ are of the same type. (Hint: Consider a network something like that shown in Figure 62.)
#### 2.4.5 Comparing general graphs with lattice graphs
We know the type of simple random walk on a lattice $`𝐙^d`$. Thus to determine the type of simple random walk on an arbitrary graph $`G`$, it is natural to try to compare G with $`𝐙^d`$. This is feasible for graphs that can be drawn in some Euclidean space $`𝐑^d`$ in a civilized manner.
Definition. A graph $`G`$ can be drawn in a Euclidean space $`𝐑^d`$ in a *civilized manner* if its vertices can be embedded in $`𝐑^d`$ so that for some $`r<\mathrm{}`$, $`s>0`$
(a) The length of each edge is $`r`$.
(b) The distance between any two points is $`>s`$.
Note that we make no requirement about being able to draw the edges of $`G`$ so they don’t intersect.
Theorem. If a graph can be drawn in $`𝐑^d`$ in a civilized manner, then it can be embedded in a $`k`$-fuzz of the lattice $`𝐙^d`$.
Proof. We carry out the proof for the case $`d=2`$. Assume that $`G`$ can be drawn in a civilized manner in $`𝐑^2`$. We want to show that $`G`$ can be embedded in a $`k`$-fuzz of $`𝐙^2`$. We have been thinking of $`𝐙^2`$ as being drawn in $`𝐑^2`$ with perpendicular lines and adjacent points a unit distance apart on these lines, but this embedding is only one particular way of representing $`𝐙^2`$. To emphasize this, let’s talk about $`L^2`$ instead of $`𝐙^2`$. Figure 63 shows another way of drawing $`L^2`$ in $`𝐑^2`$.
From a graph-theoretical point of view, this is the same as $`𝐙^2`$. In trying to compare $`G`$ to $`L`$, we take advantage of this flexibility by drawing $`L^2`$ so small that points of $`G`$ can be moved onto points of $`L^2`$ without bumping into each other.
Specifically, let $`L^2`$ be a two-dimensional rectangular lattice with lines a distance $`s/2`$ apart. In any square of $`L^2`$, there is at most one point of $`G`$. Move each point $`x`$ of $`G`$ to the southwest corner $`\overline{x}`$ of the square that it is in, as illustrated in Figure 64.
Now since any two adjacent points $`x`$, $`y`$ in $`G`$ were within $`r`$ of each other in $`𝐑^2`$, the corresponding points $`\overline{x}`$, $`\overline{y}`$ in $`L^2`$ will have Euclidean distance $`<r+2s`$. Choose $`k`$ so that any two points of $`L^2`$ whose Euclidean distance is $`<r+2s`$ can be connected by a path in $`L^2`$ of at most $`k`$ steps. Then $`\overline{x}`$ and $`\overline{y}`$ will be adjacent in $`L_k^2`$ and—since the prescription for $`k`$ does not depend on $`x`$ and $`y`$—we have embedded $`G`$ in the $`k`$-fuzz $`L_k^2`$.
Corollary. If $`G`$ can be drawn in a civilized manner in $`𝐑^1`$ or $`𝐑^2`$, then simple random walk on $`G`$ is recurrent.
Proof. Assume, for example, that $`G`$ can be drawn in a civilized manner in $`𝐑^2`$. Then $`G`$ can be embedded in a $`k`$-fuzz $`𝐙_k^2`$ of $`𝐙^2`$. If simple random walk on $`G`$ were transient, then the same would be true for $`𝐙_k^2`$ and $`𝐙^2`$. But we know that simple random walk on $`𝐙^2`$ is recurrent. Thus simple random walk on $`G`$ is recurrent. $`\mathrm{}`$
Our first proof that random walk in three dimensions is transient consisted in showing that we could embed a transient tree in $`𝐙^3`$. We now know that it would have been sufficient to show how to draw a transient tree in $`𝐑^3`$ in a civilized manner: This is easier (see Exercise 2.4.5).
The corollary implies that simple random walk on any sufficiently symmetrical graph in $`𝐑^2`$ is recurrent. For example, simple random walk on the regular graph made up of hexagons shown in Figure 65 is recurrent.
We can even consider very irregular graphs. For example, on the cover of the January 1977 Scientific American, there is an example due to Conway of an infinite non-periodic tiling using Penrose tiles of the form shown in Figure 66.
It is called the cartwheel pattern; part of it is shown in Figure 67.
A walker walking randomly on the edges of this very irregular infinite tiling will still return to his or her starting point.
Assume now that $`G`$ can be drawn in a civilized manner in $`𝐑^3`$. Then to show that simple random walk on $`G`$ is of the same type as $`𝐙^3`$, namely transient, it is sufficient to show that we can embed $`𝐙^3`$ in a $`k`$-fuzz of $`G`$. This is clearly possible for any regular lattice in $`𝐑^3`$. The three lattices that have been most studied and for which exact probabilities for return have been found are called the SC, BCC, and FCC lattices. The SC (simple cubic) lattice is just $`𝐙^3`$. The walker moves each time to a new point by adding a random choice from the six vectors
$$(\pm 1,0,0),(0,\pm 1,0),(0,0,\pm 1).$$
For the BCC (body-centered cubic) lattice, the choice is one of the eight vectors
$$(\pm 1,\pm 1,\pm 1).$$
This was the walk that resulted from three independent one-dimensional walkers. For the FCC (face-centered cubic) lattice, the random choice is made from the twelve vectors
$$(\pm 1,\pm 1,0),(\pm 1,0,\pm 1),(0,\pm 1,\pm 1).$$
For a discussion of exact calculations for these three lattices, see Montroll and West
As we have seen, once the transience of any one of these three walks is established, no calculations are necessary to determine that the other walks are transient also. Thus we have yet another way of establishing Polya’s theorem in three dimensions: Simply verify transience of the walk on the BCC lattice via the simple three-independent-walkers computation, and infer that walk on the SC lattice is also transient since the BCC lattice can be embedded in a $`k`$-fuzz of it.
###### Exercise 2.4.5
When we first set out to prove Polya’s theorem for $`d=3`$, our idea was to embed $`\text{NT}_3`$ in $`𝐙^3`$. As it turned out, what we ended up embedding was not $`\text{NT}_3`$ but $`\text{NT}_{2.5849\mathrm{}}`$, and we didn’t quite embed it at that. We tried to improve the situation by finding (in Exercise 2.2.7) an honest-to-goodness embedding of a relative of $`\text{NT}_{2.5849\mathrm{}}`$, but $`\text{NT}_3`$ was still left completely out in the cold. Now, however, we are in a position to embed $`\text{NT}_3`$, if not in $`𝐙^3`$ then at least in a $`k`$-fuzz of it. All we need to do is to draw $`\text{NT}_3`$ in $`𝐑^3`$ in a civilized manner. Describe how to do this, and thereby give one more proof of Polya’s theorem for $`d=3`$.
###### Exercise 2.4.6
Find a graph that can be embedded in a civilized manner in $`𝐑^3`$ but not in $`𝐑^2`$, but is nonetheless recurrent.
###### Exercise 2.4.7
Assume that $`G`$ is drawn in a civilized manner in $`𝐑^3`$. To show that simple random walk on $`G`$ is transient, it is enough to know that $`𝐙^3`$ can be embedded in a $`k`$-fuzz of G. Try to come up with a nice condition that will guarantee that this is possible. Can you make this condition simple, yet general enough so that it will settle all reasonably interesting cases? In other words, can you make the condition nice enough to allow us to remember only the condition, and forget about the general method lying behind it?
#### 2.4.6 Solving the type problem by flows: a variant of the cutting method
In this section we will introduce a variant of the cutting method whereby we use Thomson’s Principle directly to estimate the effective resistance of a conductor.
Thomson’s Principle says that, given any unit flow through a resistive medium, the dissipation rate of that flow gives an upper bound for the effective resistance of the medium. This suggests that to show that a given infinite network is transient, it should be enough to produce a unit flow out to infinity having finite energy dissipation.
In analogy with the finite case, we say that $`𝐣`$ is a *flow from $`\mathrm{𝟎}`$ to infinity* if
(a) $`j_{xy}=j_{yx}`$.
(b) $`_yj_{xy}=0`$ if $`x\mathrm{𝟎}`$.
We define $`j_\mathrm{𝟎}=_yj_{\mathrm{𝟎}y}`$. If $`j_\mathrm{𝟎}=1`$, we say that $`𝐣`$ is a *unit flow to infinity*. Again in analogy with the finite case, we call $`\frac{1}{2}_{x,y}j_{xy}^2R_{xy}`$ the *energy dissipation* of the flow $`𝐣`$.
Theorem. The effective resistance $`R_{\text{eff}}`$ from $`\mathrm{𝟎}`$ to $`\mathrm{}`$ is less than or equal to the energy dissipation of any unit flow from $`\mathrm{𝟎}`$ to infinity.
Proof. Assume that we have a unit flow $`𝐣`$ from $`\mathrm{𝟎}`$ to infinity with energy dissipation
$$E=\frac{1}{2}\underset{x,y}{}j_{xy}^2R_{xy}.$$
We claim that $`R_{\text{eff}}E`$. Restricting $`j_{xy}`$ to the edges of the finite graph $`G^{(r)}`$, we have a unit flow from $`\mathrm{𝟎}`$ to $`G^{(r)}`$ in $`G^{(r)}`$. Let $`i^{(r)}`$ be the unit current flow in $`G^{(r)}`$ from 0 to $`G^{(r)}`$. By the results of Section 1.3.5,
$$𝐑_{\text{eff}}^{(r)}=\frac{1}{2}\underset{G^{(r)}}{}(i_{xy}^{(r)})^2R_{xy}\frac{1}{2}\underset{G^{(r)}}{}j_{xy}^2R_{xy}\frac{1}{2}\underset{x,y}{}j_{xy}^2R_{xy}=E,$$
where $`_{G^{(r)}}`$ indicates the sum over all pairs $`x,y`$ such that $`xy`$ is an edge of $`G^{(r)}`$.
###### Exercise 2.4.8
We have billed the method of using Thomson’s Principle directly to estimate the effective resistances of a network as a variant of the cutting method. Since the cutting method was derived from Thomson’s Principle, and not vice versa, it would seem that we have got the cart before the horse. Set this straight by giving an informal (“heuristic”) derivation of Thomson’s Principle from the cutting method. (Hint: see Maxwell , Chapter VIII, Paragraph 307.) For more on this question, see Onsager .
###### Exercise 2.4.9
Let $`G`$ be an infinite graph of bounded degree and $`G_k`$ the $`k`$-fuzz of $`G`$. Using electric network arguments, show that $`R_{\text{eff}}<\mathrm{}`$ for $`G`$ if and only if $`R_{\text{eff}}<\mathrm{}`$ for $`G_k`$.
#### 2.4.7 A proof, using flows, that simple random walk in three dimensions is transient
We now apply this form of the cutting method to give another proof that simple random walk on the threedimensional lattice is transient. All we need is a flow to infinity with finite dissipation. The flow we are going to describe is not the first flow one would think of. In case you are curious, the flow described here was constructed as a side effect of an unsuccessful attempt to derive the isoperimetric inequality (see Polya ) from the “max-flow min-cut” theorem (Ford and Fulkerson ). The idea is to find a flow in the positive orthant having the property that the same amount flows through all points at the same distance from $`\mathrm{𝟎}`$.
Again, it is easiest to show the construction for the two-dimensional case. Let $`G`$ denote the part of $`𝐙^2`$ lying in the first quadrant. The graph $`G^{(4)}`$ is shown in Figure 68.
We choose our flow so that it always goes away from $`\mathrm{𝟎}`$. Into each point that is not on either axis there are two flows, one vertical and one horizontal. We want the sum of the corresponding values of $`j_{xy}`$ to be the same for all points the same distance from $`\mathrm{𝟎}`$. These conditions completely determine the flow. The flow out of the point $`(x,y)`$ with $`x+y=n`$ is as shown in Figure 69.
The values for the currents out to the fourth level are shown in Figure 70.
In general, the flow out of a point $`(x,y)`$ with $`x+y=n`$ is
$$\frac{x+1}{(n+2)(n+1)}+\frac{y+1}{(n+2)(n+1)}=\frac{1}{n+1}$$
and the flow into this point is
$$\frac{x}{n(n+1)}+\frac{y}{n(n+1)}=\frac{1}{n+1}$$
Thus the net flow at $`(x,y)`$ is 0. The flow out of $`\mathrm{𝟎}`$ is $`(1/2)+(1/2)=1`$. For this two-dimensional flow, the energy dissipation is infinite, as it would have to be. For three dimensions, the uniform flow is defined as follows: Out of $`(x,y,z)`$ with $`x+y+z=n`$ we have the flow indicated in Figure 71.
The total flow out of $`(x,y,z)`$ is then
$`{\displaystyle \frac{2(x+1)}{(n+3)(n+2)(n+1)}}+{\displaystyle \frac{2(y+1)}{(n+3)(n+2)(n+1)}}+{\displaystyle \frac{2(z+1)}{(n+3)(n+2)(n+1)}}`$
$`=`$ $`{\displaystyle \frac{2}{(n+2)(n+1)}}.`$
The flow into $`(x,y,z)`$ comes from the points $`(x1,y,z)`$, $`(x,y1,z)`$, $`(x,y,z1)`$ and, hence, the total flow into $`(x,y,z)`$ is
$$\frac{2x}{(n+2)(n+1)n}+\frac{2y}{(n+2)(n+1)n}+\frac{2z}{(n+2)(n+1)n}=\frac{2}{(n+2)(n+1)}.$$
Thus the net flow for $`(x,y,z)`$ is 0. The flow out of $`\mathrm{𝟎}`$ is $`(1/3)+(1/3)+(1/3)=1`$ We have now to check finiteness of energy dissipation. The flows coming out of the edges at the $`n`$th level are all $`2/(n+1)^2`$. There are $`(n+1)(n+2)/2`$ points a distance $`n`$ from $`\mathrm{𝟎}`$, and thus there are $`(3/2)(n+1)(n+2)3(n+1)^2`$ edges coming out of the $`n`$th level. Thus the energy dissipation $`E`$ has
$$E\underset{n}{}3(n+1)^2\left(\frac{2}{(n+1)^2}\right)^2=12\underset{n}{}\frac{1}{(n+1)^2}<\mathrm{},$$
and the random walk is transient.
#### 2.4.8 The end
We have come to the end of our labors, and it seems fitting to look back and try to say what it is we have learned.
To begin with, we have seen how phrasing certain mathematical questions in physical terms allows us to draw on a large body of physical lore, in the form of established methods and ways of thought, and thereby often leads us to the answers to those questions.
In particular, we have seen the utility of considerations involving energy. In took hundreds of years for the concept of energy to emerge and take its rightful place in physical theory, but it is now recognized as perhaps the most fundamental concept in all of physics. By phrasing our probabilistic problems in physical terms, we were naturally led to considerations of energy, and these considerations showed us the way through the difficulties of our problems.
As for Polya’s theorem and the type problem in general, we have picked up a bag of tricks, known collectively as “Rayleigh’s short-cut method”, which we may expect will allow us to determine the type of almost any random walk we are likely to embark on. In the process, we have gotten some feeling for the connection between the dimensionality of a random walk and its type. Furthermore, we have settled one of the main questions likely to occur to someone encountering Polya’s theorem, namely: “If two walks look essentially the same, and if one has been shown to be transient, must not the other also be transient?”
Another question likely to occur to someone contemplating Polya’s theorem is the question raised in Section 2.1.8: “Since the lattice $`𝐙^d`$ is in some sense a discrete analog of a resistive medium filling all of $`𝐑^d`$, should it not be possible to go quickly and naturally from the trivial computation of the resistance to infinity of the continuous medium to a proof of Polya’s theorem?” Our shorting argument allowed us to do this in the two-dimensional case; that leaves the case of three (or more) dimensions. Again, it is considerations of energy that allow us to make this connection. The trick is to start with the flow field that one gets by solving the continuous problem, and adapt it to the lattice, so as to get a lattice flow to infinity having finite dissipation rate. We leave the working out of this as an exercise, so as not to rob readers of the fun of doing it for themselves.
###### Exercise 2.4.10
Give one final proof of Polya’s theorem in 3 dimensions by showing how to adapt the $`1/r^2`$ radial flow field to the lattice. (Hint: “cubes”.)
## Acknowledgements
This work is derived from the book *Random Walks and Electric Networks*, originally published in 1984 by the Mathematical Association of American in their Carus Monographs series. We are grateful to the MAA for permitting this work to be freely redistributed under the terms of the GNU General Public License. (See Figure 72.) |
warning/0001/math0001046.html | ar5iv | text | # Quantum co-adjoint orbits of the real diamond group
## 1. Introduction
Let us first recall that it was Hermann Weyl(see \[W\]), who introduced a mapping from classical observables (i.e. functions on the phase space $`^{2n}`$) to quantum observables (i.e. normal operators in the Hilbert space $`L^2(^n)`$). The idea was to express functions on $`^{2n}`$ as Fourier transforms and then, by using the inverse Fourier transforms to correspond this correspondence to functions on characters, i.e. one-dimensional representations of the Heisenberg group, parameterized by the Planck constant $`\mathrm{}`$ \- and finally to present them as elements in the corresponding infinite-dimensional representations of the Heisenberg group. This profound idea was later retrieved by Moyal, who have seen that the symbols of the commutators or of the products of operators are of the the form of $`sine`$ (or $`exponential`$) functions of the bidifferential operators (of the Poisson brackets) of the corresponding symbols.
In the early 70’s Berezin has treated the general mathematical definition of quantization as a kind of a functor from the category of classical mechanics to a certain category of associative algebras. About the same time as F. A. Berezin, M. Flato, M. G. Fronsdal, F. Bayen, A. Lichnerowicz and D. Sternheimer considered quantization as a deformation of the commutative products of classical observables into a noncommutative $``$-products which are parameterized by the Planck constant $`\mathrm{}`$ and satisfy the correspondence principle. They systematically developed the notion of deformation quantization as a theory of $``$-products and gave an independent formulation of quantum mechanics based on this notion (see\[RT\]).
It was proved by Gerstenhaber that a formal deformation quantization exists on an arbitrary symplectic manifold, see for example \[F\] for a detailed explaination. It is however formal and quite complecate in general. We would like to simplify it in some particular cases.
From the orbit method, it is well-known that coadjoint orbits are homogeneous symplectic manifolds with respect to the natural Kirillov structure form on coadjoint orbits. A natural question is to associate in a reasonable way to these orbits some quantum objects, what could be called quantum co-adjoint orbits. In particular, in \[DH1\] and \[DH2\] we obtained “quantum half-planes” and “quantum punctured complex planes”, associated with the affine transformation groups of the real or complex straightlines. In this paper we will therefore continue to realize the problem for the real diamond Lie group. This group has a lot of nontrivial 2-dimensional coadjoint orbits, which are the half-planes, the hyperbolic cylinders and the hyperbolic paraboloids. We should find out explicit formulas for each of these orbits. Our main result therefore is the fact that by using $``$-product we can construct the corresponding quantum half-plans, quantum hyperbolic cylinders, quantum hyperbolic paraboloids and by an exact computation we can find out explicit $``$-product formulas and then, the complete list of irreducible unitary representations of this group. It is useful to do here a remark that there is a general theory for exponetial and compact groups. But our consideration concerning with non-exponential and noncompact Lie group and associated $`G`$-homogeneous symplectic manifolds.
Let us in few words describe the structure od the paper. We introduce some preliminary results in §2. Then, the adapted chart and in particular, Hamiltonian functions in canonical coordinates of the co-adjoint orbit $`\mathrm{\Omega }_F`$ are exposed in §3. The operators $`\widehat{\mathrm{}}_A`$ which define the representations of the real diamond Lie algebra are constructed in §4 and finally, by exponentiating them, we obtain the corresponding unitary representations of the real diamond Lie group $`_3`$.
## 2. Preliminary results
The so called real diamond Lie algebra is the 4-dimensional solvable Lie algebra $`𝔤`$ with basis X, Y, Z, T satisfying the following commutation relations:
$$[X,Y]=Z,[T,X]=X,[T,Y]=Y,$$
$$[Z,X]=[Z,Y]=[T,Z]=0.$$
These relations show that this real diamond Lie algebra $`𝔥_3`$ is an extension of the one-dimensional Lie algebra $`T`$ by the Heisenberg algebra $`𝔥_3`$ with basis $`X,Y,Z`$, where the action of T on Heisenberg algebra $`𝔥_3`$ is defined by the matrix
$$ad_T=\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 0\end{array}\right).$$
We introduce the following notations. The real diamond Lie algebra is isomorphic to $`^4`$ as vector spaces. The coordinates in this standard basis is denote by $`(a,b,c,d)`$. We identify its dual vector space $`𝔤^{}`$ with $`^4`$ with the help of the dual basis $`X^{},Y^{},Z^{},T^{}`$ and with the local coordinates as $`(\alpha ,\beta ,\gamma ,\delta )`$. Thus, the general form of an element of $`𝔤`$ is $`U=aX+bY+cZ+dT`$ and the general form of an element of $`𝔤^{}`$ is $`F=\alpha X^{}+\beta Y^{}+\gamma Z^{}+\delta T^{}`$. The co-adjoint action of $`G=_3`$ on $`𝔤^{}`$ is given (see e.g. \[Ki1\]) by
$$K(g)F,U=F,Ad(g^1)U,F𝔤^{},gG\text{ and }U𝔤=Lie(_3).$$
Denote the co-adjoint orbit of $`G`$ in $`𝔤`$, passing through $`F`$ by
$$\mathrm{\Omega }_F=K(G)F:=\{K(g)F|gG\}.$$
By a direct computation one obtains (see \[D\]):
* Each point of the line $`\alpha =\beta =\gamma =0`$ is a 0-dimensional co-adjoint orbit
(1)
$$\mathrm{\Omega }^1=\mathrm{\Omega }_{(0,0,0,\delta )}.$$
* The set $`\alpha 0,\beta =\gamma =0`$ is union of 2-dimensional co-adjoint orbits ,which are just the half-planes
(2)
$$\mathrm{\Omega }^2=\{(x,0,0,t)|x,t,\alpha x>0\}.$$
* The set $`\alpha =\gamma =0,\beta 0`$ is a union of 2-dimensional co-adjoint orbits, which are half-planes
(3)
$$\mathrm{\Omega }^3=\{(0,y,0,t)|y,t,\beta y>0\}.$$
* The set $`\alpha \beta 0,\gamma =0`$ is decomposed into a family of 2-dimensional co-adjoint orbits, which are hyperbolic cylinders
(4)
$$\mathrm{\Omega }^4=\{(x,y,0,t)|x,y,t\&\alpha x>0,\beta y>0,xy=\alpha \beta \}.$$
* The open set$`\gamma 0`$ is decomposed into a family of 2-dimensional co-adjoint orbits ,which are just the hyperbolic paraboloids
(5)
$$\mathrm{\Omega }^5=\{(x,y,\gamma ,t)|x,y,t\&xy\alpha \beta =\gamma (t\delta )\}.$$
Thus, the real diamond Lie algebra belongs to the class of $`MD_4`$-algebras , i.e. every K-orbit of the corresponding Lie group has dimension 0 or maximal (see \[D\]).
Let us consider now the problem of deformation quantization on half-planes, hyperbolic cylinders, hyperbolic paraboloids. In order to do this, we shall construct on each of these orbits a canonical Darboux coordinate system $`(p,q)`$ and a class of Hamiltonian functions in these coordinates.
## 3. Hamiltonian functions in canonical coordinates of the orbits $`\mathrm{\Omega }_F`$
Each element $`A𝔤`$ can be considered as the restriction of the corresponding linear functional $`\stackrel{~}{A}`$ onto co-adjoint orbits, considered as a subset of $`g^{}`$ ,$`\stackrel{~}{A}(F)=F,A`$. It is well-known that this function is just the Hamiltonnian function, associated with the Hamiltonian vector field $`\xi _A`$, defined by the formula
$$(\xi _Af)(x):=\frac{d}{dt}f(x\mathrm{exp}(tA))|_{t=0},fC^{\mathrm{}}(\mathrm{\Omega }_F).$$
It is well-known the relation $`\xi _A(f)=\{\stackrel{~}{A},f\},fC^{\mathrm{}}(\mathrm{\Omega }_F)`$. Denote by $`\psi `$ the symplectomorphism from $`^2`$ onto $`\mathrm{\Omega }_F`$
$$(p,q)^2\psi (p,q)\mathrm{\Omega }_F,$$
we have:
###### Proposition 3.1.
1. Hamiltonian function $`\stackrel{~}{A}`$ in canonical coordinates $`(p,q)`$ of the orbit $`\mathrm{\Omega }_F`$ is of the form
$$\stackrel{~}{A}\psi (p,q)=\{\begin{array}{cc}dp+a\alpha e^q,\hfill & \text{if }\mathrm{\Omega }_F=\mathrm{\Omega }^2\hfill \\ dp+b\beta e^q,\hfill & \text{if }\mathrm{\Omega }_F=\mathrm{\Omega }^3\hfill \\ dp+a\alpha e^q+b\beta e^q,\hfill & \text{if }\mathrm{\Omega }_F=\mathrm{\Omega }^4\hfill \\ (d\pm b\gamma e^q)p\pm ae^q\pm b(\alpha \beta \gamma \delta )e^q+c\gamma ,\hfill & \text{if }\mathrm{\Omega }_F=\mathrm{\Omega }^5\hfill \end{array}$$
2. In the canonical coordinates $`(p,q)`$ of the orbit $`\mathrm{\Omega }_F`$, the Kirillov form $`\omega `$ is coincided with the standard form $`dpdq`$.
Proof. 1. We adapt the diffeomorphism $`\psi `$ to each of the following cases (for 2-dimensional co-adjoint orbits, only)
* With $`\alpha 0,\beta =\gamma =0`$
$$(p,q)^2\psi (p,q)=(\alpha e^q,0,0,p)\mathrm{\Omega }^2$$
Element $`F𝔤^{}`$ is of the form $`F=\alpha X^{}+\beta Y^{}+\gamma Z^{}+\delta T^{}`$, hence the value of the function $`f_A=\stackrel{~}{A}`$ on the element $`A=aX+bY+cZ+dT`$ is $`\stackrel{~}{A}(F)=F,A=`$
$$\alpha X^{}+\beta Y^{}+\gamma Z^{}+\delta T^{},aX+bY+cZ+dT=\alpha a+\beta b+\gamma c+\delta d.$$
It follows that
(6)
$$\stackrel{~}{A}\psi (p,q)=a\alpha e^q+dp,$$
* With $`\alpha =\gamma =0,\beta 0`$,
$$(p,q)^2\psi (p,q)=(0,\beta e^q,0,p)\mathrm{\Omega }^3.$$
$`\stackrel{~}{A}(F)=F,A=\alpha a+\beta b+\gamma c+\delta d.`$ From this,
(7)
$$\stackrel{~}{A}\psi (p,q)=b\beta e^q+dp$$
* With $`\alpha \beta 0,\gamma =0`$,
$$(p,q)^2\psi (p,q)=(\alpha e^q,\beta e^q,0,p)\mathrm{\Omega }^4.$$
(8)
$$\stackrel{~}{A}\psi (p,q)=a\alpha e^q+b\beta e^q+dp$$
* At last, if $`\gamma 0`$, we consider the orbit with the first coordinate $`x>0`$
$$(p,q)^2\psi (p,q)=(e^q,(\alpha \beta +\gamma p\gamma \delta )e^q,\gamma ,p)\mathrm{\Omega }^5.$$
We have
(9)
$$\stackrel{~}{A}\psi (p,q)=ae^q+b(\alpha \beta +\gamma p\gamma \delta )e^q+c\gamma +dp=$$
$$=(d+b\gamma e^q)p+ae^q+b(\alpha \beta \gamma \delta )e^q+c\gamma .$$
The case $`x<0`$ is similarly treated:
$$(p,q)^2\psi (p,q)=(e^q,(\alpha \beta +\gamma p\gamma \delta )e^q,\gamma ,p)\mathrm{\Omega }^5.$$
(10)
$$\stackrel{~}{A}\psi (p,q)=ae^qb(\alpha \beta +\gamma p\gamma \delta )e^q+c\gamma +dp=$$
$$=(db\gamma e^q)pae^qb(\alpha \beta \gamma \delta )e^q+c\gamma .$$
2. We consider only the following case (the rest are similar):
$$(p,q)^2\psi (p,q)=(e^q,(\alpha \beta +\gamma p\gamma \delta )e^q,\gamma ,p)\mathrm{\Omega }^5.$$
$$\stackrel{~}{A}\psi (p,q)=(d+b\gamma e^q)p+ae^q+b(\alpha \beta \gamma \delta )e^q+c\gamma .$$
In canonical Darboux coordinates $`(p,q)`$ ,
$$F^{}=e^qX^{}+(\alpha \beta +\gamma p\gamma \delta )e^qY^{}+\gamma Z^{}+pT^{}\mathrm{\Omega }^5,$$
and for $`A=aX+bY+cZ+dT,B=a^{}X+b^{}Y+c^{}Z+d^{}T`$, we have $`F^{},[A,B]=`$
$$=e^qX^{}+(\alpha \beta +\gamma p\gamma \delta )e^qY^{}+\gamma Z^{}+pT^{},(ad^{}da^{})X+(db^{}bd^{})Y+(ab^{}ba^{})Z.$$
It follows therefore that
(11)
$$F^{},[A,B]=(ad^{}da^{})e^q+(db^{}bd^{})(\alpha \beta +\gamma p\gamma \delta )e^q+\gamma (ab^{}ba^{}).$$
On the other hand,
$$\xi _A(f)=\{\stackrel{~}{A},f\}=(d+b\gamma e^q)\frac{f}{q}[ae^q+b(\alpha \beta +\gamma p\gamma \delta )e^q]\frac{f}{p}$$
$$\xi _B(f)=\{\stackrel{~}{B},f\}=(d^{}+b^{}\gamma e^q)\frac{f}{q}[a^{}e^q+b^{}(\alpha \beta +\gamma p\gamma \delta )e^q]\frac{f}{p}.$$
From this, consider two vector fields
$$\xi _A=(d+b\gamma e^q)\frac{}{q}[ae^q+b(\alpha \beta +\gamma p\gamma \delta )e^q]\frac{}{p},$$
$$\xi _B=(d^{}+b^{}\gamma e^q)\frac{}{q}[a^{}e^q+b^{}(\alpha \beta +\gamma p\gamma \delta )e^q]\frac{}{p}.$$
We have
(12)
$$\xi _A\xi _B=(d+b\gamma e^q)(d^{}+b^{}\gamma e^q)\frac{}{q}\frac{}{q}+$$
$$+[(ad^{}da^{})e^q+(db^{}d^{}b)(\alpha \beta +\gamma p\gamma \delta )e^q]\frac{}{p}\frac{}{q}+$$
$$+[ae^q+b(\alpha \beta +\gamma p\gamma \delta )e^q][a^{}e^q+b^{}(\alpha \beta +\gamma p\gamma \delta )e^q]\frac{}{p}\frac{}{p}$$
From (11) and (12) we conclude that in the canonical coordinates the Kirillov form is just the standard symplectic form $`\omega =dpdq`$. The proposition is therefore proved. $`\mathrm{}`$
###### DEFINITION 3.2.
Each chart $`\psi ^1`$ on $`\mathrm{\Omega }_F`$ which satisfy 1. and 2. of proposition 3.1 is called an adapted chart on $`\mathrm{\Omega }_F`$.
In the next section we shall see that each adapted chart carries the Moyal $``$-product from $`^2`$ onto $`\mathrm{\Omega }_F`$.
## 4. Moyal $``$-product and representations of $`G=_3.`$
Let us denote by $`\mathrm{\Lambda }`$ the 2-tensor associated with the Kirillov standard form $`\omega =dpdq`$ in canonical Darboux coordinates. Let us consider the well-known Moyal $``$-product of two smooth functions $`u,vC^{\mathrm{}}(^{2n})`$ (see e.g \[AC1\],\[DH1\]), defined by
$$uv=u.v+\underset{r1}{}\frac{1}{r!}(\frac{1}{2i})^rP^r(u,v),$$
where
$$P^1(u,v)=\{u,v\}$$
$$P^r(u,v):=\mathrm{\Lambda }^{i_1j_1}\mathrm{\Lambda }^{i_2j_2}\mathrm{}\mathrm{\Lambda }^{i_rj_r}_{i_1i_2\mathrm{}i_r}^ru_{j_1j_2\mathrm{}j_r}^rv,$$
with
$$_{i_1i_2\mathrm{}i_r}^r:=\frac{^r}{x^{i_1}\mathrm{}x^{i_r}};x:=(p,q)=(p_1,\mathrm{},p_n,q^1,\mathrm{},q^n)$$
using multi-index notation. It is well-known that this series converges in the Schwartz distribution spaces $`𝒮(^{2n})`$. Furthermore, it was obtained the results (see e.g \[AC1\]): If $`u,v𝒮(^{2n})`$, then
* $`\overline{u}\overline{v}=\overline{vu}`$
* $`(uv)(\xi )𝑑\xi =uv𝑑\xi `$
* $`\mathrm{}_u:𝒮(^{2n})𝒮(^{2n})`$, defined by $`\mathrm{}_u(v)=uv`$ is continuous in $`L^2(^{2n},d\xi )`$ and then can be extended to a bounded linear operator (still denoted by $`\mathrm{}_u`$ ) on $`L^2(^{2n},d\xi )`$.
We apply this to the special case $`n=1`$, $`x=(x^1,x^2)=(p,q)`$
###### Proposition 4.1.
In the above mentioned canonical Darboux coordinates $`(p,q)`$ on the orbit $`\mathrm{\Omega }_F`$, the Moyal $``$-product satisfies the relation
$$i\stackrel{~}{A}i\stackrel{~}{B}i\stackrel{~}{B}i\stackrel{~}{A}=i\stackrel{~}{[A,B]},A,B𝔤=Lie(_3).$$
Proof. We prove the proposition for the orbit $`\mathrm{\Omega }^5`$, $`\stackrel{~}{A}=(d+b\gamma e^q)p+ae^q+b(\alpha \beta \gamma \delta )e^q+c\gamma `$ (the other cases are proved similar). Consider the elements $`A=aX+bY+cZ+dT,B=a^{}X+b^{}Y+c^{}Z+d^{}T`$, . Then as said above, the corresponding Hamiltonian functions are
$$\stackrel{~}{A}=(d+b\gamma e^q)p+ae^q+b(\alpha \beta \gamma \delta )e^q+c\gamma $$
$$\stackrel{~}{B}=(d^{}+b^{}\gamma e^q)p+a^{}e^q+b^{}(\alpha \beta \gamma \delta )e^q+c^{}\gamma $$
It is easy then to see that
$`P^0(\stackrel{~}{A},\stackrel{~}{B})=\stackrel{~}{A}.\stackrel{~}{B}`$
$`P^1(\stackrel{~}{A},\stackrel{~}{B})=\{\stackrel{~}{A},\stackrel{~}{B}\}=_p\stackrel{~}{A}_q\stackrel{~}{B}_q\stackrel{~}{A}_p\stackrel{~}{B}=`$
$`=(d+b\gamma e^q)[a^{}e^q+b^{}(\alpha \beta +\gamma p\gamma \delta )e^q]`$
$`(d^{}+b^{}\gamma e^q)[ae^q+b(\alpha \beta +\gamma p\gamma \delta )e^q]=`$
$`=[(ad^{}da^{})e^q+(db^{}d^{}b)(\alpha \beta +\gamma p\gamma \delta )e^q+(ab^{}ba^{})\gamma ]`$
$`P^2(\stackrel{~}{A},\stackrel{~}{B})=\mathrm{\Lambda }^{12}\mathrm{\Lambda }^{12}_{pp}^2\stackrel{~}{A}_{qq}^2\stackrel{~}{B}+\mathrm{\Lambda }^{12}\mathrm{\Lambda }^{21}_{pq}^2\stackrel{~}{A}_{qp}^2\stackrel{~}{B}+\mathrm{\Lambda }^{21}\mathrm{\Lambda }^{12}_{qp}^2\stackrel{~}{A}_{pq}^2\stackrel{~}{B}+`$
$`+\mathrm{\Lambda }^{21}\mathrm{\Lambda }^{21}_{qq}^2\stackrel{~}{A}_{pp}^2\stackrel{~}{B}=2bb^{}\gamma ^2e^{2q}`$
$`P^3(\stackrel{~}{A},\stackrel{~}{B})=\mathrm{\Lambda }^{12}\mathrm{\Lambda }^{12}\mathrm{\Lambda }^{12}_{ppp}^3\stackrel{~}{A}_{qqq}^3\stackrel{~}{B}+\mathrm{\Lambda }^{12}\mathrm{\Lambda }^{12}\mathrm{\Lambda }^{21}_{ppq}^3\stackrel{~}{A}_{qqp}^3\stackrel{~}{B}+`$
$`+\mathrm{\Lambda }^{12}\mathrm{\Lambda }^{21}\mathrm{\Lambda }^{12}_{pqp}^3\stackrel{~}{A}_{qpq}^3\stackrel{~}{B}+\mathrm{\Lambda }^{21}\mathrm{\Lambda }^{12}\mathrm{\Lambda }^{12}_{qpp}^3\stackrel{~}{A}_{pqq}^3\stackrel{~}{B}+`$
$`+\mathrm{\Lambda }^{21}\mathrm{\Lambda }^{21}\mathrm{\Lambda }^{12}_{qqp}^3\stackrel{~}{A}_{ppq}^3\stackrel{~}{B}+\mathrm{\Lambda }^{21}\mathrm{\Lambda }^{12}\mathrm{\Lambda }^{21}_{qpq}^3\stackrel{~}{A}_{pqp}^3\stackrel{~}{B}+`$
$`+\mathrm{\Lambda }^{12}\mathrm{\Lambda }^{21}\mathrm{\Lambda }^{21}_{pqq}^3\stackrel{~}{A}_{qpp}^3\stackrel{~}{B}+\mathrm{\Lambda }^{21}\mathrm{\Lambda }^{21}\mathrm{\Lambda }^{21}_{qqq}^3\stackrel{~}{A}_{ppp}^3\stackrel{~}{B}=0.`$
By analogy, we have
$$P^k(\stackrel{~}{A},\stackrel{~}{B})=0,k4.$$
Thus,
$$i\stackrel{~}{A}i\stackrel{~}{B}i\stackrel{~}{B}i\stackrel{~}{A}=\frac{1}{2i}[P^1(i\stackrel{~}{A},i\stackrel{~}{B})P^1(i\stackrel{~}{B},i\stackrel{~}{A})]$$
$$=i[(ad^{}da^{})e^q+(db^{}d^{}b)(\alpha \beta +\gamma p\gamma \delta )e^q+(ab^{}a^{}b)\gamma ].$$
On the other hand, as
$$[A,B]=[aX+bY+cZ+dT,a^{}X+b^{}Y+c^{}Z+d^{}T]$$
$$=(ad^{}da^{})X+(db^{}d^{}b)Y+(ab^{}a^{}b)Z$$
we obtain
$$i[(ad^{}da^{})e^q+(db^{}d^{}b)(\alpha \beta +\gamma p\gamma \delta )e^q+(ab^{}a^{}b)\gamma ]$$
$$=i\stackrel{~}{[A,B]}=i\stackrel{~}{A}i\stackrel{~}{B}i\stackrel{~}{B}i\stackrel{~}{A}.$$
The proposition is hence proved. $`\mathrm{}`$
Consequently, to each adapted chart, we associate a $`G`$-covariant $``$-product.Then there exists a representation $`\tau `$ of $`G`$ in $`AutN[[\nu ]]`$ ,(see \[G\]) such that (here $`\nu =\frac{i}{2}`$):
$$\tau (g)(uv)=\tau (g)u\tau (g)v.$$
For each $`ALie(_3)`$, the corresponding Hamiltonian function is $`\stackrel{~}{A}`$ and we can put $`\mathrm{}_A(u)=i\stackrel{~}{A}u`$,$`uL^2(^2,\frac{dpdq}{2\pi })^{\mathrm{}}`$. It is then continuated to the whole space $`L^2(^2,\frac{dpdq}{2\pi })`$. Because of the relation in Proposition (4.1), we have
###### Corollary 4.2.
(13)
$$\mathrm{}_{[A,B]}=\mathrm{}_A\mathrm{}_B\mathrm{}_B\mathrm{}_A:=[\mathrm{}_A,\mathrm{}_B]^{}$$
This implies that the correspondence $`ALie(_3)\mathrm{}_A=i\stackrel{~}{A}.`$ is a representation of the Lie algebra $`Lie(_3)`$ on the space $`N[[\frac{i}{2}]]`$ of formal power series in the parameter $`\nu =\frac{i}{2}`$(i.e $`\mathrm{}=1`$) with coefficients in $`N=C^{\mathrm{}}(M,)`$ \[G\].
Let us denote by $`_p(f)`$ the partial Fourier transform of the function $`f`$ from the variable $`p`$ to the variable $`x`$(see e.g\[MV\]), i.e.
$$_p(f)(x,q):=\frac{1}{\sqrt{2\pi }}_{}e^{ipx}f(p,q)𝑑p.$$
Let us denote by $`_p^1(f)(p,q)`$ the inverse Fourier transform.
###### Lemma 4.3.
1.$`_p_p^1(f)=i_p^1(x.f)`$
2.$`_p(p.v)=i_x_p(v)`$
3.$`k2`$ ,then $`P^k(\stackrel{~}{A},_p^1(f))=`$
$$=\{\begin{array}{cc}a\alpha e^q_{p\mathrm{}p}^k_p^1(f)\hfill & \text{if }\stackrel{~}{A}\text{ is defined by (6) }\hfill \\ (1)^kb\beta e^q_{p\mathrm{}p}^k_p^1(f)\hfill & \text{if }\stackrel{~}{A}\text{ is defined by (7)}\hfill \\ [a\alpha e^q+(1)^kb\beta e^q]_{p\mathrm{}p}^k_p^1(f)\hfill & \text{if }\stackrel{~}{A}\text{ is defined by (8)}\hfill \\ (1)^{k1}k.b\gamma e^q_{qp\mathrm{}p}^k_p^1(f)+\hfill & \\ +[ae^q+(1)^kb(\alpha \beta +\gamma p\gamma \delta )e^q]_{p\mathrm{}p}^k_p^1(f)\hfill & \text{if }\stackrel{~}{A}\text{ is defined by (9) }\hfill \end{array}$$
Proof. The first two formulas are well-known from theory of Fourier transforms.
Let us prove 3. Remark that $`\mathrm{\Lambda }=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ in the standard symplectic Darboux coordinates $`(p,q)`$ on the orbit $`\mathrm{\Omega }_F`$, then
* If $`\stackrel{~}{A}=a\alpha e^q+dp`$
$`P^2(\stackrel{~}{A},_p^1(f))=\mathrm{\Lambda }^{12}\mathrm{\Lambda }^{12}_{pp}^2\stackrel{~}{A}_{qq}^2_p^1(f))+\mathrm{\Lambda }^{12}\mathrm{\Lambda }^{21}_{pq}^2\stackrel{~}{A}_{qp}^2_p^1(f))+`$
$`\mathrm{\Lambda }^{21}\mathrm{\Lambda }^{12}_{qp}^2\stackrel{~}{A}_{pq}^2_p^1(f))+\mathrm{\Lambda }^{21}\mathrm{\Lambda }^{21}_{qq}^2\stackrel{~}{A}_{pp}^2_p^1(f))=a\alpha e^q^2_{pp}^1_p(f)=`$
$`P^3(\stackrel{~}{A},_p^1(f))=(1)^6a\alpha e^q_{ppp}^3_p^1(f)=a\alpha e^q_{ppp}^3_p^1(f)`$
and $`P^k(\stackrel{~}{A},_p^1(f))=a\alpha e^q_{p\mathrm{}p}^k_p^1(f)k4,`$
* If $`\stackrel{~}{A}=b\beta e^q+dp.`$
$$P^k(\stackrel{~}{A},_p^1(f))=(1)^kb\beta e^q_{p\mathrm{}p}^k_p^1(f)$$
with $`k2`$
* If $`\stackrel{~}{A}=a\alpha e^q+b\beta e^q+dp,`$
$`P^2(\stackrel{~}{A},_p^1(f))=\mathrm{\Lambda }^{12}\mathrm{\Lambda }^{12}_{pp}^2\stackrel{~}{A}_{qq}^2_p^1(f))+\mathrm{\Lambda }^{12}\mathrm{\Lambda }^{21}_{pq}^2\stackrel{~}{A}_{qp}^2_p^1(f))+`$
$`\mathrm{\Lambda }^{21}\mathrm{\Lambda }^{12}_{qp}^2\stackrel{~}{A}_{pq}^2_p^1(f))+\mathrm{\Lambda }^{21}\mathrm{\Lambda }^{21}_{qq}^2\stackrel{~}{A}_{pp}^2_p^1(f))=`$
$`=[a\alpha e^q+(1)^2b\beta e^q]_{pp}^2_p^1(f)`$
$`P^3(\stackrel{~}{A},_p^1(f))=[a\alpha e^q+(1)^3b\beta e^q]_{ppp}^3_p^1(f)`$
By analogy we have
$$P^k(\stackrel{~}{A},_p^1(f))=[a\alpha e^q+(1)^kb\beta e^q]_{p\mathrm{}p}^k_p^1(f)),k3.$$
* If $`\stackrel{~}{A}=(d+b\gamma e^q)p+ae^q+b(\alpha \beta \gamma \delta )e^q+c\gamma ,`$
$`P^2(\stackrel{~}{A},_p^1(f))=\mathrm{\Lambda }^{12}\mathrm{\Lambda }^{12}_{pp}^2\stackrel{~}{A}_{qq}^2_p^1(f))+\mathrm{\Lambda }^{12}\mathrm{\Lambda }^{21}_{pq}^2\stackrel{~}{A}_{qp}^2_p^1(f))+`$
$`\mathrm{\Lambda }^{21}\mathrm{\Lambda }^{12}_{qp}^2\stackrel{~}{A}_{pq}^2_p^1(f))+\mathrm{\Lambda }^{21}\mathrm{\Lambda }^{21}_{qq}^2\stackrel{~}{A}_{pp}^2_p^1(f))=`$
$`=(1)2.b\gamma e^q_{qp}_p^1(f)+[ae^q+(1)^2b(\alpha \beta +\gamma p\gamma \delta )e^q]_{pp}^2_p^1(f).`$
$`P^3(\stackrel{~}{A},_p^1(f))=(1)^2.3b\gamma e^q_{qpp}_p^1(f)+`$
$`+[ae^q+(1)^3b(\alpha \beta +\gamma p\gamma \delta )e^q]_{ppp}^3_p^1(f).`$
From this we also obtain :
$`P^k(\stackrel{~}{A},_p^1(f))=`$
$`(1)^{k1}.k.b\gamma e^q_{qp\mathrm{}p}^k_p^1(f)+`$
$`+[ae^q+(1)^kb(\alpha \beta +\gamma p\gamma \delta )e^q]_{p\mathrm{}p}^k_p^1(f).k3`$
The lemma is therefore proved. $`\mathrm{}`$
We study now the convergence of the formal power series. In order to do this, we look at the $``$-product of $`i\stackrel{~}{A}`$ as the $``$-product of symbols and define the differential operators corresponding to $`i\stackrel{~}{A}`$.
###### Theorem 4.4.
For each $`ALie(_3)`$ and for each compactly supported $`C^{\mathrm{}}`$ function $`fC_0^{\mathrm{}}(^2)`$, putting $`\widehat{\mathrm{}}_A(f):=_p\mathrm{}_A_p^1(f)`$,we have
$$\widehat{\mathrm{}}_A(f)=\{\begin{array}{cc}[d(\frac{1}{2}_q_x)+ia\alpha e^{(q\frac{x}{2})}]f\hfill & \text{if }\stackrel{~}{A}\text{ is defined by (6) }\hfill \\ [d(\frac{1}{2}_q_x)+ib\beta e^{(q\frac{x}{2})}]f\hfill & \text{if }\stackrel{~}{A}\text{ is defined by (7)}\hfill \\ [d(\frac{1}{2}_q_x)+i(a\alpha e^{(q\frac{x}{2})}+b\beta e^{(q\frac{x}{2})})]f\hfill & \text{if }\stackrel{~}{A}\text{ is defined by (8)}\hfill \\ [(d+b\gamma e^{q\frac{x}{2}})(\frac{1}{2}_q_x)]f+\hfill & \\ +i[ae^{(q\frac{x}{2})}+b(\alpha \beta \gamma \delta )e^{q\frac{x}{2}}+c\gamma ]f\hfill & \text{if }\stackrel{~}{A}\text{ is defined by (9) }\hfill \\ [(db\gamma e^{q\frac{x}{2}})(\frac{1}{2}_q_x)]f+\hfill & \\ +i[ae^{(q\frac{x}{2})}b(\alpha \beta \gamma \delta )e^{q\frac{x}{2}}+c\gamma ]f\hfill & \text{if }\stackrel{~}{A}\text{ is defined by (10) }\hfill \end{array}$$
Proof. Applying Lemma (4.3),we have :
1. If $`\stackrel{~}{A}=a\alpha e^q+dp`$ then
$$\widehat{\mathrm{}}_A(f):=_p\mathrm{}_A_p^1(f)=_p(i\stackrel{~}{A}_p^1(f))=i_p\left(\underset{r0}{}\left(\frac{1}{2i}\right)^r\frac{1}{r!}P^r(\stackrel{~}{A},_p^1(f))\right)=$$
$$\begin{array}{cc}& =i_p\{(a\alpha e^q+dp)_p^1(f)+\frac{1}{1!}\frac{1}{2i}[d_q_p^1(f)+a\alpha e^q_p_p^1(f)]+\hfill \\ & +\frac{1}{2!}(\frac{1}{2i})^2.a\alpha e^q_pp^2_p^1(f)+\mathrm{}+\frac{1}{r!}(\frac{1}{2i})^ra\alpha e^q_{p\mathrm{}p}^r_p^1(f)+\mathrm{}\}=\hfill \\ & =i\{a\alpha e^qf+d_p(p._p^1(f))+\frac{1}{1!}\frac{1}{2i}[d_qf+a\alpha e^q_p(_p_p^1(f))]+\hfill \\ & +\frac{1}{2!}(\frac{1}{2i})^2.a\alpha e^q_p(_{pp}^2_p^1(f))++\frac{1}{3!}(\frac{1}{2i})^3.a\alpha e^q_p(_{ppp}^3_p^1(f))+\mathrm{}\hfill \\ & +\frac{1}{r!}(\frac{1}{2i})^r.a\alpha e^q_p(_{p\mathrm{}p}^r_p^1(f))+\mathrm{}\}=\hfill \\ & =d(\frac{1}{2}_q_x)f+ia\alpha e^q[1+\frac{x}{2}+\frac{1}{2!}(\frac{x}{2})^2+\mathrm{}+\frac{1}{r!}(\frac{1}{x})^r+\mathrm{}]f=\hfill \\ & =d(\frac{1}{2}_q_x)f+ia\alpha e^qe^{\frac{x}{2}}f=d(\frac{1}{2}_q_x)f+ia\alpha e^{(q\frac{x}{2})}f\hfill \end{array}$$
2. If $`\stackrel{~}{A}=b\beta e^q+dp`$ then
$$\widehat{\mathrm{}}_A(f)=d(\frac{1}{2}_q_x)f+ib\beta e^{q\frac{x}{2}}f$$
3. For each $`\stackrel{~}{A}=a\alpha e^q+b\beta e^q+dp,`$ we have:
$$\begin{array}{cc}& \widehat{\mathrm{}}_A=i_p\{(a\alpha e^q+b\beta e^q+dp)_p^1(f)+\frac{1}{2i}[d_q_p^1(f)(a\alpha e^q+\hfill \\ & +b\beta e^q)_p_p^1(f)]+\frac{1}{2!}(\frac{1}{2i})^2[a\alpha e^q+(1)^2b\beta e^q]^2_{pp}_p^1(f)+\mathrm{}+\hfill \\ & +\frac{1}{r!}(\frac{1}{2i})^r[a\alpha e^q+(1)^rb\beta e^q]_{p\mathrm{}p}^r_p^1(f)+\mathrm{}\}\hfill \\ & =ia\alpha e^q.f+id_p(p._p^1(f))+ib\beta e^qf+\frac{1}{2}d_qf+\frac{1}{2}a\alpha e^q_p\left(_p_p^1(f)\right)\hfill \\ & \frac{1}{2}b\beta e^q_p\left(_p_p^1(f)\right)+\mathrm{}i\frac{1}{r!}(\frac{1}{2i})^ra\alpha e^q_p\left(_{p\mathrm{}p}^r_p^1(f)\right)+\hfill \\ & +i\frac{1}{r!}(\frac{1}{2i})^rb\beta e^q_p\left(_{p\mathrm{}p}^r_p^1(f)\right)+\mathrm{}=d(\frac{1}{2}_q_x)+ia\alpha e^q[1+\frac{x}{2}+\mathrm{}\hfill \\ & +\frac{1}{r!}(\frac{x}{2})^r+\mathrm{}]+ib\beta e^q[1+(\frac{x}{2})+\mathrm{}+\frac{1}{r!}(\frac{x}{2})^r+\mathrm{}]\hfill \\ & =d(\frac{1}{2}_q_x)+i[a\alpha e^{(q\frac{x}{2})}+b\beta e^{q\frac{x}{2}}]f.\hfill \end{array}$$
4. For each $`\stackrel{~}{A}`$ is as in (9) , remark that
$$P^0(\stackrel{~}{A},_p^1(f))=\stackrel{~}{A}._p^1(f);$$
$$P^1(\stackrel{~}{A},_p^1(f))=\{\stackrel{~}{A},_p^1(f)\}=$$
$$(d+b\gamma e^q)_q_p^1(f)[ae^q+b(\alpha \beta +\gamma p\gamma \delta )e^q]_p_p^1(f)$$
and applying Lemma (4.3), we obtain:
$$\begin{array}{cc}& \widehat{\mathrm{}}_A(f)=i\{_p\left([dp+ae^q+b(\alpha \beta +\gamma p\gamma \delta )e^q+c\gamma ]_p^1(f)\right)+\hfill \\ & +\frac{1}{2i}\frac{1}{1!}_p\left([d+b\gamma e^q]_q_p^1(f)[ae^q+b(\alpha \beta +\gamma p\gamma \delta )e^q]_p_p^1(f)\right)+\hfill \\ & +(\frac{1}{2i})^2\frac{1}{2!}_p(2b\gamma e^q_{pq}^2_p^1(f)+[ae^q+b(\alpha \beta +\gamma p\gamma \delta )e^q]_{pp}^2(_p^1(f))+\mathrm{}\hfill \\ & +(\frac{1}{2i})^r\frac{1}{r!}_p((1)^{r1}rb\gamma e^q_{p\mathrm{}pq}^r_p^1(f)+(1)^r[(1)^rae^q+b(\alpha \beta +\gamma p\gamma \delta )e^q]\times \hfill \\ & \times _{p\mathrm{}p}^r_p^1(f))+\mathrm{}\}=\hfill \end{array}$$
$$\begin{array}{cc}& =i\{ae^qf+b(\alpha \beta \gamma \delta )e^qf+d_p\left(p_p^1(f)\right)+b\gamma e^q_p\left(p_p^1(f)\right)\hfill \\ & +\frac{1}{2i}\frac{1}{1!}(d+b\gamma e^q)_qf\frac{1}{2i}[ae^qixf+b(\alpha \beta \gamma \delta )e^qixf+b\gamma e^q_p\left(p_p^1(xf)\right)]+\hfill \\ & +(\frac{1}{2i})^2\frac{1}{2!}(2b\gamma e^q)_p\left(_{pq}^2_p^1(f)\right)+(\frac{1}{2i})^2\frac{1}{2!}[ae^q(ix)^2f+b(\alpha \beta \gamma \delta )e^q(ix)^2f+\hfill \\ & +b\gamma e^q_p\left(pi^2_p^1(x^2f)\right)]+\mathrm{}+(\frac{1}{2i})^r\frac{1}{r!}(1)^{r1}rb\gamma e^q^r_{p\mathrm{}pq}^1_p(f)\hfill \\ & +(\frac{1}{2i})^r\frac{1}{r!}[ae^q(ix)^rf+(1)^rb(\alpha \beta \gamma \delta )e^q(ix)^rf+b\gamma e^q_p\left(p(ix)^r_p^1(f)\right)]+\mathrm{}\}\hfill \\ & =i[ae^q(1+\frac{1}{2!}\frac{x}{2}+\mathrm{}+\frac{1}{r!}(\frac{x}{2})^r\mathrm{})f]+i[b(\alpha \beta \gamma \delta )e^q(1\frac{1}{2!}\frac{x}{2}+\mathrm{}+\hfill \\ & +(1)^r\frac{1}{r!}(\frac{x}{2})^r\mathrm{})f]+ic\gamma f+i^2d_xf+\frac{1}{2}d_qf+ib\gamma e^q[i_xf\frac{1}{2i}_p\left(pi_p^1(xf)\right)+\hfill \\ & +\mathrm{}+(\frac{1}{2i})^r\frac{1}{r!}(1)^r_p\left(pi^r_p^1(x^rf)\right)+\mathrm{}]=d(\frac{1}{2}_q_x)f+\hfill \\ & +[iae^{(q\frac{x}{2})}+ib(\alpha \beta \gamma \delta )e^{q\frac{x}{2}}]f+ic\gamma f+\frac{1}{2}e^{\frac{x}{2}}b\gamma e^q_qfb\gamma e^qe^{\frac{x}{2}}_xf\hfill \\ & =(d+b\gamma e^{q\frac{x}{2}})(\frac{1}{2}_q_x)f+[iae^{(q\frac{x}{2})}+ib(\alpha \beta \gamma \delta )e^{q\frac{x}{2}}+ic\gamma ]f\hfill \end{array}$$
5. At last, if $`\stackrel{~}{A}`$ is defined by (10) then :
$$\widehat{\mathrm{}}_A(f)=(db\gamma e^{q\frac{x}{2}})(\frac{1}{2}_q_x)f+[iae^{(q\frac{x}{2})}ib(\alpha \beta \gamma \delta )e^{q\frac{x}{2}}+ic\gamma ]f$$
The theorem is therefore proved. $`\mathrm{}`$
###### Remark 4.5.
Setting new variables $`s=q\frac{x}{2}`$, $`t=q+\frac{x}{2}`$, we have
$$\widehat{\mathrm{}}_A(f)=\{\begin{array}{cc}\left(d_s+ia\alpha e^s\right)f|_{(s,t)}\hfill & \text{if }\stackrel{~}{A}\text{ is defined by (6) }\hfill \\ \left(d_s+ib\beta e^s\right)f|_{(s,t)}\hfill & \text{if }\stackrel{~}{A}\text{ is defined by (7)}\hfill \\ \left(d_s+i[a\alpha e^s+b\beta e^s]\right)f|_{(s,t)}\hfill & \text{if }\stackrel{~}{A}\text{ is defined by (8)}\hfill \\ ((d+b\gamma e^s)_s+\hfill & \\ i[ae^s+b(\alpha \beta \gamma \delta )e^s+c\gamma ])f|_{(s,t)}.\hfill & \text{if }\stackrel{~}{A}\text{ is defined by (9)}\hfill \\ ((db\gamma e^s)_s+\hfill & \\ i[ae^sb(\alpha \beta \gamma \delta )e^s+c\gamma ])f|_{(s,t)}.\hfill & \text{if }\stackrel{~}{A}\text{ is defined by (10)}\hfill \end{array}$$
###### Theorem 4.6.
With above notations we obtain the operators :
$$\widehat{\mathrm{}}_A=\{\begin{array}{cc}\widehat{\mathrm{}}_A^{(2)}=\left(d_s+ia\alpha e^s\right)|_{(s,t)}\hfill & \\ \widehat{\mathrm{}}_A^{(3)}=\left(d_s+ib\beta e^s\right)|_{(s,t)}\hfill & \\ \widehat{\mathrm{}}_A^{(4)}=\left(d_s+i[a\alpha e^s+b\beta e^s]\right)|_{(s,t)}\hfill & \\ \widehat{\mathrm{}}_A^{(5)}=\left((d+b\gamma e^s)_s+i[ae^s+b(\alpha \beta \gamma \delta )e^s+c\gamma ]\right)|_{(s,t)}.\hfill & \\ \widehat{\mathrm{}}_A^{(5^{})}=\left((db\gamma e^s)_s+i[ae^sb(\alpha \beta \gamma \delta )e^s+c\gamma ]\right)|_{(s,t)}\hfill & \end{array}$$
which provides the representations of the Lie algebra $`𝔤`$=$`Lie(_3)`$.
Furthermore, $`A,B𝔤`$,
$$\widehat{\mathrm{}}_A\widehat{\mathrm{}}_B\widehat{\mathrm{}}_B\widehat{\mathrm{}}_A=\widehat{\mathrm{}}_{[A,B]}$$
Proof For each compactly supported $`C^{\mathrm{}}`$ function $`fC_0^{\mathrm{}}(^2)`$ and for $`A,BLie(_3)`$,we have
$$\widehat{\mathrm{}}_{(\mu _1A+\mu _2B)}(f)=_p\mathrm{}_{(\mu _1A+\mu _2B)}_p^1(f)=_p\left(i(\stackrel{~}{\mu _1A+\mu _2B})_p^1\right)=$$
$$=\mu _1_p\mathrm{}_A_p^1(f)+\mu _2_p\mathrm{}_B_p^1(f)=\mu _1\widehat{\mathrm{}}_A(f)+\mu _2\widehat{\mathrm{}}_B(f)\mu _1,\mu _2.$$
Moreover,
$$\widehat{\mathrm{}}_A\widehat{\mathrm{}}_B(f)\widehat{\mathrm{}}_B\widehat{\mathrm{}}_A(f)=\widehat{\mathrm{}}_A\left(_p\mathrm{}_B_p^1(f)\right)\widehat{\mathrm{}}_B\left(_p\mathrm{}_A_p^1(f)\right)=$$
$$=_p(i\stackrel{~}{A}(i\stackrel{~}{B}_p^1(f))_p(i\stackrel{~}{B}(i\stackrel{~}{A}_p^1(f))=_p(i\stackrel{~}{[A,B]}_p^1(f))=\widehat{\mathrm{}}_{[A,B]}(f)$$
$`\mathrm{}`$
###### DEFINITION 4.7.
Let $`\mathrm{\Omega }_F^\lambda `$ be K-orbits of the real diamond Lie group $`G`$. With $`A`$ runs over the Lie algebra $`𝔤=Lie(G)`$,
* $`(\mathrm{\Omega }^2,\widehat{\mathrm{}}_A^{(2)});(\mathrm{\Omega }^3,\widehat{\mathrm{}}_A^{(3)})`$ are called the quantum half-planes,
* $`(\mathrm{\Omega }^4,\widehat{\mathrm{}}_A^{(4)})`$ \- quantum hyperbolic cylinder,
* $`(\mathrm{\Omega }^5,\widehat{\mathrm{}}_A^{(5)},\widehat{\mathrm{}}_A^{(5^{})})`$ \- quantum hyperbolic paraboloid,
with respect to the co-adjoint action of Lie group $`G`$. In the other words, $`(\mathrm{\Omega }_F,\widehat{\mathrm{}}_A)`$, with $`A`$ running over the Lie algebra $`𝔤`$ is called a quantum co-adjoint orbit of Lie group $`G`$.
As G=$`_3`$ is connected and simply connected, we obtain a unitary representations $`T`$ of G defined by the following formula
$$T(\mathrm{exp}A):=\mathrm{exp}(\widehat{\mathrm{}}_A);A𝔤$$
More detail,
$$\mathrm{exp}(\widehat{\mathrm{}}_A)=\{\begin{array}{cc}\mathrm{exp}(d_s+ia\alpha e^s)|_{(s,t)}\hfill & \text{if }\stackrel{~}{A}\text{ is defined by (6) }\hfill \\ \mathrm{exp}(d_s+ib\beta e^s)|_{(s,t)}\hfill & \text{if }\stackrel{~}{A}\text{ is defined by (7)}\hfill \\ \mathrm{exp}(d_s+i[a\alpha e^s+b\beta e^s])|_{(s,t)}\hfill & \text{if }\stackrel{~}{A}\text{ is defined by (8)}\hfill \\ \mathrm{exp}((d+b\gamma e^s)_s+\hfill & \\ i[ae^s+b(\alpha \beta \gamma \delta )e^s+c\gamma ])|_{(s,t)}.\hfill & \text{if }\stackrel{~}{A}\text{ is defined by (9)}\hfill \\ \mathrm{exp}((db\gamma e^s)_s+\hfill & \\ i[ae^sb(\alpha \beta \gamma \delta )e^s+c\gamma ])|_{(s,t)}.\hfill & \text{if }\stackrel{~}{A}\text{ is defined by (10)}\hfill \end{array}$$
This means that we refind all the representations $`T(\mathrm{exp}A)`$ of the real diamond Lie group $`_3`$, those could implicitly obtained from (induction) orbit method induction. What we did here gives us more precise analytic formulas in this case for orbit method induction.
ACKNOWLEDGMENT
The author would like to express his gratitude to Professor Do Ngoc Diep for all his helpfulness and for suggesting many of the topics considered in this paper. The author also thanks Dr. Nguyen Viet Dung for his encouragement. |
warning/0001/math0001190.html | ar5iv | text | # Multi-variable Polynomial Solutions to Pell’s Equation and Fundamental Units in Real Quadratic Fields
## 1. Introduction
Solving Pell’s equation is of relevance in finding fundamental units in real quadratic fields and for this reason polynomial solutions are interesting in that they can supply the fundamental units in infinite families of such fields.
There have been several papers written over the past thirty years which describe certain polynomials whose square roots have periodic continued fraction expansions which can be written down explicitly in terms of the coefficients and variables of the polynomials. See for example the papers of Bernstein , Levesque and Rhin , Madden , Van der Poorten and Van der Poorten and Williams .
In this paper an algorithm is described which allows one to construct, for each positive integer $`n`$, a finite collection of multi-variable Fermat-Pell polynomials which have *all* positive integers whose square-roots have a continued fraction expansion of period $`n+1`$ in their range. If $`F_i:=F_i(t_0,t_1,\mathrm{},t_{\frac{n+1}{2}})`$ is any one of these polynomials, the fundamental polynomial solution to the equation
(1.1)
$$C_i^2F_iH_i^2=(1)^{n1}$$
(where $`C_i`$ and $`H_i`$ are polynomials in the variables $`t_0,t_1,\mathrm{},t_{\frac{n+1}{2}}`$) can be found. Moreover, the continued fraction expansion of $`\sqrt{F_i}`$ can be written down when $`t_1,\mathrm{},t_{\frac{n+1}{2}}0`$ and $`t_0>g_i(t_1,\mathrm{},t_{\frac{n+1}{2}})`$, a certain rational function of these variables. Some implications for single-variable Fermat-Pell polynomials are discussed as are the implications for writing down the fundamental units in a wide class of real quadratic number fields.
Definition: a multi-variable polynomial
$$F:=F(t_0,t_1,\mathrm{},t_k)[t_0,t_1,\mathrm{},t_k],k1$$
is called a *multi-variable Fermat-Pell polynomial* <sup>1</sup><sup>1</sup>1These polynomials are called “Fermat-Pell polynomials” here to avoid confusion with “Pell Polynomials” and also because Fermat investigated the “Pell” equation. if there exists polynomials
$$C:=C(t_0,t_1,\mathrm{},t_k)\text{ and }H:=H(t_0,t_1,\mathrm{},t_k)[t_0,t_1,\mathrm{},t_k]$$
such that either
(1.2) $`C^2FH^2=1,\text{for all }t_i,\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}ik\text{,\hspace{0.17em}\hspace{0.17em}\hspace{0.17em} or}`$
$`C^2FH^2=1,\text{for all }t_i,\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}ik\text{.}`$
Such a triple of polynomials $`\{C,H,F\}`$ satisfying equation (1.2) constitute a *multi-variable polynomial solution* to Pell’s equation.
Definition: The multi-variable Fermat-Pell polynomial $`F`$ (as above) is said to have a *multi-variable polynomial continued fraction expansion* if there exists a positive integer $`n`$, a real constant $`T`$, a rational function $`g(t_1,\mathrm{},t_k)(t_1,\mathrm{},t_k)`$ and polynomials $`a_0:=a_0(t_0,t_1,\mathrm{},t_k)[t_0,t_1,\mathrm{},t_k]`$ and $`a_j:=a_j(t_1,\mathrm{},t_k)[t_1,\mathrm{},t_k],\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}1}jn`$, which take only positive integral values for integral $`t_iT,1ik`$ and (possibly half-) integral $`t_0>g(t_1,\mathrm{},t_k)`$ such that
$$\sqrt{F}=[a_0;\overline{a_1,\mathrm{},a_n,2a_0}],\text{ for all }t_i\text{’s in the ranges stated},\mathrm{\hspace{0.17em}0}ik.$$
Remarks:
(1) From the point of view of simplicity it would be desirable to replace the condition $`t_0g(t_1,\mathrm{},t_k)`$ by $`t_0T`$ but it will be seen that for the polynomials examined here that the former condition is more natural and indeed cannot be replaced by the latter condition.
(2)The restriction that the $`a_i(t_1,\mathrm{},t_k)0`$, $`1in`$ may also seem artificial to some since negative terms can easily be removed from a continued fraction expansion (see, for example ) but this changes the period of the continued fraction so is avoided here.
(3) It may also seem artificial to have $`a_0`$ depend on a variable $`t_0`$ while the other $`a_i`$’s do not but this will also be seen to occur naturally.
(4)Finally, allowing $`t_0`$ to take half-integral values in some circumstances may also seem strange but this also will be seen to be natural and indeed necessary.
Definition: If, for all sets of integers $`\{t_0^{},t_1^{},\mathrm{},t_k^{}\}`$ satisfying $`t_0^{}g(t_1^{},\mathrm{},t_k^{})`$ and $`t_i^{}T`$, $`1ik`$,
$$X=C_i(t_0^{},t_1^{},\mathrm{},t_k^{}),Y=H_i(t_0^{},t_1^{},\mathrm{},t_k^{})$$
constitutes the fundamental solution (in integers) to
$$X^2F_i(t_0^{},t_1^{},\mathrm{},t_k^{})Y^2=(1)^{n1}$$
then $`(C_i(t_0,t_1,\mathrm{},t_k),H_i(t_0,t_1,\mathrm{},t_k))`$ is termed the *fundamental polynomial solution* to equation (1.1).
Standard notations are used:
$$a_0+\frac{1}{a_1+}\frac{1}{a_2+}\frac{1}{a_3+}\mathrm{}\frac{1}{a_N}:=a_0+\frac{1}{a_1+{\displaystyle \frac{1}{a_2+{\displaystyle \frac{1}{a_3+\mathrm{}{\displaystyle \frac{1}{a_N}}}}}}}.$$
To save space this continued fraction is usually written $`[a_0;a_1,\mathrm{},a_n]`$. The infinite periodic continued fraction with initial non-periodic part $`a_0`$ and periodic part $`a_1,\mathrm{},a_n,2a_0`$ is denoted by $`[a_0;\overline{a_1,\mathrm{}\mathrm{}.,a_n,2a_0}]`$. The $`i`$-th approximant of the continued fraction $`[a_0;a_1,\mathrm{},]`$ is denoted by $`P_i/Q_i`$.
Repeated use will be made of some basic facts about continued fractions, such as:
(1.3) $`P_nQ_{n1}P_{n1}Q_n=(1)^{n1},`$
$`P_{n+1}=a_{n+1}P_n+P_{n1},`$
$`Q_{n+1}=a_{n+1}Q_n+Q_{n1},`$
each of these relations being valid for $`n=1,2,3\mathrm{}`$.
Before coming to the main problem, it is necessary to first solve a related problem on symmetric strings of positive integers.
## 2. A problem concerning Symmetric Sequences
Question: For which symmetric sequences of positive integers $`a_1,\mathrm{},a_n`$ do there exist positive integers $`a_0`$ and $`D`$ such that
(2.1)
$$\sqrt{D}=[a_0;\overline{a_1,\mathrm{}\mathrm{}.,a_n,2a_0}]\mathrm{?}$$
Let $`P_i/Q_i`$ denote the $`i`$th approximant of the continued fraction
(2.2)
$$0+\frac{1}{a_1+}\frac{1}{a_2+}\frac{1}{a_3+}\mathrm{}\frac{1}{a_n}.$$
By the well known correspondence between convergents and matrices
$`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\left(\begin{array}{cc}a_1& 1\\ 1& 0\end{array}\right)\mathrm{}\mathrm{}\left(\begin{array}{cc}a_n& 1\\ 1& 0\end{array}\right)`$ $`=\left(\begin{array}{cc}P_n& P_{n1}\\ Q_n& Q_{n1}\end{array}\right)`$
$`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\left(\begin{array}{cc}a_1& 1\\ 1& 0\end{array}\right)\mathrm{}\mathrm{}\left(\begin{array}{cc}a_n& 1\\ 1& 0\end{array}\right)\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ $`=\left(\begin{array}{cc}P_{n1}& P_n\\ Q_{n1}& Q_n\end{array}\right).`$
Since the left side in the second equation is a symmetric sequence of symmetric matrices it follows that
(2.3)
$$P_n=Q_{n1}.$$
Suppose $`\sqrt{D}=[a_0;\overline{a_1,\mathrm{}\mathrm{}.,a_n,2a_0}]=a_0+\beta `$, where $`\beta =`$$`[\mathrm{\hspace{0.17em}0};\overline{a_1,\mathrm{}\mathrm{}.,a_n,2a_0}]`$ so that
$`\beta =[\mathrm{\hspace{0.17em}0};a_1,\mathrm{}\mathrm{}.,a_n,2a_0+\beta ],`$
$`\beta ={\displaystyle \frac{(2a_0+\beta )P_n+P_{n1}}{(2a_0+\beta )Q_n+Q_{n1}}}={\displaystyle \frac{\beta P_n+(2a_0P_n+P_{n1})}{\beta Q_n+(2a_0Q_n+Q_{n1})}},`$
$`\beta ^2Q_n+(2a_0Q_n+Q_{n1}P_n)\beta (2a_0P_n+P_{n1})=0,`$
$`\beta ^2Q_n+(2a_0Q_n)\beta (2a_0P_n+P_{n1})=0,\text{(by (}\text{2.3}\text{))}`$
$`\sqrt{D}=a_0+\beta =\sqrt{a_0^2+{\displaystyle \frac{2a_0P_n+P_{n1}}{Q_n}}}`$
The problem now becomes one of determining for which symmetric sequences of positive integers $`a_1,\mathrm{},a_n`$ does there exist positive integers $`a_0`$ such that $`(2a_0P_n+P_{n1})/Q_n`$ is an integer.
###### Theorem 1.
There exists a positive integer $`a_0`$ such that $`(2a_0P_n+P_{n1})/Q_n`$ is an integer if and only if $`P_{n1}Q_{n1}`$ is even.
###### Proof.
$``$ Suppose first of all that $`P_{n1}Q_{n1}`$ is even. By equation (1.3)
$$P_nQ_{n1}+(1)^n=P_{n1}Q_n$$
(i) Suppose $`n`$ is even. Then $`P_nQ_{n1}P_{n1}+P_{n1}=P_{n1}^2Q_n`$. Choose $`t`$ to be any integer or half-integer such that $`tQ_n`$ is an integer and $`a_0:=Q_{n1}P_{n1}/2+tQ_n>0`$. Then
$$\frac{2a_0P_n+P_{n1}}{Q_n}=\frac{Q_{n1}P_{n1}P_n+2tP_nQ_n+P_{n1}}{Q_n}=2tP_n+P_{n1}^2$$
(ii)Similarly, in the case $`n`$ is odd, $`P_nQ_{n1}P_{n1}+P_{n1}=P_{n1}^2Q_n`$. Choose $`t`$ to be any integer or half-integer such that $`tQ_n`$ is an integer and $`a_0:=Q_{n1}P_{n1}/2+tQ_n>0`$. In this case
$$\frac{2a_0P_n+P_{n1}}{Q_n}=2tP_nP_{n1}^2$$
$``$ Suppose next that $`P_{n1}`$ and $`Q_{n1}`$ are both odd and that there exists a positive integer $`a_0`$ such that $`(2a_0P_n+P_{n1})/Q_n`$ is a positive integer, $`m`$, say. Using (1.3) and (2.3) it follows that $`Q_n`$ is even. Then $`2a_0P_n+P_{n1}=mQ_n`$ implies $`P_{n1}`$ is even - a contradiction. ∎
Remarks:
(i)Note that this process gives all $`a_0`$ such that $`(2a_0P_n+P_{n1})/Q_n`$ is an integer. Indeed,
$`(2a_0P_n+P_{n1})/Q_n=k,\text{ an integer}`$
$`2a_0P_nQ_{n1}=P_{n1}Q_{n1}+kQ_nQ_{n1},`$
$`2a_0(1)^{n1}=2a_0(P_nQ_{n1}P_{n1}Q_n),\text{( by (}\text{1.3}\text{))}`$
$`=P_{n1}Q_{n1}+Q_n(kQ_{n1}2a_0P_{n1}),`$
$`a_0=(1)^{n1}\left({\displaystyle \frac{P_{n1}Q_{n1}}{2}}+Q_n{\displaystyle \frac{kQ_{n1}2a_0P_{n1}}{2}}\right).`$
Notice also that if there is one such $`a_0`$ that there are infinitely many of them.
(ii)Notice that, with $`P_n,P_{n1},Q_n`$ and $`Q_{n1}`$ as defined above, if there exists a positive integer $`D`$ satisfying (2.1) then $`D=p(t_0)`$, for some allowed $`t_0`$, where
$`p(t)=\left({\displaystyle \frac{Q_{n1}P_{n1}}{2}}+tQ_n\right)^2+2tP_n+P_{n1}^2,t>{\displaystyle \frac{Q_{n1}P_{n1}}{2Q_n}},\text{(}n\text{ even),}`$
$`p(t)=\left({\displaystyle \frac{Q_{n1}P_{n1}}{2}}+tQ_n\right)^2+2tP_nP_{n1}^2,t>{\displaystyle \frac{Q_{n1}P_{n1}}{2Q_n}},\text{(}n\text{ odd).}`$
The above theorem suggests a simple algorithm for deciding if, for a given symmetric sequence of positive integers $`a_1,\mathrm{},a_n`$, there exist positive integers $`a_0`$ and $`D`$ such that (2.1) holds. Notice that all that matters is the parity of the $`a_i`$ so all calculations can be done in $`_2`$. First of all define the following matrices:
$$J=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),K=\left(\begin{array}{cc}1& 1\\ 1& 0\end{array}\right)\text{and}I=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
Convert the sequence $`a_1,a_2,\mathrm{},a_n`$ to a sequence of $`J`$\- and $`K`$-matrices, according to whether each $`a_i`$ is odd (replace by a $`K`$) or even (replace by a $`J`$). Prefix a $`J`$-matrix (to account for the initial $`0`$ in the continued fraction (2.2)). Multiply this sequence together (modulo $`2`$) using the facts that $`J^2=K^3=I,\text{ and }JK=K^2J`$.
The final matrix $`\left(\begin{array}{cc}& 1\\ & 1\end{array}\right)mod2`$ there do not exist positive integers $`a_0`$ and $`D`$ such that (2.1) holds.
###### Example 1.
Do there exist positive integers $`a_0`$ and $`D`$ such that
$$\sqrt{D}=[a_0;\overline{22,34,97,32,15,17,17,15,32,97,34,22,2a_0}]\mathrm{?}$$
As described above convert the sequence $`22,34,97,32,15,17,17,15,32,97`$, $`34,22`$ to a sequence of $`J`$\- and $`K`$-matrices, prefix a $`J`$-matrix and multiply the sequence together:
$`\underset{}{JJ}JKJ\underset{}{KKK}KJK\underset{}{JJ}=JK(JK)JK`$ $`=J\underset{}{K(K^2}\underset{}{J)J}K`$
$`=JK=\left(\begin{array}{cc}1& 0\\ 1& 1\end{array}\right).`$
Therefore there do exist positive integers $`a_0`$ and $`D`$ such that
$$\sqrt{D}=[a_0;\overline{22,34,97,32,15,17,17,15,32,97,34,22,2a_0}].$$
## 3. Multi-variable Fermat-Pell Polynomials
Definition: If $`\{a_1,\mathrm{},a_n\}`$ is a symmetric zero-one sequence such that
$$\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\underset{i=1}{\overset{n}{}}\left(\begin{array}{cc}a_i& 1\\ 1& 0\end{array}\right)\left(\begin{array}{cc}& 1\\ & 1\end{array}\right)mod2$$
then the sequence $`\{a_1,\mathrm{},a_n\}`$ is termed a *permissible* sequence. Let $`r(n)`$ denote the number of permissible sequences of length $`n`$.
Note: It is not difficult to show that $`r(2m)=((1)^m+2^{m+1})/3`$ and that $`r(2m+1)=((1)^m+5\times 2^m)/3`$.
If $`D`$ is a positive integer such that $`\sqrt{D}=[a_0;\overline{a_1,\mathrm{}\mathrm{}.,a_n,2a_0}]`$ then $`\{a_1,\mathrm{},a_n\}mod2`$ must equal one of the above permissible sequences and $`D`$ is said to be *associated* with this permissible sequence . The collection of all positive integers associated with a particular permissible sequence is termed the *parity class* of this permissible sequence. Sometimes, if there is no danger of ambiguity, these collections of positive integers will be referred to simply as *parity classes*.
###### Theorem 2.
(i)For each positive integer $`n`$ there exists a finite collection of multi-variable Fermat-Pell polynomials $`F_j(t_0,t_1,\mathrm{},t_{\frac{n+1}{2}}),1jr(n)`$, such that each positive integer whose square root has a continued fraction expansion with period $`n+1`$ lies in the range of exactly one of these polynomials. Moreover, these polynomials can be constructed;
(ii) These polynomials have a polynomial continued fraction expansion which can be explicitly determined;
(iii) The fundamental polynomial solution
$`C=C_j(t_0,t_1,\mathrm{},t_{\frac{n+1}{2}}),H=H_j(t_0,t_1,\mathrm{},t_{\frac{n+1}{2}})`$ to
(3.1)
$$C^2F_j(t_0,t_1,\mathrm{},t_{\frac{n+1}{2}})H^2=(1)^{n1}$$
exists and can be explicitly determined.
###### Proof.
(i)The proof will be by construction.
*Step 1*: Find all permissible sequences. This will involve checking $`2^{\frac{n+1}{2}}`$ zero-one sequences in a way similar to the example (1) above.
*Step 2*: For each permissible sequence $`\{a_1,\mathrm{},a_n\}`$ create a new symmetric polynomial sequence $`\{a_1(t_1),a_2(t_2),\mathrm{},a_{n1}(t_2),a_n(t_1)\}`$ by replacing each $`a_i`$ and its partner $`a_{n+1i}`$ in the symmetric sequence by $`a_i(t_i)=a_{n+1i}(t_i)=2t_i+1`$ if $`a_i=1`$ and by $`a_i(t_i)=a_{n+1i}(t_i)=2t_i+2`$ if $`a_i=0`$. This new sequence will sometimes be referred to as the sequence $`\{a_1,\mathrm{},a_n\}`$, if there is no danger of ambiguity. Each of the integer variables $`t_i`$ (in the polynomial being constructed) will be allowed to vary independently over the range $`0t_i<\mathrm{}`$ and each of the new $`a_i`$’s will keep the same parity and stay positive.
*Step 3* As in (2.2), form the continued fraction
$$0+\frac{1}{a_1(t_1)+}\frac{1}{a_2(t_2)+},\mathrm{},\frac{1}{a_{n1}(t_2)+}\frac{1}{a_n(t_1)}$$
and calculate $`P_n,Q_n,P_{n1}`$ and $`Q_{n1}`$ for this polynomial continued fraction, where these expressions are now polynomials in the $`t_i`$’s.
*Step 4* Construct $`F_j:=F_j(t_0,t_1,\mathrm{},t_{\frac{n+1}{2}})`$, the multi-variable Fermat-Pell polynomial corresponding to the particular parity sequence under consideration. This is simply done by defining
(3.2) $`F_j:=\{\begin{array}{cc}\left(\frac{Q_{n1}P_{n1}}{2}+t_0Q_n\right)^2+2t_0P_n+P_{n1}^2,\text{(}n\text{ even)}\hfill & \\ \left(\frac{Q_{n1}P_{n1}}{2}+t_0Q_n\right)^2+2t_0P_nP_{n1}^2,\text{(}n\text{ odd)}\hfill & \end{array}`$
where $`(1)^{n+1}Q_{n1}P_{n1}/(2Q_n)<t_0<\mathrm{}`$ and $`t_0`$ can take half-integral values if $`Q_n`$ is even and otherwise takes integral values.
Every positive integer whose square root has a continued fraction expansion with period $`n+1`$ lies in the range of exactly one of these polynomials. That these polynomials are multi-variable Fermat-Pell polynomials follows from equation (3.4) below.
(ii) With $`t_0`$ in the range given, then
$$\sqrt{F_j}=[a_0(t_0,t_1,\mathrm{},t_{\frac{n+1}{2}});\overline{a_1(t_1),\mathrm{},a_n(t_1),2a_0(t_0,t_1,\mathrm{},t_{\frac{n+1}{2}})}],$$
for all $`t_i0`$. Here
(3.3) $`a_0=a_0(t_0,t_1,\mathrm{},t_{\frac{n+1}{2}}):=\{\begin{array}{cc}\frac{Q_{n1}P_{n1}}{2}+t_0Q_n,\text{(}n\text{ even)}\hfill & \\ \frac{Q_{n1}P_{n1}}{2}+t_0Q_n,\text{(}n\text{ odd)}.\hfill & \end{array}`$
(iii) Notice (using (1.3) and (2.3)) that
(3.4)
$$(a_0Q_n+P_n)^2(a_0^2+(2a_0P_n+P_{n1})/Q_n)Q_n^2=(1)^{n1}.$$
To see that $`(a_0Q_n+P_n,Q_n)`$ is the fundamental solution to (3.1), notice that
$$\sqrt{F_j}=[a_0(t_0,t_1,\mathrm{},t_{\frac{n+1}{2}});\overline{a_1(t_1),\mathrm{},a_n(t_1),2a_0(t_0,t_1,\mathrm{},t_{\frac{n+1}{2}})}].$$
This has period $`n+1`$ and the $`n`$th approximant is $`a_0+P_n/Q_n=(a_0Q_n+P_n)/Q_n`$ and by the theory of the Pell equation $`(a_0Q_n+P_n,Q_n)`$ is the fundamental solution to (3.1). ∎
As regards fundamental units in quadratic fields there is the following theorem on page 119 of :
###### Theorem 3.
Let $`D`$ be a square-free, positive rational integer and let $`K=(\sqrt{D})`$. Denote by $`ϵ_0`$ the fundamental unit of $`K`$ which exceeds unity, by $`s`$ the period of the continued fraction expansion for $`\sqrt{D}`$, and by $`P/Q`$ the ($`s1`$)-th approximant of it.
If $`D1mod4`$ or $`D1mod8`$, then
$$ϵ_0=P+Q\sqrt{D}.$$
However, if $`D5mod8`$, then
$$ϵ_0=P+Q\sqrt{D}.$$
or
$$ϵ_0^3=P+Q\sqrt{D}.$$
Finally, the norm of $`ϵ_0`$ is positive if the period $`s`$ is even and negative otherwise.
It is easy, working modulo $`4`$, to determine simple conditions (on $`t_0`$ ) which make $`F_j2\text{ or}\mathrm{\hspace{0.17em}\hspace{0.17em}3}mod4`$ and thus to say further, for a particular set of choices of $`t_1,\mathrm{},t_{\frac{n+1}{2}}`$ and for all odd or even $`t_0`$, that if $`F_j`$ is square-free, then $`a_0Q_n+P_n+\sqrt{F_j}Q_n`$ is the fundamental unit in $`[\sqrt{F_j}]`$. For example, suppose that $`n`$ is even and that the original $`Q_{n1}`$ determined from the permissible zero-one sequence is also even (so that $`P_{n1}`$ and $`Q_n`$ are both odd and $`P_n=Q_{n1}`$ is even). Then the multi-variable form of $`Q_{n1}`$ evaluated in *Step 3* will also have all even coefficients. Suppose $`\frac{Q_{n1}}{2}c_0+t_i^{^{}}mod2`$. (Here $`c_0`$ may be $`0`$ and the sum $`t_i^{^{}}`$ may contain some, all or none of the $`t_i`$’s ) It is easy to see that $`F_j(t_0,t_1,\mathrm{},t_{\frac{n+1}{2}})(c_0+t_i^{^{}}+t_0)^2+1mod4`$. Even more simply, if the original original $`Q_{n1}`$ as in *Step 1* is odd (here also the case $`n`$ is even is considered) then $`P_{n1}`$ as evaluated in *Step 3* is even and it is not difficult to show that in fact $`P_{n1}2mod4`$ (since for $`n`$ even $`P_nQ_{n1}P_{n1}Q_n=1`$) and that $`Q_n`$ is odd, which leads to $`F_j(t_0,t_1,\mathrm{},t_{\frac{n+1}{2}})t_0^2+1mod4`$. Similar relations hold in the case where $`n`$ is odd.
The polynomials constructed in theorem (2) take values in only one parity class, if all the variables are positive. However, given any two parity classes, there are multi-variable Fermat-Pell polynomials that take values in those two classes.
###### Theorem 4.
Let $`n`$ be any fixed positive integer large enough so that the set of positive integers whose square roots have a continued fraction expansion of period $`n+1`$ can be divided into more than one parity class.
(i) Given any two parity classes of integers whose square roots have continued fraction expansions of period $`n+1`$, there are multi-variable Fermat -Pell polynomials, which can be constructed, that take values in both parity classes;
(ii) These polynomials have a polynomial continued fraction expansion which can be explicitly determined;
(iii) If $`F=F(t_0,c,t_1,\mathrm{},t_{\frac{n+1}{2}})`$ is any such polynomial then the fundamental polynomial solution
$$C=C(t_0,c,t_1,\mathrm{},t_{\frac{n+1}{2}}),H=H(t_0,c,t_1,\mathrm{},t_{\frac{n+1}{2}})$$
to
(3.5)
$$C^2FH^2=(1)^{n1}$$
can be explicitly determined.
###### Proof.
As in *Step 2* in theorem (2) a polynomial sequence $`\{a_1,\mathrm{},a_n\}`$ is created. Suppose $`L_1=\{b_1\mathrm{},b_n\}`$ and $`L_2=\{c_1,\mathrm{},c_n\}`$ are the permissible sequences associated with the two parity classes. Let $`i_1,\mathrm{},i_k`$ be those positions $`\frac{n+1}{2}`$ at which the sequences agree. For each of these $`i_r`$’s set $`a_{i_r}(t_{i_r})=a_{n+1i_r}(t_{i_r})=2t_{i_r}+1`$, if $`c_{i_r}`$ is odd and set $`a_{i_r}(t_{i_r})=a_{n+1i_r}(t_{i_r})=2t_{i_r}+2`$, if $`c_{i_r}`$ is even. Subdivide the remaining positions (those positions $`\frac{n+1}{2}`$ at which $`L_1`$ and $`L_2`$ differ) into two subsets: those at which $`L_1`$ has a $`0`$ and $`L_2`$ has a $`1`$ and those at which $`L_1`$ has a $`1`$ and $`L_2`$ has a $`0`$.
Suppose $`i_j`$ is a position of the first kind. Let $`a_{i_j}(c,t_{i_j})=a_{n+1i_j}(c,t_{i_j})`$ $`=c+2+2t_{i_j}`$. Repeat this for all the positions $`i_j`$ in this first set. Likewise, Suppose $`i_j`$ is a position of the second kind. In this case let $`a_{i_j}(c,t_{i_j})=a_{n+1i_j}(c,t_{\frac{n+1}{2}})`$ $`=c+1+2t_{i_j}`$.This is also repeated for all the positions $`i_j`$ in this second set. *Step 3* and *Step 4* are then carried out as above. The rest of the proof is identical to theorem (2). Denote the polynomial produced by
(3.6)
$$F:=F(t_0,c,t_1,\mathrm{},t_{\frac{n+1}{2}}).$$
As in theorem (2), if $`c`$ and all the $`t_i`$’s are non-negative, $`1i\frac{n+1}{2}`$ and $`t_0>(1)^{n+1}Q_{n1}P_{n1}/(2Q_n)`$ then
$$\sqrt{F}=[a_0;\overline{a_1,\mathrm{}\mathrm{}.,a_n,2a_0}],$$
where the $`a_i`$’s, $`1in`$ are as defined just above and $`a_0`$ is as defined in equation (3.3).
Under these conditions also the parity class of $`F(t_0,c,t_1,\mathrm{},t_{\frac{n+1}{2}})`$ will depend only on the parity of $`c`$. As in theorem (2) the fundamental polynomial solution to
$$C^2F(t_0,c,t_1,\mathrm{},t_{\frac{n+1}{2}})H^2=(1)^{n1}$$
is given by $`C=a_0Q_n+P_n,H=Q_n`$. ∎
## 4. A Worked Example
As an example, consider those positive integers whose square-roots have continued fraction expansion with period of length 9. Thus the symmetric part of the period has length $`8`$ and it is necessary to check the $`2^4=16`$ zero-one sequences to determine which are permissible. (This checking is done in essentially the same way as in Example 1 above.) There are $`11`$ valid sequences:
$`0,0,0,0,0,0,0,0`$
$`0,0,0,1,1,0,0,0`$
$`0,0,1,1,1,1,0,0`$
$`0,1,0,0,0,0,1,0`$
$`0,1,0,1,1,0,1,0`$
$`0,1,1,1,1,1,1,0`$
$`1,0,0,1,1,0,0,1`$
$`1,0,1,0,0,1,0,1`$
$`1,0,1,1,1,1,0,1`$
$`1,1,0,0,0,0,1,1`$
$`1,1,1,0,0,1,1,1`$
The ninth of these is considered in more detail (Each of the others can be dealt with in a similar way). For clarity the letters $`a,b,c`$ and $`d`$ are used instead of $`t_1,t_2,t_3`$ and $`t_4`$. Evaluating the continued fraction
(4.1)
$$0+\frac{1}{2a+1+}\frac{1}{2b+2+}\frac{1}{2c+1+}\frac{1}{2d+1+}\frac{1}{2d+1+}\frac{1}{2c+1+}\frac{1}{2b+2+}\frac{1}{2a+1}$$
it is found that
$`P_8`$ $`=Q_7=12d+`$
$`2(3+4a+2b+4ab)(4+3b+4c+4bc+6d+4bd+8cd+8bcd)`$
$`+4(3+2b+4(1+b)c)(2+b+3c+2bc+a(3+2b+4(1+b)c))\times `$
$`(1+2d+2d^2),`$
$`P_7`$ $`=4(1+b)(4+3b+4c+4bc+6d+4bd+8cd+8bcd)+`$
$`2(3+2b+4(1+b)c)^2(1+2d+2d^2)\text{ and}`$
$`Q_8`$ $`=8(2+b+3c+2bc+a(3+2b+4(1+b)c))^2(1+2d+2d^2)+`$
$`(3+4a+2b+4ab)(3+4c+4d+8cd+(2+4a)(4+3b+`$
$`4c+4bc+6d+4bd+8cd+8bcd)).`$
Since $`n`$ is $`8`$ (even) and $`Q_8`$ is odd (so $`t_0`$ cannot take half-integer values), in this case $`F_9(t_0,a,b,c,d)`$ is defined by
(4.2)
$$F_9(t_0,a,b,c,d)=(Q_7P_7/2+t_0Q_8)^2+2t_0P_8+P_7^2$$
and
$$\begin{array}{c}\sqrt{F_9(t_0,a,b,c,d)}=[Q_7P_7/2+t_0Q_8;\overline{2a+1,2b+2,2c+1,2d+1,2d+1,}\hfill \\ \hfill \overline{2c+1,2b+2,2a+1,2(Q_7P_7/2+t_0Q_8)}],\end{array}$$
this expansion being valid for all $`a,b,c,d0`$ and all $`t_0>Q_7P_7/(2Q_8)`$ and in particular for all $`t_00`$. In these ranges
$$C=(Q_7P_7/2+t_0Q_8)Q_8+P_8,H=Q_8$$
gives the fundamental polynomial solution to
$$C^2F_9H^2=1.$$
$$F_9(t_0,a,b,c,d)=(Q_7P_7/2+t_0Q_8)^2+2t_0P_8+P_7^2(1+t_0^2)mod4$$
so that if $`(Q_7P_7/2+t_0Q_8)^2+2t_0P_8+P_7^2`$ is a square-free number for some particular $`a,b,c,d\mathrm{\hspace{0.17em}0}`$ and some odd $`t_0>Q_7P_7/(2Q_8)`$, then
$$(Q_7P_7/2+t_0Q_8)Q_8+P_8+\sqrt{(Q_7P_7/2+t_0Q_8)^2+2t_0P_8+P_7^2}Q_8$$
is the fundamental unit in $`\left(\sqrt{(Q_7P_7/2+t_0Q_8)^2+2t_0P_8+P_7^2}\right)`$.
## 5. Mystification, Fermat-Pell polynomials of a single variable and more on odd-even
Clearly it is possible to “mystify” this process by replacing each $`t_i`$ by some polynomial $`g_i(t_i)`$ taking only positive values or by replacing $`2t_i`$ (recalling that the continued fraction expansion contains only terms like $`2t_i+1`$ or $`2t_i+2`$) by some polynomial $`g_i(t_i)`$ taking only even non-negative values or by setting $`t_i=t_i(X_1,X_2,\mathrm{},X_k),\mathrm{\hspace{0.17em}\hspace{0.17em}1}i\frac{n+1}{2}`$, a polynomial in the $`X_j`$’s taking only positive values, where the $`X_j`$’s can be independent variables and $`k`$ can be as large as desired and so on.
Finally of course one can obtain single-variable Fermat-Pell polynomials by replacing the original variables $`t_0,t_i,1i\frac{n+1}{2}`$ by polynomials in a single variable. If it is desired that the period of the continued fraction expansion of the new single-variable Fermat-Pell polynomial should stay the same as that of the originating multi-variable polynomial then the domain of the single variable should be restricted so that the polynomials replacing each of the $`t_i`$’s take only positive values as in the multi-variable case and the polynomial replacing $`t_0`$ must be such that the $`a_0`$ term stays positive for all allowed values of the new single variable.
For example, letting $`a=s,b=0,c=s,d=0`$ and $`t_0=s`$ in the polynomial (4.2) above produces the single-variable Fermat-Pell polynomial
$$\begin{array}{c}g(s)=639557+6858268s+33078145s^2+\hfill \\ \hfill 94534688s^3+177380352s^4+228442240s^5+204593408s^6+\\ \hfill 125870080s^7+50925568s^8+12238848s^9+1327104s^{10}\end{array}$$
which has the continued fraction expansion (valid for all $`s0`$)
$$\begin{array}{c}\sqrt{g(s)}=[799+4289s+9184s^2+9856s^3+5312s^4+1152s^5;\hfill \\ \hfill \overline{2s+1,2,2s+1,1,1,2s+1,2,2s+1,}\\ \hfill \overline{2(799+4289s+9184s^2+9856s^3+5312s^4+1152s^5)}].\end{array}$$
$`g(s)(1+s^2)mod4`$ so when $`s`$ is odd and positive and $`g(s)`$ is square-free
$$\begin{array}{c}51982+534625s+2429840s^2+6408000s^3+\hfill \\ \hfill 10812928s^4+12115200s^5+9019392s^6+4304896s^7+1196032s^8+\\ \hfill 147456s^9+\sqrt{g(s)}(65+320s+576s^2+448s^3+128s^4)\end{array}$$
is the fundamental unit in $`[\sqrt{g(s)}]`$. For example, letting $`s=1`$ gives that $`47020351+1537\sqrt{935888258}`$ is the fundamental unit in $`[\sqrt{935888258}]`$.
Starting with the continued fraction
$$0+\frac{1}{2a+1+}\frac{1}{2b+2+}\frac{1}{c+2e+}\frac{1}{2d+1+}\frac{1}{2d+1+}\frac{1}{c+2e+}\frac{1}{2b+2+}\frac{1}{2a+1}$$
and following the same steps as above with the continued fraction (4.1) a multi-variable Fermat-Pell polynomial is developed which takes values in the parity classes associated with permissible sequences 7 and 9. Letting $`a=b=d=e=t=0`$ one gets the single-variable Fermat-Pell polynomial
$$\begin{array}{c}g(c)=4325+28140c+83652c^2+147440c^3+168000c^4+\hfill \\ \hfill 126528c^5+61504c^6+17664c^7+2304c^8\end{array}$$
with continued fraction expansion
$$\begin{array}{c}\sqrt{g(c)}=[65+214c+288c^2+184c^3+48c^4;\hfill \\ \hfill \overline{1,2,c,1,1,c,2,1,2(65+214c+288c^2+184c^3+48c^4)}],\end{array}$$
valid for $`c\mathrm{\hspace{0.17em}1}`$.
## 6. Concluding Remarks
Every Fermat-Pell polynomial in one variable, $`s`$ say, that eventually has a continued fraction expansion of fixed period length can be found from (3.2), if it takes values in only one parity class for all sufficiently large $`s`$, and from (3.6), if it takes values in two parity class for all sufficiently large $`s`$. (Recall remark (i) after theorem (1))
Of course none of this does anything to answer Schinzel’s question of whether every Fermat-Pell polynomial in one variable has a continued fraction expansion. Neither does it provide a criterion (such as Schinzel’s in the degree-two case) for deciding if a polynomial of arbitrarily high even degree is a Fermat-Pell polynomial. Perhaps it raises another question - Does every multi-variable Fermat-Pell polynomial have a continued fraction expansion? Does every multi-variable Fermat-Pell polynomial have a continued fraction expansion, assuming every Fermat-Pell polynomial in one variable does? |
warning/0001/hep-th0001014.html | ar5iv | text | # Contents
## 1. Introduction
Four fundamental interactions in nature are known today: Gravitation, electrodynamics, weak and strong nuclear forces. The latter three are in present day elementary particle physics successfully described by quantum gauge field theories. Successfully means in this context that there are no experimental data that do not agree with the predictions of these theories and that the agreement is very good e.g. in quantum electrodynamics (QED). The general theory of relativity describes gravitation classically. It is also a gauge theory in a wider sense of the word. A sound quantum theory for gravitation is still missing.
The distinguishing feature of these gauge theories is their gauge group: $`SU(2)\times U(1)`$ for the combined theory of electric and weak interactions and $`SU(3)`$ for the strong interaction. Both gauge groups are non-Abelian Lie groups. Therefore a comprehensive understanding of non-Abelian quantum gauge theories is needed to understand nature at the quantum level.
Originally the conceptual and mathematical framework of quantum field theory was developed for Abelian theories and in particular for QED. This was already an established theory in perfect agreement with the experimental data when physicists directed their attention towards non-Abelian gauge theories. They realized that quantum field theory required a modification of its mathematical description before it could be applied to non-Abelian theories.
The first study of a non-Abelian model — motivated by the isospin $`SU(2)`$ group — which attained wide reception was done by Yang and Mills \[YM54\] in 1954<sup>1</sup><sup>1</sup>1The first who studied non-Abelian models was O. Klein in 1938 \[Kle39\]. The interest of elementary particle physicists in non-Abelian quantum field theories grew strongly when in the next two decades several other such models were proposed to explain various phenomena. These include e.g. the Salam–Weinberg model \[Sal68, Wei67\] and the $`SU(3)`$ colour model for the strong interaction \[GM64, Zwe64\], but also attempts to quantize gravitation, e.g. \[Fey63\] or \[DeW67a, DeW67b, DeW67c\].
There was a series of obstacles to a satisfactory quantum theory for non-Abelian gauge theories due to the self coupling of the gauge bosons. A naive application of the methods developed for QED leads to serious difficulties, like an S-matrix that fails to be unitary \[Fey63, DeW67b\].
A major step to overcome these obstacles was made by Faddeev and Popov \[FP67, Fad69\]. They defined a unitary S-matrix in the functional integral approach, but for that they had to introduce unphysical fields that violate the spin-statistics theorem — the famous Faddeev Popov ghosts. In the mid seventies Becchi, Rouet and Stora \[BRS74, BRS76\] and independently of them Tyutin \[Tyu75\] found that the Faddeev Popov Lagrangian is invariant under a rigid symmetry transformation that mixes the ghosts with the other fields — the BRS transformation. Kugo and Ojima \[KO79\] gave an operator formulation<sup>2</sup><sup>2</sup>2Curci and Ferrari \[CF76\] gave already an operator formulation, but they postulated wrong hermiticity properties for the ghosts for this BRS theory, and Scharf and collaborators \[DHKS94a\] found with the operator gauge invariance a criterion of BRS symmetry for operator theories that needs no recurrence to an underlying classical theory<sup>3</sup><sup>3</sup>3Originally operator gauge invariance was postulated for theories of the Yang-Mills type. Recent results of Scharf and Wellmann \[SW99\] that it also a suitable criterion for spin two models.
Quantum field theory is plagued with two sources of possible infinities: the ultraviolet and the infrared divergences. Ultraviolet divergences are due to the distributional character of the field operators. In perturbation theory they can be removed by numerous renormalization procedures — so they are under control in this framework. Unfortunately these renormalization procedures are not unique — there remains the freedom of finite renormalization.
Infrared divergences occur since the asymptotic behaviour of incoming and outgoing interacting fields is not under control. This problem is particularly severe for non-Abelian gauge theories. It may be overcome by a replacement of the coupling constant by a spacetime dependent switching function so that the theory becomes free at finite times in the past and in the future. But as the real physical coupling is constant, one must in general perform the adiabatic limit, i.e. let the switching function tend to a constant. This limit does not exist in general.
In QED the infrared divergences are logarithmic, and Blanchard and Seneor \[BS75\] proved that the adiabatic limit exists for Green’s and Wightman functions. Unfortunately this is no longer true for non-Abelian theories. Their infrared behaviour is in general worse.
For strongly interacting fields this comes from the experimental observation of confinement. This means the fact that strongly interacting particles always combine to hadrons. Even after a high energy scattering process that breaks up the hadron structure the particles recombine immediately into new hadrons (hadronization). So the fields are not asymptotically free but constitute bound states. Moreover confinement cannot be described perturbatively.
In the electroweak theory confinement does not occur, but the model contains unstable, observable particles — the vector bosons $`W^\pm `$ and $`Z`$. These cannot occur as asymptotic states.
A solution for the infrared problem is to consider local theories, i.e. theories where all fields are localized in a finite region of spacetime. If the coupling is constant within this region and if the algebra of fields remains unaltered when the coupling is modified outside that region, the adiabatic limit needs not to be performed. Brunetti and Fredenhagen \[BF97\] proved that such a modification induces merely a unitary transformation on the algebra of fields. So the physical content is not changed by that modification and there is consequently no need for the adiabatic limit. Therefore no infrared problems occur in the construction of the local algebras.
One common problem of gauge theories — already encountered in QED — is that the algebra of fields must be quantized in an indefinite inner product space. Therefore positivity must be assured, i.e. the algebra of observables must be non trivially represented in a Hilbert space. Dütsch and Fredenhagen \[DF99\] succeeded in proving positivity for perturbation theories quantized in the BRS framework, provided the underlying free theory is also positive. In their view the interacting theory is regarded as a deformation of that underlying free model. They also constructed a local perturbation theory for QED.
The first to examine Yang-Mills theories in the causal framework were Scharf and collaborators \[DHKS94a\] \- \[DHS95b\], see also \[Sch95\]. They investigated the operator gauge invariance in the Yang-Mills case and found that it can be accomplished, provided a weak assumption concerning the infrared behaviour of the Green’s functions is fulfilled.
The aim of this thesis is to construct local perturbative gauge theories as operator theories in the BRS framework. The design is as general as possible, the motivation is Yang-Mills theory which serves as an example throughout the thesis.
The result is that the construction can always be performed, provided the generalized operator gauge invariance holds. It could not be proven that the latter can be accomplished in general models. A similar set of equations are the descent equations in the framework of algebraic renormalization — see, e.g. \[PS95\]. It may be possible to prove generalized operator gauge invariance by translating these results into our language, but this seems to be a tedious task and is not done here.
We use the renormalization scheme of causal perturbation theory as it was developed by Epstein and Glaser \[EG73\] following ideas proposed by Bogoliubov, Shirkov \[BS59\] and Stückelberg. It avoids divergent expressions throughout the entire procedure. Scharf and collaborators as well as Dütsch and Fredenhagen formulated their results in the same framework. This makes it easy to use their results for our construction and to compare them with our results.
Moreover our approach is local in order to avoid infrared divergences and to be able to define observables and physical states.
Like Dütsch and Fredenhagen we use normalization conditions for the time ordered products as an essential tool to establish desired relations in the field algebra. Their normalization conditions are generalized to include fields that contain a spacetime derivative. Ward identities for the ghost and BRS current are introduced as new normalization conditions with regard to the definition of observables and physical states. We introduce an algebra of auxiliary variables for the fields containing a spacetime derivative and define a linear representation of the polynomials in this algebra as operators acting on the Fock space. We present a reformulation of time ordering. It is formally a multi linear generalization of the linear representation mentioned above to multiple arguments. This and the definition of propagator functions for the fields with a spacetime derivative allows us to generalize the normalization conditions in the desired manner. It is proven that all these conditions — except the BRS Ward identities — can be accomplished simultaneously. The existence of a solution for the BRS Ward identities and its compatibility with the other conditions must be proven in individual models. The proof for QED is presented.
There are relations for the local field algebra that are determined by the normalization conditions, e.g. renormalized field equations and the BRS algebra. The latter allows for a definition of observables and a construction of a positive physical state space.
The thesis is organized as follows: In chapter (2) we set up the algebraic framework of BRS theory, following Kugo and Ojima \[KO79\]. The definition of observables and the construction of the Hilbert space are performed using certain algebraic relations between the interacting operators. The rest of the thesis will be devoted to the construction of models in which these relations hold.
In chapter (3) the free model underlying our perturbation theory is put up. The algebra of auxiliary variables is constructed and its linear representation as Fock space operators is defined. Then the propagator functions are examined. Finally the proof of Razumov and Rybkin \[RR90\] for the positivity of theories with certain BRS charges is presented.
The new definition of time ordering is given in chapter (4). It contains also the formulation of six normalization conditions and the proof that the first five have simultaneous solutions. The sixth, the BRS Ward identities, is shown to be equivalent with a generalized version of operator gauge invariance.
Local causal perturbation theory is introduced in chapter (5) along the lines of Epstein, Glaser \[EG73\], Dütsch and Fredenhagen \[DF99\]. Conditions for a polynomial to be a candidate for a Lagrangian are given.
The local field algebra is constructed in chapter (6). The conserved currents and charges, the ghost number of an interacting field and the interacting BRS transformation are defined, field equations and the BRS algebra are derived. The chapter concludes with a reflection on the correspondence between the quantum theory defined above and its classical counterpart.
The inspection of gauge theories is deepened in chapter (7) for two exemplary models: QED and Yang-Mills theory. The BRS Ward identities are proven for QED, and we compare the relations between the interacting fields with those between the corresponding classical fields.
At the end a conclusion and an outlook for possible further developments are included.
## 2. BRS theory — algebraic considerations
In this chapter canonical BRS theory according to Kugo and Ojima \[KO79\] is carried through on a purely algebraic level. The availability of suitable BRS and ghost charges is formulated as assumptions. Then perturbative theories — i.e. theories where the operators and the state vectors are formal power series — are examined in this framework. Dütsch and Fredenhagen \[DF99, DF98\] prove that the positivity structure of a theory can be maintained during deformation. Their proof is presented here.
### 2.1. Why BRS theory?
All quantum gauge theories share one common difficulty: There is no positive definite Hilbert space in which the field algebra can be represented and which possesses a nontrivial unitary representation of the Poincaré group. Nakanishi and Ojima \[NO90\] proved that there exists no nontrivial Hilbert space representation for manifestly covariant theories with massless gauge bosons. This could be circumvented by non covariant gauges, but this means abandoning manifest covariance.
The field algebra is not observable, so a direct physical interpretation of the theory which requires a Hilbert space representation is not possible. But the algebra of observables must have a Hilbert space representation, and the Hilbert space must carry a unitary representation of the Poincaré group.
For QED Gupta \[Gup50\] and Bleuler \[Ble50\] found an elegant way out of this dilemma. They retain manifest covariance at the prize of representing the field algebra in an indefinite inner product space. Then there exists a non trivial, pseudo unitary<sup>4</sup><sup>4</sup>4Pseudo unitary, pseudo hermitian etc. means unitary, hermitian etc. w.r.t. the indefinite inner product. representation of the Poincaré group. This space is too big: It contains vectors with negative norm that have no physical interpretation — they would lead to negative transition probabilities. Consequently the physical state vectors form a distinguished proper subspace of the inner product space. This subspace is selected by a linear subsidiary condition, and it is found to be positive semidefinite. It becomes a Hilbert space with unitary action of the Poincaré group when all state vectors differing by a zero norm vector are identified with each other and the space is subsequently completed.
Unfortunately this strategy breaks down in non Abelian gauge theories because there is no appropriate subsidiary condition available. This is due to the nonlinear self interaction of the gauge fields.
BRS theory is a solution for that problem. The canonical BRS formalism of Kugo and Ojima \[KO79\] follows the same ideas as Gupta and Bleuler but it can also be applied to non-Abelian theories. Initially the algebra of fields is again represented in an indefinite inner product space. The presence of the ghosts in the BRS approach makes it possible to define a suitable subsidiary condition for the physical subspace which is a Hilbert space. The formalism provides also a definition of an algebra of observables that is represented in this Hilbert space. There exists a pseudo unitary action of the Poincaré group on the indefinite space. This action is lifted to a unitary one on the Hilbert space.
### 2.2. Canonical BRS theory
The construction starts in the following situation: There is an initial Hilbert space $`\{𝒱,(,)\}`$ with a positive scalar product $`(,)`$ that encompasses all fields including the unphysical ones (scalar vector bosons, ghosts etc.). This scalar product has no direct physical meaning. It does not describe the transition amplitudes, in particular it is not Poincaré covariant. The adjoint in this Hilbert space is denoted as <sup>+</sup>, i.e. $`(\varphi ,A\psi )=(A^+\varphi ,\psi )`$ for every<sup>5</sup><sup>5</sup>5$`(𝒱)`$ is the space of endomorphisms on $`𝒱`$ $`A𝒱`$. It is possible to find a Krein operator $`J(𝒱)`$ in the Hilbert space with the following three properties:
* $`J`$ is hermitian, i.e. $`J^+=J`$
* It is idempotent, i.e. $`J^2=1\mathrm{l}`$
* It defines a new inner product on $`𝒱`$ via
$$\varphi ,\psi \stackrel{\text{def}}{=}(\varphi ,J\psi )$$
(2.1)
such that the new inner product is Poincaré covariant.
The new inner product is assumed to describe the correct transition probabilities. Therefore it is referred to as the physical inner product. The vector space $`𝒱`$ forms a Krein space with the physical product $`,`$. Since $`(,)`$ was not covariant while $`,`$ was, $`J=1\mathrm{l}`$ can be excluded. Then the physical inner product is always indefinite, because there must exist a vector $`|\varphi `$ such that $`(1\mathrm{l}J)|\varphi 0`$, and then $`(1\mathrm{l}J)|\varphi `$ has negative norm. The adjoint w.r.t. the physical inner product is defined as an involution denoted by , namely $`A^{}\stackrel{\text{def}}{=}JA^+J`$, such that $`\varphi ,A\psi =A^{}\varphi ,\psi `$ for every $`A𝒱`$.
For the canonical BRS theory the following assumption is essential:
A1: There exists an operator $`Q_B(𝒱)`$ — the BRS charge — with the following properties:
* $`Q_B`$ is a conserved charge.
* It is pseudo hermitian, i.e. $`(Q_B)^{}=Q_B`$.
* It is nilpotent<sup>6</sup><sup>6</sup>6Nilpotent means throughout this thesis nilpotent of order two., i.e. $`(Q_B)^2=0`$.
* It annihilates the vacuum, i.e. $`Q_B|\omega =0`$
where $`|\omega `$ is the vacuum vector. This assumption is highly non trivial, and the appearance of ghosts in $`𝒱`$ is necessary for it. It has to be verified in the concrete model.
It is easily verified that the image of $`Q_B`$ contains only zero norm vectors w.r.t. the physical scalar product:
$$Q_B\varphi ,Q_B\varphi =\varphi ,(Q_B)^2\varphi =0.$$
(2.2)
With the second assumption a grading is introduced on $`𝒱`$ by means of the ghost charge $`Q_c`$.
A2: There exists an operator $`Q_c(𝒱)`$ — the ghost charge — with the following properties:
* $`Q_c`$ is a conserved charge.
* It is anti pseudo hermitian, i.e. $`(Q_c)^{}=Q_c`$.
* It has integer eigenvalues, i.e. $`Q_c|\psi =q|\psi q\mathrm{𝖹𝖹}`$.
* It satisfies the commutator relation $`[Q_c,Q_B]_{}=Q_B`$.
* It annihilates the vacuum, i.e. $`Q_c|\omega =0`$ .
The eigenvalue of a state vector w.r.t. the ghost charge is called its ghost number. For the physical inner product of two vectors to be non zero they must have opposite ghost numbers: Let $`Q_c\psi =q\psi `$ and $`Q_c\varphi =p\varphi `$, then
$$0=\psi ,Q_c\varphi \psi ,Q_c\varphi =\psi ,Q_c\varphi +Q_c\psi ,\varphi =(q+p)\psi ,\varphi ,$$
(2.3)
so $`(q+p)=0`$ or $`\psi ,\varphi =0`$. This implies in particular that only states with vanishing ghost number can have non zero norm w.r.t. the physical inner product.
The commutator relation $`[Q_c,Q_B]_{}=Q_B`$ forms together with the nilpotency of the BRS charge, $`(Q_B)^2=0`$, the BRS algebra.
Like in the Gupta-Bleuler scheme the negative norm states are excluded by a subsidiary condition. The kernel of $`Q_B`$ is regarded as a candidate for the physical Hilbert space. It contains necessarily zero norm states from the image of $`Q_B`$ — due to $`(Q_B)^2=0`$ we have $`Q_B\mathrm{ker}Q_B`$ — and possibly also vectors with non vanishing ghost number. Therefore the following definition for the Hilbert space $`_{\mathrm{ph}}`$ of physical state vectors is given:
$$_{\mathrm{ph}}\stackrel{\text{def}}{=}\overline{(\mathrm{ker}Q_B,𝒱)/(Q_B,𝒱)}^{}.$$
(2.4)
Completion is understood in the norm topology. Now it must be verified that the physical state vectors form a positive definite inner product space. This is guaranteed if the following positivity assumption is valid.
A3:
* The kernel of $`Q_B`$ contains only positive semidefinite vectors, i.e. $`Q_B|\varphi =0\varphi ,\varphi 0`$
* Its image encompasses all zero norm vectors in its kernel, i.e. $`|\varphi \mathrm{ker}(Q_B,𝒱)`$ and $`\varphi ,\varphi =0|\varphi (Q_B,𝒱)`$.
The second point guarantees in particular that all elements in $`(\mathrm{ker}Q_B,𝒱)`$ with nonvanishing ghost number are in $`(Q_B,𝒱)`$. The scalar product is well defined on these equivalence classes, so it does not depend on the representative of a class:
$$\varphi +Q_B\chi ,\psi =\varphi ,\psi +\chi ,Q_B\psi =\varphi ,\psi .$$
(2.5)
It is also positive definite by construction — if assumption A3 holds —, so the quotient space is a pre Hilbert space and becomes a Hilbert space after completion. The structure above is called a state cohomology.
The ghost charge induces a derivation on $`(𝒱)`$,
$$s_c(A)\stackrel{\text{def}}{=}[Q_c,A]_{}A(𝒱).$$
(2.6)
Its eigenvalue for an operator $`A(𝒱)`$ is called the ghost number of $`A`$ and is always an integer.
The BRS charge induces an graded derivation on $`(𝒱)`$, namely the BRS transformation<sup>7</sup><sup>7</sup>7Here $`[,]_{}`$ denotes the graded commutator. Suppose, $`A,B𝒱`$ have ghost numbers $`a,b\mathrm{𝖹𝖹}`$. Then $`[A,B]_{}\stackrel{\text{def}}{=}AB(1)^{ab}BA`$.
$$s(A)\stackrel{\text{def}}{=}[Q_B,A]_{}A(𝒱).$$
(2.7)
It is nilpotent because $`Q_B`$ is also nilpotent and the Jacobi-identity holds for the graded commutators.
With these definitions the algebra of observables $`𝒜_{\mathrm{ph}}`$ can be defined as
$$𝒜_{\mathrm{ph}}\stackrel{\text{def}}{=}\left((\mathrm{ker}s,(𝒱))(\mathrm{ker}s_c,(𝒱))\right)/\left((s,(𝒱))(\mathrm{ker}s_c,(𝒱))\right).$$
(2.8)
This structure is called an operator cohomology. Its elements are well defined operators on $`_{\mathrm{ph}}`$, i.e. $`𝒜_{\mathrm{ph}}(\mathrm{ker}Q_B,𝒱)(\mathrm{ker}Q_B,𝒱)`$ and $`𝒜_{\mathrm{ph}}[0]=[0]`$, where $`[0]`$ is the equivalence class of zero.
There is a -involution induced on the algebra of observables by the -involution on the representatives. But unlike the original involution this one acts on operators on a Hilbert space, so the notions hermitian, unitary and so on must be used without the prefix pseudo.
There is also a unitary action of the Poincaré group defined on $`_{\mathrm{ph}}`$, namely the lift of the initial pseudo-unitary action on the representatives to the equivalence classes. This induces a unitary representation on $`_{\mathrm{ph}}`$.
There is a physical interpretation available for the cohomologies. Initially the model is not characterized in terms of the algebra of field operators described here but in terms of the sub algebra without the ghosts — these were only introduced to make possible the definition of the BRS charge. In the picture above physics is invariant under local gauge transformations, i.e. gauge transformations generated by spacetime dependent functions. Then the BRS transformation, restricted to the sub algebra, may be regarded as the infinitesimal local gauge transformation. The role of the spacetime dependent functions is played by the ghosts. For them the BRS transformation is defined such that it is nilpotent on the entire algebra.
So the restriction to the kernel of $`s`$ singles out fields that are invariant under infinitesimal gauge transformations. Fields in the same equivalence class are regarded as physically indistinguishable. In this interpretation the physical Hilbert space contains equivalence classes of states that are invariant under infinitesimal gauge transformations.
### 2.3. Interacting theories and deformation stability
In perturbation theory field operators are represented by formal power series of linear operators. This makes it necessary to recapitulate the BRS formalism for formal power series of state vectors and operators, since e.g. the notion of positivity is not defined a priori for formal power series. This situation has been examined by Dütsch and Fredenhagen \[DF98, DF99\] and we present their results here.
In their picture the interacting theory is a deformation of an underlying free theory. In some models positivity — i.e. assumption A3 — can be proven by direct computation for the underlying free theory. Dütsch and Fredenhagen found a construction for the deformed — i.e. interacting — state space such that positivity holds also there in a sense defined below.
In the interacting theory both the state space and the operators acting on it are modules over the ring $`\stackrel{~}{\text{ }\mathrm{C}}`$ of formal power series of complex numbers:
$$\stackrel{~}{\text{ }\mathrm{C}}\stackrel{\text{def}}{=}\{\stackrel{~}{a}=\underset{n=0}{\overset{\mathrm{}}{}}g^na_n:a_n\text{ }\mathrm{C}\}$$
(2.9)
where $`g`$ is the deformation parameter. The element $`\stackrel{~}{1\mathrm{l}}\stackrel{\text{def}}{=}(1,0,0,\mathrm{})`$ is the identity in this ring. An element $`\stackrel{~}{a}\stackrel{~}{\text{ }\mathrm{C}}`$ is only invertible<sup>8</sup><sup>8</sup>8Bordemann and Waldmann \[BW96\] consider Laurent series instead. These are invertible if $`\stackrel{~}{a}0`$, so they form a field. if $`a_00`$. The interacting indefinite inner product space is defined as the $`\stackrel{~}{\text{ }\mathrm{C}}`$-module $`\stackrel{~}{𝒱}\stackrel{\text{def}}{=}\{\stackrel{~}{\psi }=_ng^n\psi _n:\psi _n𝒱\}`$ which has the inner product $`,`$ induced from $`𝒱`$. For $`\stackrel{~}{\psi }=_ng^n\psi _n`$ and $`\stackrel{~}{\chi }=_ng^n\chi _n`$ this means
$$\begin{array}{cc}& ,:\stackrel{~}{𝒱}\times \stackrel{~}{𝒱}\stackrel{~}{\text{ }\mathrm{C}}\hfill \\ & \stackrel{~}{\psi },\stackrel{~}{\chi }=\underset{n}{}g^n\left(\underset{k=1}{\overset{n}{}}\psi _k,\chi _{nk}\right).\hfill \end{array}$$
(2.10)
This is sesquilinear in $`\stackrel{~}{\text{ }\mathrm{C}}`$, i.e. $`\stackrel{~}{a}\stackrel{~}{\chi },\stackrel{~}{b}\stackrel{~}{\psi }=\stackrel{~}{a}^{}\stackrel{~}{b}\stackrel{~}{\chi },\stackrel{~}{\psi }`$. The means complex conjugation, where the deformation parameter $`g`$ is real, so
$$\stackrel{~}{a}^{}=\underset{n=0}{\overset{\mathrm{}}{}}g^n\overline{a}_n$$
(2.11)
where $`\overline{}`$ denotes complex conjugation in $`\mathrm{C}`$.
The operators in $`(\stackrel{~}{𝒱})`$ acting on $`\stackrel{~}{𝒱}`$ can be written as
$$(\stackrel{~}{𝒱})=\{\stackrel{~}{A}=\underset{n}{}g^nA_n:A_n(𝒱)\}$$
(2.12)
and form a $`\stackrel{~}{\text{ }\mathrm{C}}`$-module, too. The multiplication law in this algebra is
$$\stackrel{~}{A}\stackrel{~}{B}=\underset{n}{}g^n\left(\underset{k=1}{\overset{n}{}}A_kB_{nk}\right)\stackrel{~}{A},\stackrel{~}{B}(\stackrel{~}{𝒱}).$$
(2.13)
The interacting BRS-charge and the interacting ghost charge are such operators,
$$\begin{array}{cc}& \stackrel{~}{Q}_B=\underset{n}{}g^nQ_{B,n},Q_{B,n}(𝒱)\hfill \\ \hfill \text{and}& \stackrel{~}{Q}_c=\underset{n}{}g^nQ_{c,n},Q_{c,n}(𝒱),\hfill \end{array}$$
(2.14)
where $`\stackrel{~}{Q}_{B,0}`$ and $`\stackrel{~}{Q}_{c,0}`$ agree with the free charges. $`\stackrel{~}{Q}_B`$ must be chosen such that it is nilpotent, $`\stackrel{~}{Q}_B^2=0`$, and pseudo hermitian, $`(\stackrel{~}{Q}_B)^{}=\stackrel{~}{Q}_B`$, and $`\stackrel{~}{Q}_c`$ must be anti pseudo hermitian, $`(\stackrel{~}{Q}_c)^{}=\stackrel{~}{Q}_c`$. The involution is the one induced from $`(𝒱)`$. The charges must satisfy the BRS algebra $`[\stackrel{~}{Q}_c,\stackrel{~}{Q}_B]_{}=\stackrel{~}{Q}_B`$.
The interacting state space can be defined as in the general case,
$$\stackrel{~}{}_{\mathrm{ph}}\stackrel{\text{def}}{=}(\mathrm{ker}\stackrel{~}{Q}_B,\stackrel{~}{𝒱})/(\stackrel{~}{Q}_B,\stackrel{~}{𝒱}),$$
(2.15)
with the only difference that the space is not completed since there is no convenient topology in the space of formal power series.
The question is whether this space has a positive scaler product, and above all what positivity means for formal power series.
Following Dütsch and Fredenhagen \[DF99\] we adopt here Steinmann’s \[Ste89\] point of view<sup>9</sup><sup>9</sup>9Here Bordemann and Waldmann \[BW96\] follow again a different prescription: They define a real formal power series as positive if its first non vanishing coefficient is a positive number. With this definition the field of real Laurant series becomes ordered. The notion of positivity presented here is a stricter one: Every positive series in Steinmann’s sense is also positive in their sense, but not converse. that a formal power series $`\stackrel{~}{b}=_nb_ng^n\stackrel{~}{\text{ }\mathrm{C}}`$ is positive if it is the absolute square of another power series $`\stackrel{~}{c}\stackrel{~}{\text{ }\mathrm{C}}`$, i.e. $`\stackrel{~}{b}=\stackrel{~}{c}^{}\stackrel{~}{c}`$. Dütsch and Fredenhagen define also that a class of state vectors $`[\stackrel{~}{\phi }]\stackrel{~}{}_{\mathrm{ph}}`$ can be normalized if there exists an $`\stackrel{~}{a}\stackrel{~}{\text{ }\mathrm{C}}`$ and $`[\stackrel{~}{\psi }]\stackrel{~}{}_{\mathrm{ph}}`$ such that $`[\stackrel{~}{\phi }]=\stackrel{~}{a}[\stackrel{~}{\psi }]`$ and $`[\stackrel{~}{\psi }],[\stackrel{~}{\psi }]=\stackrel{~}{1\mathrm{l}}`$.
With these notions of positivity and normalizability they prove in \[DF99\] the following results:
Let the positivity assumption A3 be fulfilled for the undeformed theory. Then
$`(i)`$ $`\stackrel{~}{\psi },\stackrel{~}{\psi }0\stackrel{~}{\psi }(\mathrm{ker}\stackrel{~}{Q}_B,\stackrel{~}{𝒱})`$ (2.16)
$`(ii)`$ $`\stackrel{~}{\psi }(\mathrm{ker}\stackrel{~}{Q}_B,\stackrel{~}{𝒱})\stackrel{~}{\psi },\stackrel{~}{\psi }=0\stackrel{~}{\psi }(\stackrel{~}{Q}_B,\stackrel{~}{𝒱}),`$ (2.17)
$`(iii)`$ $`\psi (\mathrm{ker}Q_B,𝒱)\stackrel{~}{\psi }(\mathrm{ker}\stackrel{~}{Q}_B,\stackrel{~}{𝒱}):(\stackrel{~}{\psi })_0=\psi `$ (2.18)
$`(iv)`$ $`\text{Every }[\stackrel{~}{\psi }]0\stackrel{~}{}_{\mathrm{ph}}\text{ is normalizable in the sense above}.`$ (2.19)
For the proofs of these results we refer to their article.
So assumption A3 is fulfilled for the interacting theory if it is fulfilled for the free theory underlying it. Therefore the interacting physical state space defined above is a pre Hilbert space. Result $`(iii)`$ implies that an interacting vacuum state $`|\stackrel{~}{\omega }`$ can be defined that is annihilated by $`\stackrel{~}{Q}_B`$ such that $`|\stackrel{~}{\omega }_0=|\omega `$, provided that the free charge annihilates the free vacuum.
The interacting BRS-transformation is the formal power series
$$\stackrel{~}{s}=\underset{n}{}g^ns_n,\stackrel{~}{s}(\stackrel{~}{A})\stackrel{\text{def}}{=}[\stackrel{~}{Q}_B,\stackrel{~}{A}]_{}\stackrel{~}{A}\stackrel{~}{𝒱}.$$
(2.20)
Each $`s_n`$ is an anti-derivation on $`\stackrel{~}{𝒱}`$ and $`s_0`$ agrees with the free BRS-transformation. $`\stackrel{~}{s}_c`$ is analogously defined as
$$\stackrel{~}{s}_c=\underset{n}{}g^ns_{c,n},\stackrel{~}{s}_c(\stackrel{~}{A})\stackrel{\text{def}}{=}[\stackrel{~}{Q}_c,\stackrel{~}{A}]_{}\stackrel{~}{A}\stackrel{~}{𝒱}$$
(2.21)
where each $`s_{c,n}`$ is a derivation on $`\stackrel{~}{𝒱}`$ and $`s_{c,0}`$ agrees with $`s_c`$. The interacting observable algebra is defined as
$$\stackrel{~}{𝒜}_{\mathrm{ph}}\stackrel{\text{def}}{=}\left((\mathrm{ker}\stackrel{~}{s},\stackrel{~}{𝒱})(\mathrm{ker}\stackrel{~}{s}_c,\stackrel{~}{𝒱})\right)/\left((\stackrel{~}{s},\stackrel{~}{𝒱})(\mathrm{ker}\stackrel{~}{s}_c,\stackrel{~}{𝒱})\right).$$
(2.22)
So in the framework of BRS theory an algebra of interacting observables can be defined and represented in a (pre) Hilbert space if the following conditions can be accomplished:
1. In the underlying free theory a ghost charge $`Q_c`$ and a BRS charge $`Q_B`$ can be defined that fulfill the assumptions $`\mathrm{𝐀𝟏}`$ \- $`\mathrm{𝐀𝟑}`$.
2. A conserved interacting BRS charge $`\stackrel{~}{Q}_B`$ can be constructed such that $`(\stackrel{~}{Q}_B)_0=Q_B`$ with the properties $`\stackrel{~}{Q}_B^2=0`$ and $`(\stackrel{~}{Q}_B)^{}=\stackrel{~}{Q}_B`$.
3. A conserved interacting ghost charge $`\stackrel{~}{Q}_c`$ with integer eigenvalues can be constructed such that $`(\stackrel{~}{Q}_c)_0=Q_c`$ with the property $`(\stackrel{~}{Q}_c)^{}=\stackrel{~}{Q}_c`$ .
4. The BRS algebra $`[\stackrel{~}{Q}_c,\stackrel{~}{Q}_B]_{}=\stackrel{~}{Q}_B`$ holds.
## 3. The free theory
We start our considerations concerning BRS theory with free theories. The treatment of free theories in the BRS framework is not a goal in its own but provides us with definitions that will become important for the interacting theory in the next chapters. Furthermore positivity is proven for the underlying free model in order to take advantage of deformation stability for the interacting theory.
We already pointed out the essential significance of normalization conditions for the time ordered products in our construction. For some of these normalization conditions it is necessary to give a precise meaning to expressions like $`\frac{A}{\phi _j}(x)`$, the derivative of a Wick monomial $`A`$ w.r.t. a free field operator $`\phi _j`$. Dütsch and Fredenhagen \[DF99\] solve this problem for QED by an implicit definition,
$$[A(x),\phi _j(y)]_{}=i\underset{k}{}\mathrm{\Delta }_{jk}(xy)\frac{A}{\phi _k}(x)$$
(3.1)
where $`\mathrm{\Delta }_{jk}(x)`$ is a commutator function. This equation is indeed a definition for the partial derivative on the right hand side if the theory contains no derivated fields<sup>10</sup><sup>10</sup>10Derivated fields means here and below fields containing a spacetime derivative. like QED. But for theories that do contain such derivated fields — like Yang-Mills theory — there is no such definition available.
The natural attempt to include derivated fields would be the replacement of the partial derivative by a functional derivative, where the latter would be defined by means of
$$[A(x),\phi _j(y)]_{}=i\underset{k}{}d^4z\mathrm{\Delta }_{jk}(zy)\frac{\delta A(x)}{\delta \phi _k(z)}.$$
(3.2)
Unfortunately the equation above is no definition. This can be seen as follows: Let $`D`$ be the differential operator that implements the field equations for the free fields such that
$$\underset{j}{}D_{ij}^x\mathrm{\Delta }_{jk}(x)=0.$$
(3.3)
Such an operator exists in general, it will be explicitely constructed later in this chapter. We can define an operator $`D^{}`$ according to
$$d^4zf(z)\left(D_{ij}^{,z}g(z)\right)=d^4z\left(D_{ij}^zf(z)\right)g(z).$$
(3.4)
Then an expression of the form $`_mD_{mj}^{,z}\mathrm{\Phi }_m(x,z)`$ with arbitrary operators $`\mathrm{\Phi }_m(x,z)`$ can be added to the functional derivative without altering the equation.
Our strategy to solve this problem is the following: We introduce an algebra that is generated by symbols for the basic and derivated fields. These symbols serve as auxiliary variables. For this algebra the derivative w.r.t. a generator is defined. The polynomials in this algebra are then linearly represented as operator valued distributions acting on the Fock space. Time ordering is defined as a multi linear representation of several such polynomials as distributional Fock space operators in the next chapter. The derivative that we needed to formulate the normalization condiditons occurs only in the arguments of time ordered products. With the definition of time ordered products introduced above these arguments are polynomials in the algebra. For the algebra the derivative is well defined, and therefore the normalization conditions can be formulated.
The chapter is organized as follows: In the first section we define the algebra $`𝒫`$ of auxiliary variables. In the second section the Fock space $``$ and operators therein are constructed. This construction will be completely standard and is included here to establish our notation.
In the third section the linear representation of the algebra $`𝒫`$ as (distributional) Fock space operators is defined. This definition includes commutator functions for basic and derivated fields.
In the fourth section propagator functions for basic and derivated fields are constructed that have a differential operator as their inverse.
The chapter concludes with a section concerning the free model underlying Yang-Mills theory where the essential operators — ghost charge, BRS charge etc. — are defined. In particular we present Razumovs and Rybkins \[RR90\] proof for the positivity of that theory.
### 3.1. The algebra of auxiliary variables
The algebra $`𝒫`$ is the graded commutative $`\mathrm{C}`$-algebra generated by auxiliary variables for the basic and derivated fields. At first we specify its generators. Therefore we determine which basic fields and which derivatives of the basic fields we wish to deal with in the model to be defined. For example, with respect to Yang-Mills theory we include Lie algebra valued vector bosons $`A_\mu `$ and its first derivatives $`(_\nu A_\mu )`$, since in the interaction Yang-Mills Lagrangian both non derivated and derivated vector bosons appear. For the same reason ghosts and anti-ghosts $`u,\stackrel{~}{u}`$ and their derivatives $`(_\nu u),(_\nu \stackrel{~}{u})`$ are added. If we whish to include fermionic matter, coloured spinor fields $`\psi ,\overline{\psi }`$ must be incorporated, but no derivated spinors because these do not appear in the interaction Lagrangian. The set of fields is then completed by the double derivated vector bosons $`(_\nu _\rho A_\mu )`$ which do not appear in the Lagrangian but in the BRS current (see below). The non derivated fields are referred to as basic fields.
Then we define one symbol for each of these fields — with a distinct symbol for each derivative of the basic fields that is included in the list above. These symbols are the generators of $`𝒫`$. The generators corresponding to the basic fields are called the basic generators, those corresponding to derivated fields are called the higher generators. We adopt the following notation: the generators are written as $`\phi _i`$ where the index $`i`$ numerates the basic and higher generators. Sometimes it is desirable to distinguish basic and higher generators. Then the generators are denoted as $`\phi _i^\alpha `$, where the index $`i`$ numerates here the basic generators, and $`\alpha `$ is a multi index,
$$\alpha =(\left|\alpha \right|,\mu _1,\mathrm{},\mu _{\left|\alpha \right|}).$$
(3.5)
The degree $`\left|\alpha \right|`$ of a generator $`\phi _i^\alpha `$ is is the number of spacetime derivatives on the corresponding field operator. Basic generators are therefore denoted as $`\phi _i^{(0)}`$. The indices $`\mu _1,\mathrm{},\mu _{\left|\alpha \right|}`$ are Lorentz indices corresponding to the Lorentz indices of the spacetime derivatives on the field operator. The symbols $`\phi _i`$ may carry additional Lorentz (e.g. if $`\phi _i=A_\mu `$) or spinor (e.g. if $`\phi _i=\psi `$) indices. We will define a representation of the Lorentz group on $`𝒫`$ at the end of this section. To give an example for the multi indices, we relate some generators $`\phi _i`$ to the corresponding field operators:
$$\begin{array}{cc}\hfill \phi _i^{(0)}\phi _i(x),& \phi _i^{(1,\mu )}_x^\mu \phi _i(x)\hfill \\ & \phi _i^{(2,\mu \nu )}_x^\mu _x^\nu \phi _i(x)\mathrm{}.\hfill \end{array}$$
(3.6)
The symbols are symmetric under permutation of the Lorentz-indices stemming from the multi-indices, e.g. $`\phi _i^{(2,\mu \nu )}=\phi _i^{(2,\nu \mu )}`$.
The set $`𝒢`$ of all generators of $`𝒫`$ is defined as
$$𝒢\stackrel{\text{def}}{=}\{\phi _i^\alpha :\phi _i^\alpha \text{ has a counterpart in the desired set of fields}\}.$$
(3.7)
Sometimes the set of basic generators will become important:
$$𝒢_b\stackrel{\text{def}}{=}\{\phi _i^\alpha 𝒢:\alpha =(0)\}𝒢.$$
(3.8)
Now $`𝒫`$ is defined as the unital<sup>11</sup><sup>11</sup>11This means that there is an identity operator $`1\mathrm{l}`$ included in $`𝒫`$ algebra generated by $`𝒢`$. In addition, $`𝒫`$ is graded symmetric. There are two gradings involved here: the ghost number $`g`$ and the (physical) fermion number $`f`$,
$$f,g:\left\{\text{monomials in }𝒫\right\}\mathrm{𝖹𝖹}.$$
(3.9)
They are additive quantum numbers,
$$g(AB)=g(A)+g(B)\text{ and }f(AB)=f(A)+f(B)A,B𝒫,$$
(3.10)
and are defined as
$$\begin{array}{cc}\hfill g(u^\alpha )=1,g(\stackrel{~}{u}^\alpha )=1,g(\phi _i^\alpha )=0\text{otherwise}& \\ \hfill f(\psi ^\alpha )=1,f(\overline{\psi }^\alpha )=1,f(\phi _i^\alpha )=0\text{otherwise}& .\hfill \end{array}$$
(3.11)
Polynomials in $`𝒫`$ that have a definite ghost or fermion number are called homogeneous w.r.t. the ghost or fermion number.
Graded symmetric means that for any two elements of the algebra $`A,B𝒫`$ the following commutation relation holds:
$$AB=(1)^{g(A)g(B)+f(A)f(B)}BAA,B𝒫.$$
(3.12)
This means that $`𝒫`$ is the unital algebra freely generated by $`𝒢`$ with the equivalence relation $`AB(1)^{g(A)g(B)+f(A)f(B)}BA`$ divided out. The commutation relation above implies that ghosts fulfill commutation relations with physical fermions.
It is important that the elements of $`𝒫`$ are only symbols. In particular they are no operators in a Hilbert space and no functions on a manifold. The higher generators have no relation with the basic ones and the symbols do not satisfy field equations — e.g. $`g_{\mu \nu }u^{(2,\mu \nu )}0`$, where $`g_{\mu \nu }`$ is the metric tensor, although the ghost $`u`$ is a massless Klein-Gordon field in our example. Only the representation of the symbols as operator valued distributions in Fock space will restore these relations.
On $`𝒫`$ a derivative w.r.t. its generators is defined as a graded derivation according to
$$\begin{array}{cc}\hfill \frac{}{\phi _i}\left(AB\right)=\left(\frac{A}{\phi _i}\right)B+(1)^{f(A)f(\phi _i)+g(A)g(\phi _i)}A\left(\frac{B}{\phi _i}\right)& \\ \hfill \frac{\phi _i}{\phi _j}=\delta _{ij}1\mathrm{l}A,B𝒫,\phi _i,\phi _j𝒢& .\hfill \end{array}$$
(3.13)
The representation $``$ of the Lorentz group (or its covering group $`SL(2,\text{ }\mathrm{C})`$ for the spinors) on $`𝒫`$ is defined as follows: It acts as an algebra homomorphism, i.e. a linear mapping for which
$$_\mathrm{\Lambda }\left(\underset{i}{}\phi _i\right)=\underset{i}{}\left(_\mathrm{\Lambda }(\phi _i)\right),\mathrm{\Lambda }𝔏_+^{},\phi _i𝒢.$$
(3.14)
where $`𝔏_+^{}`$ is the homogeneous proper Lorentz group. The action of $``$ on the generators is the same as for the corresponding field operators. For the basic generators this means the following: Suppose the generator $`(\phi _i)^{(0)}`$ corresponds to the basic field $`\phi _i(x)`$, $`(\phi _i)^{(0)}\phi _i(x)`$, and the basic field transforms according to<sup>12</sup><sup>12</sup>12The field operators and the action $`U`$ of the Lorentz group on them are constructed in the next chapter
$$U(\mathrm{\Lambda })\phi _i(x)U^1(\mathrm{\Lambda })=\underset{j}{}\left(_\mathrm{\Lambda }\right)_{ij}\phi _j(\mathrm{\Lambda }^1x),\mathrm{\Lambda }𝔏_+^{}$$
(3.15)
for some numerical matrix $`\left(_\mathrm{\Lambda }\right)`$. Then the basic generator transforms according to
$$_\mathrm{\Lambda }\left((\phi _i)^{(0)}\right)\stackrel{\text{def}}{=}\underset{j}{}\left(_\mathrm{\Lambda }\right)_{ij}(\phi _j)^{(0)}$$
(3.16)
with the same numerical matrix $`\left(_\mathrm{\Lambda }\right)`$. In our standard example of Yang-Mills theory we have e.g.
$$_\mathrm{\Lambda }(A^\mu )=(\mathrm{\Lambda })_\nu ^\mu A^\nu ,_\mathrm{\Lambda }(u)=u,_\mathrm{\Lambda }(\stackrel{~}{u})=\stackrel{~}{u},\mathrm{\Lambda }𝔏_+^{}$$
(3.17)
Here $`(\mathrm{\Lambda })_\nu ^\mu `$ is the representative of $`\mathrm{\Lambda }`$ in the defining representation of $`𝔏_+^{}`$. The higher generators transform according to
$$_\mathrm{\Lambda }\left((\phi _i)^{(n,\mu _1\mathrm{}\mu _n)}\right)\stackrel{\text{def}}{=}\underset{j}{}\underset{\nu _1\mathrm{}\nu _n}{}(\mathrm{\Lambda })_{\nu _1}^{\mu _1}\mathrm{}(\mathrm{\Lambda })_{\nu _n}^{\mu _n}\left(_\mathrm{\Lambda }\right)_{ij}(\phi _j)^{(n,\nu _1\mathrm{}\nu _n)}.$$
(3.18)
with $`(\mathrm{\Lambda })`$ like above, and this completes the definition of $``$.
There is also an anti-linear involution defined on $`𝒫`$. It acts on products according to
$$(aAB)^{}=\overline{a}B^{}A^{}A,B𝒫,a\text{ }\mathrm{C},$$
(3.19)
where $`\overline{}`$ denotes complex conjugation in $`\mathrm{C}`$.
The involution is to implement the Krein adjoint for the fields in $`𝒫`$. So take a basic generator $`\phi _i`$ and a basic field $`\phi _i(x)`$ like above and let $`\left(\phi _i(x)\right)^{}=_ja_{ij}\phi _j(x)`$, where the -operation on the left hand side is the Krein adjoint on the fields. Then we define for this basic generator and its corresponding higher generators
$$\left((\phi _i)^{(n,\mu _1\mathrm{}\mu _n)}\right)^{}\stackrel{\text{def}}{=}\underset{j}{}a_{ij}(\phi _j)^{(n,\mu _1\mathrm{}\mu _n)}.$$
(3.20)
In anticipation of the results presented in the next section we state what this means for the basic generators in the standard example:
$$(A^\mu )^{}=A^\mu ,(u)^{}=u,(\stackrel{~}{u})^{}=\stackrel{~}{u}(\psi )^{}=\overline{\psi }\gamma ^0,(\overline{\psi })^{}=\gamma ^0\psi .$$
(3.21)
### 3.2. Fock space and Fock space operators
In this section we will construct the field operators already mentioned as operator valued distributions in the Fock space. We begin with some notations: A four-vector $`p`$ on the forward light cone $`\overline{V}_+`$ will be denoted as<sup>13</sup><sup>13</sup>13The construction is outlined here for massless fields, for simplicity.
$$\widehat{p}\stackrel{\text{def}}{=}(E_p,𝐩),E_p\stackrel{\text{def}}{=}\sqrt{𝐩^2}.$$
(3.22)
The invariant volume measure on the light cone and its Dirac distribution are defined as usual:
$$d\widehat{p}\stackrel{\text{def}}{=}\frac{d^3𝐩}{2(2\pi )^3E_p}\delta (\widehat{p})\stackrel{\text{def}}{=}2(2\pi )^3E_p\delta (𝐩).$$
(3.23)
At first we must construct the Fock space $`_{\phi _i}`$ for each basic field that corresponds to a basic generator $`\phi _i𝒢_b`$. That means for our standard example $`\phi _i=(A_\mu ^a),(u^a,\stackrel{~}{u}_a)`$ or $`(\psi ^r,\overline{\psi }_r)`$, where $`a`$ and $`r`$ are possible internal indices. To this end we begin with the $`n`$-particle Hilbert space $`_{\phi _i}^n`$. It is the Hilbert space of $`L^2(d\widehat{p}_1\mathrm{}d\widehat{p}_n,M^n)`$ functions of $`n`$ momenta and $`n`$ sets of indices (group-, colour- and Lorentz indices, for example) which are collectively written as $`a_i`$:
$$\phi _{(a_1\mathrm{}a_n)}^n(\widehat{p}_1,\mathrm{},\widehat{p}_n)_{\phi _i}^n$$
(3.24)
These functions are completely symmetric or antisymmetric under transposition of momenta and indices, $`(\widehat{p}_i,a_i)(\widehat{p}_j,a_j)`$, depending on the bosonic or fermionic character of $`\phi _i`$.
The scalar product on $`_{\phi _i}^n`$ is then defined as
$$(\psi ^n,\varphi ^n)\stackrel{\text{def}}{=}\underset{a_1\mathrm{}a_n}{}𝑑\widehat{p}_1\mathrm{}𝑑\widehat{p}_n\overline{\psi }_{(a_1\mathrm{}a_n)}^n(\widehat{p}_1,\mathrm{},\widehat{p}_n)\varphi _{(a_1\mathrm{}a_n)}^n(\widehat{p}_1,\mathrm{},\widehat{p}_n)$$
(3.25)
This scalar product is positive and allows to define a norm $`\varphi ^n\stackrel{\text{def}}{=}(\varphi ^n,\varphi ^n)^{\frac{1}{2}}`$.
With $`_{\phi _i}^0\stackrel{\text{def}}{=}\text{ }\mathrm{C}`$ and $`(\varphi ^0,\psi ^0)\stackrel{\text{def}}{=}\overline{\varphi ^0}\psi ^0`$ we can define the Fock space $`_{\phi _i}`$ for the field $`\phi _i`$ as
$$_{\phi _i}\stackrel{\text{def}}{=}\underset{n=0}{\overset{\mathrm{}}{}}_{\phi _i}^n,(\varphi ,\psi )=\underset{n=0}{\overset{\mathrm{}}{}}(\varphi ^n,\psi ^n),$$
(3.26)
where $`_{\phi _i}`$ contains only sequences $`\varphi `$ with $`(\varphi ,\varphi )<\mathrm{}`$. The vector $`|\omega _{\phi _i}\stackrel{\text{def}}{=}(1,0,0,\mathrm{})`$ is the vacuum for this Fock space.
Next we define $`𝒟_{\phi _i}`$ as the dense subspace of $`_{\phi _i}`$ that includes only elements with a finite particle number and whose wave functions are Schwartz’ test functions:
$$\varphi 𝒟_{\phi _i}_{\phi _i}m\mathrm{I}\mathrm{N}:\varphi \underset{n=0}{\overset{m}{}}𝒮(M^n)_{\phi _i}$$
(3.27)
where $`𝒮(M^n)`$ is the space of Schwartz’ test functions on $`M^n`$. This subspace has the advantage that Wick products are well defined operators acting on it \[GW64\]. It is the common domain of all operators on $`_{\phi _i}`$ defined below. Recently Brunetti and Fredenhagen \[BF99\] have found a definition of Wick products that is well posed on a bigger dense subspace than $`𝒟_{\phi _i}`$, but we will stick in this thesis to the space $`𝒟_{\phi _i}`$ defined above.
Annihilation operators may be defined on $`𝒟_{\phi _i}`$ according to
$$\begin{array}{cc}& v_a(\widehat{p}):𝒟_{\phi _i}𝒟_{\phi _i},\hfill \\ & [v_a(\widehat{p})\varphi ]_{(a_1\mathrm{}a_n)}^{(n)}(\widehat{p}_1,\mathrm{},\widehat{p}_n)=\sqrt{n+1}\varphi _{(a,a_1\mathrm{}a_n)}^{(n+1)}(\widehat{p},\widehat{p}_1,\mathrm{},\widehat{p}_n).\hfill \end{array}$$
(3.28)
Their adjoint — w.r.t. the scalar product defined above — operators $`v_i^+(p)`$, the creation operators $`v_a^+(\widehat{p})`$, are defined as
$$\begin{array}{cc}& [v_a^+(\widehat{p})\varphi ]_{(a_1\mathrm{}a_n)}^{(n)}(\widehat{p}_1,\mathrm{},\widehat{p}_n)=\hfill \\ & =\sqrt{n}(\delta _{a,a_1}\delta (\widehat{p}_1\widehat{p})\varphi _{(a_2\mathrm{}a_n)}^{(n1)}(\widehat{p}_2,\mathrm{},\widehat{p}_n)\hfill \\ & \pm \underset{k=2}{\overset{n}{}}\delta _{a,a_k}\delta (\widehat{p}_k\widehat{p})\varphi _{(a_1\mathrm{}\stackrel{ˇ}{a}_k\mathrm{}a_n)}^{(n1)}(\widehat{p}_2\mathrm{}\stackrel{ˇ}{p}_k\mathrm{}\widehat{p}_n)).\hfill \end{array}$$
(3.29)
Here $`\stackrel{ˇ}{}`$ means omission of the corresponding argument and the plus sign occurs if the field is bosonic, the minus sign if it is fermionic. The creation operators are no endomorphisms of $`𝒟_{\phi _i}`$ but map $`𝒟_{\phi _i}`$ to $`𝒟_{\phi _i}^{}`$, the dual space of $`𝒟_{\phi _i}`$. This is due to the appearance of the delta function in their definition.
Creation and annihilation operators fulfill the usual (anti-) commutation relations,
$$[v_a^+(\widehat{p}),v_b(\widehat{q})]_{}=\delta _{ab}\delta (\widehat{p}\widehat{q}),[v_a^+(\widehat{p}),v_b^+(\widehat{q})]_{}=[v_a(\widehat{p}),v_b(\widehat{q})]_{}=0$$
(3.30)
for bosons and ghosts (where the commutator above is the graded one) and
$$\{v_a^+(\widehat{p}),v_b(\widehat{q})\}_+=\delta _{a,b}\delta (\widehat{p}\widehat{q}),\{v_a^+(\widehat{p}),v_b^+(\widehat{q})\}_+=\{v_a(\widehat{p}),v_b(\widehat{q})\}_+=0$$
(3.31)
for spinors. Here $`v_a(\widehat{p})`$ is the annihilator for the field that is conjugate to the field with the annihilator $`v_a(\widehat{p})`$.
The normal ordering — or Wick ordering — of an arbitrary product of creation and annihilation operators is defined as the same product with all the annihilation operators on the right and all the creation operators on the left. The normal product of a product $`W`$ is denoted as $`:W:`$.
Operators on the Fock space can be defined from these distributional operators according to
$$v_a(f)=𝑑\widehat{p}\overline{f(\widehat{p})}v_a(\widehat{p}),v_a^+(f)=𝑑\widehat{p}f(\widehat{p})v_a^+(\widehat{p}).$$
(3.32)
With this smearing also the Wick products become operators in $`(𝒟_{\phi _i})`$.
The field operators defined below are operator valued distributions acting on the dense subspace $`𝒟_{\phi _i}`$. To give a precise meaning to that expression, we define the $`n^{\mathrm{th}}`$ order operator valued distributions on an arbitrary subspace $`𝒟`$ of a Fock space, abbreviated as $`_n(𝒟)`$, as $`\mathrm{C}`$-linear strongly continuous mappings
$$_n(𝒟)\stackrel{\text{def}}{=}\{A:𝒟(M^n)(𝒟)\}.$$
(3.33)
where $`M`$ is the Minkowski space and $`𝒟M^n`$ the space of test functions on $`M^n`$ with compact support.
The field operators defined below are in $`_1(𝒟_{\phi _i})`$.
We begin the definition of the field operators that correspond to the basic generators with the vector bosons. The corresponding Fock space is denoted as $`_A`$, its dense subspace as $`𝒟_A`$. The creation and annihilation operators are denoted as $`a_\mu ^{a,+}(\widehat{p})`$ and $`a_\mu ^a(\widehat{p})`$. They fulfill the commutation relations
$$[a_\mu ^{+,a}(\widehat{p}),a_\nu ^b(\widehat{q})]_{}=\delta ^{ab}\delta _{\mu \nu }\delta (\widehat{p}\widehat{q}),[a_\mu ^{+,a}(\widehat{p}),a_\nu ^{+,b}(\widehat{q})]_{}=[a_\mu ^a(\widehat{p}),a_\nu ^b(\widehat{q})]_{}=0.$$
(3.34)
The vector boson field is defined as
$$\begin{array}{cc}& A_0^a(x)\stackrel{\text{def}}{=}d\widehat{p}[a_0^a(\widehat{p})e^{i\widehat{p}x}a_0^{a,+}(\widehat{p})e^{i\widehat{p}x}]_1(𝒟_A),\hfill \\ & A_i^a(x)\stackrel{\text{def}}{=}d\widehat{p}[a_i^a(\widehat{p})e^{i\widehat{p}x}+a_i^{a,+}(\widehat{p})e^{i\widehat{p}x}]_1(𝒟_A).\hfill \end{array}$$
(3.35)
It satisfies the commutation relation
$$[A_\mu ^a(x),A_\nu ^b(y)]_{}=i\delta ^{ab}g_{\mu \nu }D(xy)$$
(3.36)
and the massless Klein-Gordon equation
$$\mathrm{}^xA_\mu ^a(x).$$
(3.37)
Here $`D(x)`$ is the massless Pauli-Jordan function
$$D(x)\stackrel{\text{def}}{=}2i𝑑\widehat{p}\mathrm{sin}(\widehat{p}x).$$
(3.38)
It has causal support. It may be split into a positive and a negative frequency part according to
$$D^+(x)\stackrel{\text{def}}{=}𝑑\widehat{p}e^{i\widehat{p}x},D^{}(x)\stackrel{\text{def}}{=}D^+(x).$$
(3.39)
Its corresponding retarded, advanced and Feynman propagators $`D^R,D^A`$ and $`D^F`$ are defined as
$$D^R(x)\stackrel{\text{def}}{=}\theta (x^0)D(x),D^A(x)\stackrel{\text{def}}{=}\theta (x^0)D(x),D^F(x)\stackrel{\text{def}}{=}D^R(x)D^{}(x).$$
(3.40)
They are the inverse of the massless Klein-Gordon operator:
$$\mathrm{}^xD^{R,A,F}(x)=\delta (x).$$
(3.41)
Clearly $`D^R`$ has retarded and $`D^A`$ has advanced support.
The $`0`$-component of the vector bosons is anti hermitian, $`(A_0^a)^+=A_0^a`$. Furthermore the scalar product is not Lorentz invariant as can be easily verified already in the one-particle space. This is a typical situation in gauge theories as described in the last chapter. To find a physical inner product on $`_A`$ one must find a suitable Krein operator $`J_A`$ acting on it and define
$$\varphi ,\psi \stackrel{\text{def}}{=}(\varphi ,J_A\psi ).$$
(3.42)
This suitable Krein operator is
$$J_A\stackrel{\text{def}}{=}(1)^{N_0},N_0=\underset{b}{}𝑑\widehat{p}a_0^{b,+}(\widehat{p})a_0^b(\widehat{p}),$$
(3.43)
where $`N_0`$ is the number operator for $`A^0(x)`$ with eigenvalues in $`\mathrm{I}\mathrm{N}`$. It is obviously hermitian, $`J=J^+`$, and idempotent, $`J^2=1\mathrm{l}`$. With the -involution
$$B^{}\stackrel{\text{def}}{=}J_AB^+J_A,B(𝒟_A),$$
(3.44)
also called the Krein adjoint, the vector bosons become pseudo-hermitian, $`(A_\mu ^a)^{}=A_\mu ^a`$. Furthermore we find the inner product $`,`$ to define a Lorentz invariant norm, but it is indefinite.
The definition of the spinor Fock space $`_\psi `$ and the field operators $`\psi (x),\overline{\psi }(x)`$ acting therein proceeds in the same way and can be found in textbooks on quantum field theory. The fermions satisfy the commutation relations
$$[\psi (x),\overline{\psi }(y)]_{}=i(i/_x+m)D(xy)$$
(3.45)
and the field equations
$$(i/_xm)\psi =0,\overline{\psi }(i\stackrel{}{/}_xm)=0.$$
(3.46)
The Krein operator $`J_\psi `$ on the spinor Fock space is trivial, $`J_\psi =1\mathrm{l}`$.
On the Fock space for the ghosts, $`_u`$ with its dense subspace $`𝒟_u`$, creation and annihilation operators are denoted by $`b^{a,+}(\widehat{p})`$, $`c^{a,+}(\widehat{p})`$, $`b^a(\widehat{p})`$ and $`c^a(\widehat{p})`$, respectively. They fulfill the anti-commutation relations
$$\{b^{+,a}(\widehat{p}),b^b(\widehat{q})\}_+=\delta ^{ab}\delta (\widehat{p}\widehat{q}),\{c^{+,a}(\widehat{p}),c^b(\widehat{q})\}_+=\delta ^{ab}\delta (\widehat{p}\widehat{q})$$
(3.47)
and all other anti-commutators vanish. The ghost field $`u^a(x)`$ and the anti-ghost field $`\stackrel{~}{u}^a(x)`$ are defined as
$$\begin{array}{cc}& u^a(x)\stackrel{\text{def}}{=}d\widehat{p}[b^a(\widehat{p})e^{i\widehat{p}x}+c^{a,+}(\widehat{p})e^{i\widehat{p}x}]_1(𝒟_u),\hfill \\ & \stackrel{~}{u}^a(x)\stackrel{\text{def}}{=}d\widehat{p}[c^a(\widehat{p})e^{i\widehat{p}x}+b^{a,+}(\widehat{p})e^{i\widehat{p}x}]_1(𝒟_u).\hfill \end{array}$$
(3.48)
Then we get for the anti-commutators of the ghosts
$$\begin{array}{cc}& \{u^a(x),\stackrel{~}{u}^b(y)\}_+=i\delta ^{ab}D(xy).\hfill \\ & \{u^a(x),u^b(y)\}_+=\{\stackrel{~}{u}^a(x),\stackrel{~}{u}^b(y)\}_+=0.\hfill \end{array}$$
(3.49)
The Krein operator for the ghosts was explicitely determined by Krahe \[Kra95\] and reads
$$J_u=\mathrm{exp}\left(\frac{i\pi }{2}𝑑\widehat{p}\left[b^+(\widehat{p})b(\widehat{p})b^+(\widehat{p})c(\widehat{p})+c^+(\widehat{p})c(\widehat{p})c^+(\widehat{p})b(\widehat{p})\right]\right).$$
(3.50)
For us it is only important that this implies for the field operators
$$\left(u^a(x)\right)^{}=u^a(x)\text{and}\left(\stackrel{~}{u}^a(x)\right)^{}=\stackrel{~}{u}^a(x),$$
(3.51)
so the ghosts are pseudo-hermitian and the anti-ghosts are anti-pseudo-hermitian.
Now we introduce the pseudo-unitary representation of the proper Poincaré group $`𝔓_+^{}`$ in the individual Fock spaces. It reads for scalar fields like the ghosts
$$\begin{array}{cc}& [U(p)\varphi ]^{(0)}=\varphi ^{(0)},p=(a,\mathrm{\Lambda })𝔓_+^{}\hfill \\ & \begin{array}{cc}\hfill [U(p)\varphi ]_{(a_1\mathrm{}a_n)}^{(n)}(\widehat{q}_1,\mathrm{},\widehat{q}_n)=\mathrm{exp}(i(\widehat{q}_1+\mathrm{}+\widehat{q}_n)a)\times & \\ \hfill \times \varphi _{(a_1\mathrm{}a_n)}^{(n)}(\mathrm{\Lambda }\widehat{q}_1,\mathrm{},\mathrm{\Lambda }\widehat{q}_n)& .\hfill \end{array}\hfill \end{array}$$
(3.52)
For vector fields like the vector bosons it reads
$$\begin{array}{cc}& [U(p)\varphi ]^{(0)}=\varphi ^{(0)},p=(a,\mathrm{\Lambda })𝔓_+^{}\hfill \\ & \begin{array}{cc}\hfill \left([U(p)\varphi ]^{(n)}\right)_{(a_1\mathrm{}a_n)}^{(\mu _1\mathrm{}\mu _n)}(\widehat{q}_1,\mathrm{},\widehat{q}_n)=& \mathrm{exp}(i(\widehat{q}_1+\mathrm{}+\widehat{q}_n)a)\times \hfill \\ & \times (\mathrm{\Lambda })_{\nu _1}^{\mu _1}\mathrm{}(\mathrm{\Lambda })_{\nu _n}^{\mu _n}\left(\varphi ^{(n)}\right)_{(a_1\mathrm{}a_n)}^{(\nu _1\mathrm{}\nu _n)}(\mathrm{\Lambda }\widehat{q}_1,\mathrm{},\mathrm{\Lambda }\widehat{q}_n)\hfill \end{array}\hfill \end{array}$$
(3.53)
where the Lorentz indices $`\mu _i`$ have been separated from the other indices $`a_i`$ and summation over repeated indices is understood. The matrices $`(\mathrm{\Lambda })`$ are the representatives of $`\mathrm{\Lambda }`$ in the defining representation of $`𝔏_+^{}𝔓_+^{}`$, like above. For the spinors an analogous definition holds. The vacuum vector $`|\omega _{\phi _i}`$ is clearly Poincaré invariant. As was pointed out by Krahe \[Kra95\], it is also cyclic w.r.t. the field operators defined above.
The field operators transform according to
$$\begin{array}{cc}& U(p)u^a(x)U^1(p)=u^a(\mathrm{\Lambda }^1xa),U(p)\stackrel{~}{u}^a(x)U^1(p)=\stackrel{~}{u}^a(\mathrm{\Lambda }^1xa)\hfill \\ & U(p)A_\mu ^a(x)U^1(p)=\left(\mathrm{\Lambda }\right)_\mu ^\nu A_\nu ^a(\mathrm{\Lambda }^1xa).\hfill \end{array}$$
(3.54)
With the Fock spaces for the individual fields the Fock space of the entire theory $``$, its dense subspace $`𝒟`$ and the Krein operator $`J`$ acting on $``$ are defined as
$$\stackrel{\text{def}}{=}\underset{i}{}_{\phi _i}𝒟\stackrel{\text{def}}{=}\underset{i}{}𝒟_{\phi _i}J\stackrel{\text{def}}{=}\underset{i}{}J_{\phi _i}.$$
(3.55)
The vacuum vector of the Fock space $``$ is denoted by $`|\omega `$. We introduce the notation $`\omega _0\left(A\right)`$ for $`\omega \left|A\right|\omega `$ for every $`A(𝒟)`$. Here $`(𝒟)`$ is the algebra of endomorphisms of $`𝒟`$. An important fact concerning this algebra is that it has trivial centre. Even more, for an arbitrary element $`A(𝒟)`$ the following equivalence holds:
$$\begin{array}{cc}\hfill [A,T\left(\phi _i\right)(x)]_{}& =0\phi _i𝒢_b,\hfill \\ \hfill A& =a1\mathrm{l},a\text{ }\mathrm{C}.\hfill \end{array}$$
(3.56)
For the proof of this assertion see Scharf \[Sch95\], for example.
In chapter (6) it will turn out that spacetime must be compactified in spacelike directions for the BRS charge to be a well defined operator. Therefore it is important to construct the Fock space and the operators acting on it also for a quantum field theory in the compactified spacetime. This has been done by Dütsch and Fredenhagen in \[DF99, appendix A\]. We refer to their results, especially concerning the choice of boundary conditions, but we do not go here into details.
### 3.3. The linear representation of $`𝒫`$
In this section we define the $`\mathrm{C}`$-linear representation $`T`$ of the polynomials in $`𝒫`$ as operator valued distributions
$$T:𝒫_1(𝒟).$$
(3.57)
Linear representation means that the linear structure of $`𝒫`$ is preserved, but not its structure as an algebra. This comes from the fact that a pointwise product of distributions is in general no well defined operation.
The precise definition of $`T`$ will take three steps: at first it is defined for the basic generators, then for the higher generators and finally for composed elements of $`𝒫`$.
The first end has already been achieved with the definition of an operator valued distribution $`\phi _i(x)_1(𝒟)`$ for each basic generator $`\phi _i𝒢_b`$. The representative of the basic generator is defined as:
$$T\left(\phi _i\right)(x)\stackrel{\text{def}}{=}\phi _i(x),\phi _i𝒢_b,\phi _i(x)_1(𝒟).$$
(3.58)
This definition can work only for the basic generators since for the higher ones there are no corresponding free field operators.
For these generators we define
$$T\left((\phi _i)^{(n,\nu _1\mathrm{}\nu _n)}\right)(x)\stackrel{\text{def}}{=}_x^{\nu _1}\mathrm{}_x^{\nu _n}\phi _i(x),(\phi _i)^{(\mathrm{})}𝒢.$$
(3.59)
We remind the reader that there are no relations between the basic generators and the higher generators in $`𝒫`$, and that there are no field equations in $`𝒫`$. But with the definition above there is a relation established between the representatives of the basic and those of the higher generators, and the former clearly satisfy field equations. So the linear representation is not faithful.
For the representation of the composed elements in $`𝒫`$ we define at first the commutator function
$$i\mathrm{\Delta }_{ij}(xy)=[T(\phi _i)(x),T(\phi _j)(y)]_{},\phi _i,\phi _j𝒢.$$
(3.60)
Here $`i`$ and $`j`$ take on values also for the higher generators. With this commutator function we give an implicit definition of the representation of monomials in $`𝒫`$, namely
$$\begin{array}{cc}\hfill [T(W)(x),T(\phi _i)(y)]_{}& =i\underset{j}{}T\left(\frac{W}{\phi _j}\right)(x)\mathrm{\Delta }_{ij}(xy)\hfill \\ \hfill \omega _0\left(T(W)(x)\right)& =0W𝒫.\hfill \end{array}$$
(3.61)
The existence of a solution is guaranteed by the observation that the normal products solve both equations. Suppose $`A=_i\phi _i𝒫,\phi _i𝒢`$, then the normal product $`:_iT\left(\phi _i\right)(x):_1(𝒟)`$ is indeed a searched for solution.
The uniqueness of this solution can be seen inductively. Suppose, the representation for all monomials containing at most $`k1`$ generators is defined. Then the commutator condition determines the solution for monomials of $`k`$ generators up to a $`\mathrm{C}`$-number distribution — this is due to eqn. (3.56). The $`\mathrm{C}`$-number part is determined by the second condition — it is zero. So the equations above give indeed a definition for the representation of monomials in $`𝒫`$. This completes the definition of the representation.
As for the Pauli-Jordan function $`D(x)`$ we can find a positive and a negative frequency solution for the commutator function $`\mathrm{\Delta }_{ij}(x)`$. The two point function — or positive frequency part of $`\mathrm{\Delta }`$ — is denoted as $`\mathrm{\Delta }^+`$ and defined as
$$i\mathrm{\Delta }_{ij}^+(xy)\stackrel{\text{def}}{=}\omega _0\left(T\left(\phi _i\right)(x)T\left(\phi _j\right)(y)\right),$$
(3.62)
the negative frequency part of $`\mathrm{\Delta }`$ is defined as
$$\mathrm{\Delta }_{ij}^{}(x)\stackrel{\text{def}}{=}\mathrm{\Delta }_{ij}(x)\mathrm{\Delta }_{ij}^+(x).$$
(3.63)
### 3.4. The propagator functions
In this section we define propagator functions $`\mathrm{\Delta }_{ij}^R(x),\mathrm{\Delta }_{ij}^A(x)`$ analogous to $`D^R(x)`$ and $`D^A(x)`$ that are restrictions of $`\mathrm{\Delta }_{ij}(x)`$ to the past and future light cone, such that $`\mathrm{\Delta }_{ij}^R(x)\mathrm{\Delta }_{ij}^A(x)=\mathrm{\Delta }_{ij}(x)`$. Simultaneously we search for a differential operator $`D_{ij}^x`$ that takes over the part of the Klein-Gordon operator, i.e. that fulfills the equations
$$\underset{j}{}D_{ij}^x\mathrm{\Delta }_{jk}^{R,A}(x)=\delta _{ik}\delta (x)\underset{j}{}D_{ij}^x\mathrm{\Delta }_{jk}(x)=0.$$
(3.64)
This means in particular that the propagators must be invertible with $`D_{ij}^x`$ as their inverse.
To see what form the propagators might have we take a closer look at the commutator function. If $`𝒢`$ contains $`r`$ generators, this is an $`r\times r`$-matrix. It has the following block diagonal structure:
$$\mathrm{\Delta }(x)=\left(\begin{array}{cccc}\mathrm{\Delta }^{\phi _1}(x)& 0& 0& \mathrm{}\\ 0& \mathrm{\Delta }^{\phi _2}(x)& 0& \mathrm{}\\ 0& 0& \mathrm{\Delta }^{\phi _3}(x)& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right).$$
(3.65)
Here the matrices $`\mathrm{\Delta }^{\phi _i}`$ are of two different types. The first type corresponds to field operators that have no distinct conjugate field like the uncharged vector bosons. Then the index $`\phi _i`$ corresponds to the field, e.g. $`\phi _i=A_\mu `$. The other type corresponds to field operators $`\phi _i`$ that do have such a distinct conjugate field $`\stackrel{~}{\phi }_i`$ like ghosts with the anti-ghosts. In this case the field and the conjugated field form one common block in the matrix and the index $`\phi _i`$ corresponds to the field and its conjugated field, e.g. $`\phi _i=u,\stackrel{~}{u}`$. All blocks include also the commutators of the derivatives as far as higher generators exist in $`𝒢`$ that correspond to these derivatives. In our standard example it has the form
$$\mathrm{\Delta }(x)=\left(\begin{array}{ccc}\mathrm{\Delta }^A(x)& 0& 0\\ 0& \mathrm{\Delta }^{u,\stackrel{~}{u}}(x)& 0\\ 0& 0& \mathrm{\Delta }^{\psi ,\overline{\psi }}(x)\end{array}\right)$$
(3.66)
if QED is treated where no internal indices appear. In Yang-Mills theory, where internal indices do appear, there is an individual block for each index on the diagonal. From now on we will disregard internal indices for their inclusion is straightforward. The vector boson part has the form
$$\mathrm{\Delta }^A(x)\stackrel{\text{def}}{=}g_{\mu \nu }\left(\begin{array}{ccc}D(x)& _x^{\nu _1}D(x)& _x^{\nu _1}_x^{\nu _2}D(x)\\ _x^{\rho _1}D(x)& _x^{\nu _1}_x^{\rho _1}D(x)& _x^{\nu _1}_x^{\rho _1}_x^{\nu _2}D(x)\\ _x^{\rho _2}_x^{\rho _1}D(x)& _x^{\nu _1}_x^{\rho _2}_x^{\rho _1}D(x)& _x^{\nu _1}_x^{\rho _2}_x^{\rho _1}_x^{\nu _2}D(x)\end{array}\right),$$
(3.67)
the ghost part
$$\mathrm{\Delta }^{u,\stackrel{~}{u}}(x)\stackrel{\text{def}}{=}\left(\begin{array}{cccc}0& 0& D(x)& _x^{\nu _1}D(x)\\ 0& 0& _x^{\rho _1}D(x)& _x^{\rho _1}_x^{\nu _1}D(x)\\ D(x)& _x^{\nu _1}D(x)& 0& 0\\ _x^{\rho _1}D(x)& _x^{\rho _1}_x^{\nu _1}D(x)& 0& 0\end{array}\right)$$
(3.68)
and the spinor part
$$\mathrm{\Delta }^{\psi ,\overline{\psi }}(x)\stackrel{\text{def}}{=}\left(\begin{array}{cc}0& i(i/+m)D_m(x)\\ i(i/m)D_m(x)& 0\end{array}\right).$$
(3.69)
The distribution $`D_m`$ is the Pauli-Jordan function for mass $`m`$. The matrices are given here in the basis
$$((A_\mu )^{(0)},(A_\mu )^{(1,\nu _1)},(A_\mu )^{(2,\nu _1\nu _2)})^t$$
(3.70)
for the vector bosons,
$$((u)^{(0)},(u)^{(1,\nu _1)},(\stackrel{~}{u})^{(0)},(\stackrel{~}{u})^{(1,\nu _1)})^t$$
(3.71)
for the ghosts and anti-ghosts and
$$((\psi )^{(0)},(\overline{\psi })^{(0)})^t,$$
(3.72)
for the spinors, where <sup>t</sup> denotes transposition.
The natural attempt would be to replace the Pauli-Jordan function $`D(x)`$ by its retarded, $`D^R(x)`$, or advanced, $`D^A(x)`$, propagator in each entry to define the matrices $`\mathrm{\Delta }_{ij}^R(x)`$ and $`\mathrm{\Delta }_{ij}^A(x)`$. These would clearly be well defined distributions with the desired support properties, but they would not be invertible. This comes from the fact that with this definition each row would be the derivative of the row above, and therefore the determinant — w.r.t. convolution — of these matrices would vanish.
To improve the definition above we observe that the matrices $`\mathrm{\Delta }_{ij}^{R,A}(x)`$ are defined by their desired support properties — $`\mathrm{\Delta }_{ij}^R(x)\overline{V}_+`$ and $`\mathrm{\Delta }_{ij}^A(x)\overline{V}_{}`$ — and their relation to the commutator function everywhere but in the origin. That means that we may alter the propagator functions only at the origin, i.e. by delta distributions or its derivatives at the individual entries. As a further restriction of possible propagators we demand that this modification does not increase the scaling degree (see below) of the individual entries and that it does not change the Lorentz transformation property of that entry.
Scaling degree means the following: For every numerical distribution $`d`$ one can define a dilated distribution
$$d_\lambda (x)=d(\lambda x)\lambda \mathrm{I}\mathrm{R}^+\left\{0\right\}.$$
(3.73)
Clearly $`d_\lambda `$ is a numerical distribution, too. Then the scaling degree $`(d)`$ of $`d`$ w.r.t. the origin is defined, according to Steinmann \[Ste71\], as
$$(d)\stackrel{\text{def}}{=}inf\{\beta \mathrm{I}\mathrm{R}:\underset{\lambda 0}{lim}\lambda ^\beta d_\lambda =0\},$$
(3.74)
where the equation in the bracket holds in the sense of distributions.
The restriction on the scaling degree fixes some entries uniquely, e.g. $`\mathrm{\Delta }_{ij}^R(x)=D^R(x)`$ if $`\phi _i=u`$ and $`\phi _j=\stackrel{~}{u}`$. For others there remains a certain ambiguity, e.g.
$$\mathrm{\Delta }_{ij}^R(x)=g_{\mu \nu }^\rho ^\sigma D^R(x)Cg_{\mu \nu }g^{\rho \sigma }\delta (x)$$
(3.75)
if $`\phi _i=(A_\mu )^{(1,\rho )}`$ and $`\phi _j=(A_\nu )^{(1,\sigma )}`$. The numerical constant $`C`$ is then arbitrary. In the following we define propagator functions with an inverse that is a differential operator and we will give later the explicit form of these differential operators.
The propagators have the same block diagonal structure as the commutator function:
$$\mathrm{\Delta }^{R,A}(x)=\left(\begin{array}{cccc}\mathrm{\Delta }_{R,A}^{\phi _1}(x)& 0& 0& \mathrm{}\\ 0& \mathrm{\Delta }_{R,A}^{\phi _2}(x)& 0& \mathrm{}\\ 0& 0& \mathrm{\Delta }_{R,A}^{\phi _3}(x)& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right).$$
(3.76)
In the following we will consider only the construction of the retarded propagator. The advanced propagator is defined as $`\mathrm{\Delta }^A=\mathrm{\Delta }^R\mathrm{\Delta }`$. For the determination of the individual blocks we notice that usually the $`(0,0)`$-component<sup>14</sup><sup>14</sup>14We start the numbering of columns and rows with zero, such that the index of a column or row agrees with the degree of the corresponding generator of the commutator function has a scaling degree smaller than the spacetime dimension, so that its retarded solution is uniquely determined by the following condition:
$$\mathrm{\Delta }_{00}^{R,\phi _i}(x)=\mathrm{\Delta }_{00}^{\phi _i}(x)x\overline{V}_+,\mathrm{\Delta }_{00}^R\overline{V}_+.$$
(3.77)
With this the general matrix element of a block $`\mathrm{\Delta }_R^{\phi _i}`$ of the retarded propagator can be written as
$$\mathrm{\Delta }_{jk}^{R,\phi _i}(x)=(1)^k^{\mu _1}\mathrm{}^{\mu _k}^{\nu _1}\mathrm{}^{\nu _j}\mathrm{\Delta }_{00}^{R,\phi _i}(x)+(1)^k\delta _{jk}C_{\phi _i,k}\delta (x).$$
(3.78)
(no summation over $`k`$ in the last term). The constants $`C_{\phi _i,k}`$ are non zero real numbers, $`C_{\phi _i,k}\mathrm{I}\mathrm{R}\left\{0\right\}`$. As these constants will determine the normalization of higher order time ordered products (c.f. next chapter), they will be called normalization constants.
In our standard example the propagator has the form
$$\mathrm{\Delta }^R(x)=\left(\begin{array}{ccc}\mathrm{\Delta }_R^A(x)& 0& 0\\ 0& \mathrm{\Delta }_R^{u,\stackrel{~}{u}}(x)& 0\\ 0& 0& \mathrm{\Delta }_R^{\psi ,\overline{\psi }}(x)\end{array}\right).$$
(3.79)
For the vector boson block of the retarded propagator, $`\mathrm{\Delta }_R^A`$, we omit all spacetime arguments because it otherwise would not fit into the line. Then it reads
$$\mathrm{\Delta }_R^A=g_{\mu \nu }\left(\begin{array}{ccc}D^R& ^\nu D^R& ^\nu ^\rho D^R\\ ^\sigma D^R& ^\nu ^\sigma D^RC_{A,1}g^{\nu \sigma }\delta & ^\nu ^\rho ^\sigma D^R\\ ^\sigma ^\tau D^R& ^\nu ^\sigma ^\tau D^R& ^\nu ^\rho ^\sigma ^\tau D^RC_{A,2}g^{\nu \sigma }g^{\rho \tau }\delta \end{array}\right),$$
(3.80)
For the ghosts we get the contribution,
$$\mathrm{\Delta }_R^{u,\stackrel{~}{u}}(x)=\left(\begin{array}{cc}0& d_R^u(x)\\ d_R^u(x)& 0\end{array}\right)$$
(3.81)
with the $`2\times 2`$-matrices
$$d_R^u(x)=\left(\begin{array}{cc}D^R(x)& _x^{\nu _1}D^R(x)\\ _x^{\rho _1}D^R(x)& _x^{\nu _1}_x^{\rho _1}D^R(x)C_{u,1}g^{\nu _1\rho _1}\delta (x)\end{array}\right).$$
(3.82)
The spinors finally give
$$\mathrm{\Delta }_R^{\psi ,\overline{\psi }}(x)=\left(\begin{array}{cc}0& (i/+m)D_m^R(x)\\ (i/m)D_m^R(x)& 0.\end{array}\right)$$
(3.83)
$`D_m^R`$ is the retarded part of $`D_m`$. All the matrices are given in the same basis as for the commutator function. The retarded propagator function has obviously retarded support and agrees with the commutator function outside the forward light cone.
The advanced propagator $`\mathrm{\Delta }^A=\mathrm{\Delta }^R\mathrm{\Delta }`$ has advanced support and agrees with the commutator function outside the backward light cone. In the example above the respective advanced propagators can be derived from the retarded propagators by a substitution of $`D^R`$ with $`D^A`$.
We define also a Feynman propagator
$$\mathrm{\Delta }^F:\mathrm{\Delta }_{ij}^F(x)\stackrel{\text{def}}{=}\mathrm{\Delta }_{ij}^R(x)\mathrm{\Delta }_{ij}^{}(x).$$
(3.84)
Now we come to the differential operator valued matrix $`D^x`$ that inverts the propagators defined above, i.e. for which the equation
$$\underset{j}{}D_{ij}^x\mathrm{\Delta }_{jk}^{R,A,F}(x)=\delta _{ik}\delta (x)$$
(3.85)
holds. It is an $`r\times r`$ matrix, where $`r`$ was the number of generators in $`𝒢`$. It has the usual block diagonal form:
$$D^x=\left(\begin{array}{cccc}D^{\phi _1,x}& 0& 0& \mathrm{}\\ 0& D^{\phi _2,x}& 0& \mathrm{}\\ 0& 0& D^{\phi _3,x}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right)$$
(3.86)
or, in our standard example,
$$D^x=\left(\begin{array}{ccc}D^{A,x}& 0& 0\\ 0& D^{u,\stackrel{~}{u},x}& 0\\ 0& 0& D^{\psi ,\overline{\psi },x}\end{array}\right).$$
(3.87)
Like for the propagators the individual blocks correspond to field operators or pairs of conjugated fields. We define the blocks for single fields as $`(s+1)\times (s+1)`$-matrices if higher generators up to degree $`s`$ are included in $`𝒢`$ for that field. Let $`K^{\phi _i,x}`$ be the differential operator that defines the field equation for $`\phi _i(x)`$, i.e. which fulfills the equation
$$K^{\phi _i,x}\mathrm{\Delta }_R^{\phi _i}(x)=\delta (x),$$
(3.88)
e.g. $`K^{A,x}=\mathrm{}^x`$. Then the corresponding block is written in the basis
$$((\phi _i)^{(0)},(\phi _i)^{(1,\nu _1)},\mathrm{},(\phi _i)^{(s,\nu _1\mathrm{}\nu _s)})^t$$
(3.89)
as the matrix with the components
$$\begin{array}{cc}\hfill D_{00}^{\phi _i,x}& =\left(K^{\phi _i,x}\underset{k=1}{\overset{n}{}}(1)^kC_{\phi _i,k}^1\mathrm{}^k\right)\hfill \\ \hfill D_{0k}^{\phi _i,x}& =(1)^kC_{\phi _i,k}^1\left(^{\nu _1}\mathrm{}^{\nu _k}\right)\hfill \\ \hfill D_{j0}^{\phi _i,x}& =C_{\phi _i,j}^1\left(^{\sigma _1}\mathrm{}^{\sigma _j}\right)\hfill \\ \hfill D_{kk}^{\phi _i,x}& =C_{\phi _i,k}^1\left(g^{\nu _1\sigma _1}\mathrm{}g^{\nu _k\sigma _k}\right)\hfill \\ \hfill D_{jk}^{\phi _i,x}& =0\text{otherwise}.\hfill \end{array}$$
(3.90)
where the constants $`C_{\phi _i,k}`$ are those determined in the propagator functions.
Again we exemplify the definition above for our standard example. The only uncharged fields there are the vector bosons. In the same basis as for the commutator function, the block $`D^{A,x}`$ has the form
$$D^A\stackrel{\text{def}}{=}\left(\begin{array}{ccc}(1+C_{A,1}^1)\mathrm{}C_{A,2}^1\mathrm{}^2& C_{A,1}^1^{\nu _1}& C_{A,2}^1^{\nu _1}^{\nu _2}\\ C_{A,1}^1^{\rho _1}& C_{A,1}^1g^{\rho _1\nu _1}& 0\\ C_{A,2}^1^{\rho _1}^{\rho _2}& 0& C_{A,2}^1g^{\rho _1\nu _1}g^{\rho _2\nu _2}\end{array}\right).$$
(3.91)
For the charged fields we construct according to the rules above one block $`D^{\phi _i,x}`$ for the fields $`\phi _i`$ and one block $`D^{\stackrel{~}{\phi }_i,x}`$ for the conjugated fields $`\stackrel{~}{\phi }_i`$. Like for the propagators, the combined block for the fields and conjugated fields reads then
$$D^{\phi _i,\stackrel{~}{\phi }_i,x}\stackrel{\text{def}}{=}\left(\begin{array}{cc}0& D^{\stackrel{~}{\phi }_i,x}\\ D^{\phi _i,x}& 0\end{array}\right)$$
(3.92)
in the basis
$$((\phi _i)^{(0)},\mathrm{},(\phi _i)^{(s,\nu _1\mathrm{}\nu _s)},(\stackrel{~}{\phi }_i)^{(0)},\mathrm{},(\stackrel{~}{\phi }_i)^{(s,\nu _1\mathrm{}\nu _s)})^t,$$
(3.93)
if higher generators up to degree $`s`$ are included. The expressions for our standard example, i.e. for the ghosts and the spinors, read then
$$D^{u,\stackrel{~}{u}}\stackrel{\text{def}}{=}\frac{1}{C_{u,1}}\left(\begin{array}{cccc}0& 0& (1+C_{u,1})\mathrm{}& ^{\nu _1}\\ 0& 0& ^{\rho _1}& g^{\nu _1\rho _1}\\ (1+C_{u,1})\mathrm{}& ^{\nu _1}& 0& 0\\ ^{\rho _1}& g^{\nu _1\rho _1}& 0& 0\end{array}\right).$$
(3.94)
For the spinors no higher generators are included in our example, so they contribute the expression
$$D^\psi \stackrel{\text{def}}{=}\left(\begin{array}{ccc}0& (i/+m)& \\ (i/m)& 0\end{array}\right).$$
(3.95)
where the operator in the second line acts from the right.
It is easily verified by direct calculation that this differential operator really inverts the propagators. Furthermore, the representatives of the generators satisfy the following free field equations:
$$\underset{j}{}D_{ij}^xT(\phi _j)(x)=0.$$
(3.96)
Here the sum runs over all generators. This equation holds independently of the choice of the normalization constants $`C_{\phi _i,k}`$. From its definition it is already clear that the commutator function is annihilated by $`D^x`$:
$$\underset{j}{}D_{ij}^x\mathrm{\Delta }_{jk}(x)=0.$$
(3.97)
If $`D^x`$ is determined, the propagator functions $`\mathrm{\Delta }^R,\mathrm{\Delta }^A`$ and $`\mathrm{\Delta }^F`$ are uniquely determined by the following conditions:
* $`\mathrm{\Delta }^{R,A,F}(x)`$ must fulfil eqn. (3.85)
* $`\mathrm{\Delta }^R(x)=\mathrm{\Delta }(x)x\overline{V}_{}`$ and $`\mathrm{\Delta }^R(x)=0x\overline{V}_{}\left\{0\right\}`$
* $`\mathrm{\Delta }^A(x)=\mathrm{\Delta }^R(x)\mathrm{\Delta }(x)`$
* $`\mathrm{\Delta }^F(x)=\mathrm{\Delta }^R(x)\mathrm{\Delta }^+(x)`$ .
So $`D^x`$ is a relativistically covariant, hyperbolic differential operator with a unique solution for the Cauchy problem. In particular the normalization constants $`C_{\phi _i,k}`$ that appear in the propagators are uniquely determined by their choice in the differential operator $`D^x`$.
We do not claim that our choice for the operator $`D^x`$ or the propagators is the most general one. But we point out that there are serious restrictions to the choice of the propagators. As we already saw, the apparently easiest choice is not invertible, and all other choices we tried proved to be invertible, but with pseudodifferential operators as their inverse instead of differential operators. We do not examine the question whether field equations with pseudodifferential operators are suitable choices within the general framework of quantum field theory. Instead we stick to differential operators as one is used to, the more so as the propagators we have defined above are completely sufficient for our purposes.
### 3.5. The free BRS theory
We examine in this section the free theory that includes vector bosons, spinors and ghosts. This is the theory that served as an example throughout the considerations above. The generators for the algebra $`𝒫`$ are in this model $`A_\mu ^a,(A_\mu ^a)^{(1,\nu )}`$ and $`(A_\mu ^a)^{(2,\nu \rho )}`$ for the Lie algebra valued vector bosons, $`u^a,\stackrel{~}{u}^a,(u^a)^{(1,\mu )}`$ and $`(\stackrel{~}{u}^a)^{(1,\mu )}`$ for the respective ghosts and anti-ghosts and $`\psi ^r`$ and $`\overline{\psi }^r`$ for the coloured spinors. The field operators that correspond to the generators $`A_\mu ^a,u^a,\stackrel{~}{u}_a,\psi ^r`$ and $`\overline{\psi }^r`$ are already constructed as operators in the Fock space $``$ with a common dense domain $`𝒟`$. As we already mentioned when we constructed the Fock space, the inner product $`,`$ is indefinite. To perform the BRS construction, we must define a BRS charge and a ghost charge and prove that the state space is positive.
At first we define the ghost current
$$k^\mu \stackrel{\text{def}}{=}i\underset{a}{}[(u^a)^{(0)}(\stackrel{~}{u}^a)^{(1,\mu )}(u^a)^{(1,\mu )}(\stackrel{~}{u}^a)^{(0)}]𝒫.$$
(3.98)
and the BRS current
$$j_B^\mu \stackrel{\text{def}}{=}\underset{a}{}[(u^a)^{(1,\mu )}(A_\nu ^a)^{(1,\nu )}(u^a)^{(0)}(A_\nu ^a)^{(2,\nu \mu )}]𝒫.$$
(3.99)
as elements of $`𝒫`$. Then their definitions as operators in the Fock space follow immediately as
$$k^\mu (x)=T\left(k^\mu \right)(x)\text{and}j_B^\mu (x)=T\left(j_B^\mu \right)(x).$$
(3.100)
Taking into account the field equations, we note that both operators are conserved,
$$_\mu ^xk^\mu (x)=_\mu ^xj_B^\mu (x)=0.$$
(3.101)
Now it is possible to define the corresponding charges, the ghost charge $`Q_c`$ and the BRS charge $`Q_B`$, as
$$Q_c\stackrel{\text{def}}{=}\underset{\lambda 0}{lim}d^4xh_\lambda (x)k^0(x)\text{and}Q_B\stackrel{\text{def}}{=}\underset{\lambda 0}{lim}d^4xh_\lambda (x)j_B^0(x).$$
(3.102)
Here $`h_\lambda 𝒟(M),\lambda \mathrm{I}\mathrm{R}^+\left\{0\right\}`$ is a test function that has the following structure:
$$\begin{array}{cc}\hfill h_\lambda (x)=\lambda h^t(\lambda x_0)b(\lambda 𝐱),& h^t𝒟(\mathrm{I}\mathrm{R})b𝒟(\mathrm{I}\mathrm{R}^3),\hfill \\ & 𝑑x_0h^t(x_0)=1,\hfill \end{array}$$
(3.103)
with $`b=1`$ on an open domain including the origin of $`\mathrm{I}\mathrm{R}^3`$. Due to a general argument of Requardt \[Req76\] the limit $`\lambda 0`$ exists and it is independent of the choice of $`h_\lambda `$. So the charges define well posed operators in the Fock space. The charges have no counterpart in the symbolic algebra, because the integrals would make no sense there.
The ghost transformation and the BRS transformation are (anti-) derivations on the algebra $`(𝒟)`$:
$$s_c(A)\stackrel{\text{def}}{=}[Q_c,A]_{},\text{and}s_0(A)\stackrel{\text{def}}{=}[Q_B,A]_{}A(𝒟).$$
(3.104)
The derivations give for the basic fields the following results:
$$\begin{array}{cc}& s_c(u^a(x))=u^a(x),s_c(\stackrel{~}{u}^a(x))=\stackrel{~}{u}^a(x),\hfill \\ & s_c(\phi _i(x))=0\text{otherwise},\hfill \\ & s_0(A_\mu ^a(x))=i_\mu ^xu^a(x),s_0(\stackrel{~}{u}^a(x))=i_x^\mu A_\mu ^a(x),\hfill \\ & s_0(\phi _i(x))=0\text{otherwise}.\hfill \end{array}$$
(3.105)
Finally we must prove that for the physical state space, defined as the state cohomology of $``$ w.r.t. the BRS charge $`Q_B`$ above, the positivity assumption holds. This has already been done by Kugo and Ojima \[KO79\], but we present here a modern version that is due to Razumov and Rybkin \[RR90\]. We collect here only the essential points of their proof.
At first they note that the entire space $`𝒟`$ can be decomposed as
$$=Q_B\left(Q_BQ_B^+\right)Q_B^+$$
(3.106)
with
$$\begin{array}{cc}& Q_B\left(Q_BQ_B^+\right)=\mathrm{ker}Q_B\hfill \\ \hfill \text{and}& \left(Q_BQ_B^+\right)Q_B^+=\mathrm{ker}Q_B^+.\hfill \end{array}$$
(3.107)
Then they propose an alternative definition of the physical (pre-) Hilbert space according to
$$_{\mathrm{phys}}=\left(Q_BQ_B^+\right)$$
(3.108)
or, which is the same,
$$_{\mathrm{phys}}=\mathrm{ker}\left(\{Q_B,Q_B^+\}_+\right).$$
(3.109)
This definition of the physical (pre-) Hilbert space deviates from the original one in the way that it selects from each equivalence class there exactly one representative. Now a direct calculation of the operator $`\{Q_B,Q_B^+\}_+`$ reveals
$$\{Q_B,Q_B^+\}_+=N_0+N_L+N_g$$
(3.110)
where $`N_0`$ is the number operator of scalar vector bosons introduced above, $`N_L`$ is the corresponding operator for the longitudinal vector bosons and $`N_g`$ the operator that counts the total number ghosts and anti-ghosts. Comparison with the definition of the Krein operator
$$J=(1)^{N_0}1\mathrm{l}J_g$$
(3.111)
reveals that $`J=1\mathrm{l}`$ on $`_{\mathrm{phys}}=\mathrm{ker}\left(\{Q_B,Q_B^+\}_+\right)`$. Therefore the inner product must be positive on $`_{\mathrm{phys}}`$ since the original scalar product was. It is not necessary to restrict the physical Hilbert space to the kernel of $`Q_c`$ since this Hilbert space is already contained in $`\mathrm{ker}N_g\mathrm{ker}Q_c`$.
The result ensuring positivity holds also for our definition of $`_{\mathrm{phys}}`$ as
$$_{\mathrm{phys}}=\overline{(\mathrm{ker}Q_B,𝒟)/(Q_B,𝒟)}^{},$$
(3.112)
since in this definition each equivalence class modulo $`(Q_B)`$ corresponds to exactly one element of $`\mathrm{ker}\left(\{Q_B,Q_B^+\}_+\right)`$, and the inner product does not depend on the choice of the representative within the equivalence class.
Then the algebra of observables is defined as usual,
$$𝒜_{\mathrm{ph}}\stackrel{\text{def}}{=}\left((\mathrm{ker}s,(𝒟))(\mathrm{ker}s_c,(𝒟))\right)/(s,(𝒟)).$$
(3.113)
As was pointed out by Dütsch and Fredenhagen, the algebra is faithfully represented in the physical Hilbert space.
## 4. Time ordered products and their normalization
In this chapter the construction of time ordered products, antichronological products and their respective properties are presented. Since the construction is in general not unique, normalization conditions are postulated that restrict the ambiguity. The construction of time ordered products is the central point for causal perturbation theory, which is presented in the next chapter. In particular this is the point where renormalization takes place in this framework.
The time ordering of $`n`$ arbitrary Wick polynomials $`W_1(x_1),\mathrm{},W_n(x_n)`$, $`W_i_1(𝒟)`$, can be done by the following prescription
$$\begin{array}{cc}\hfill T(W_1(x_1)\mathrm{}W_n(x_n))\stackrel{\text{def}}{=}& (1)^{f(\pi )+g(\pi )}\underset{\pi 𝒫_{\underset{¯}{n}}}{}\theta (x_{\pi (1)}^0x_{\pi (2)}^0)\mathrm{}\hfill \\ & \mathrm{}\theta (x_{\pi (n1)}^0x_{\pi (n)}^0)W_{\pi (1)}(x_{\pi (1)})\mathrm{}W_{\pi (n)}(x_{\pi (n)})\hfill \end{array}$$
(4.1)
if all the points $`x_i`$ are different. Here $`𝒫_{\underset{¯}{n}}`$ is the set of permutations of $`\underset{¯}{n}\stackrel{\text{def}}{=}\{1,\mathrm{},n\}`$, $`f(\pi )`$ is the number of transpositions in $`\pi 𝒫_{\underset{¯}{n}}`$ that involve arguments with an odd fermion number and $`g(\pi )`$ is the number of those that involve arguments with odd ghost number. $`\theta `$ is the Heaviside step function,
$$\theta (x)=\{\begin{array}{cc}1\hfill & \text{if }x^0>0\hfill \\ 0\hfill & \text{otherwise.}\hfill \end{array}$$
(4.2)
The crucial point is that this prescription is not defined for coinciding points, because the Wick polynomials $`W_i`$ are distributions that “do not like to be multiplied by discontinuous functions” \[Sto93\]. This is the origin of the ultraviolet divergences of quantum field theories. The prescription above gives, as it stands, well defined distributions only on a smaller space of test functions than $`𝒟(M^n)`$. This is the space of test functions in $`𝒟(M^n)`$ that vanish with all their derivatives if two or more of their spacetime arguments coincide. To form time ordered products these distributions on the smaller space of test functions must be extended to elements of $`_n(𝒟)`$.
The time ordering of $`n`$ arguments is usually regarded as a mapping of $`n`$ operator valued distributions in $`_1(𝒟)`$ to an operator valued distribution in $`_n(𝒟)`$. We however define the time ordering of $`n`$ arguments as a mapping of $`n`$ polynomials in $`𝒫`$ to an operator valued distribution in $`_n(𝒟)`$. As already mentioned this has the advantage that the normalization conditions can be formulated also for derivated fields. Beside that technical point the extension of the distributions follows the method of Epstein and Glaser \[EG73\]. The extension exists always but is in general not unique. Therefore for each combination of arguments one element in $`_n(𝒟)`$ must be chosen as the time ordered product of these arguments. This choice is called the normalization of that time ordered product according to Scharf \[Sch95\].
The normalization conditions implement various properties of the time ordered products that are desired from the physical point of view. The postulation of the normalization conditions restricts the number of possible normalizations, but the extension is in general still not unique.
This chapter is organized as follows: The first section presents the properties of time ordered products that are required for their construction. In the next section this construction is performed. Antichronological products are defined in the third section. The chapter concludes with a section in which the normalization conditions are formulated.
### 4.1. Properties of time ordered products
The construction of time ordered products proceeds by induction. The time ordered products of a number of arguments are built out of the time ordered products with fewer arguments. This construction works only if the time ordered products with fewer arguments have certain properties. These properties are presented here. They are
P1 (Well posedness): The time ordering operator for $`n`$ arguments, $`T_n`$, is a multi linear mapping of $`n`$ polynomials in $`𝒫`$ to the operator valued distributions of order $`n`$ on the dense subspace $`𝒟`$:
$$T_n:\underset{n\mathrm{times}}{\underset{}{𝒫\times \mathrm{}\times 𝒫}}_n(𝒟).$$
(4.3)
If the arguments are explicitely given, the index $`n`$ indicating the number of arguments will be omitted.
From the physical interpretation of time ordering we would expect that the time ordering operator must have at least two arguments, for otherwise there is nothing to be put in order. But it turns out to be useful to extend the mapping defined above formally also to the cases $`n=0`$ and $`n=1`$. This is achieved by the following definitions:
$$T_0\stackrel{\text{def}}{=}1\mathrm{l},1\mathrm{l}(𝒟)$$
(4.4)
and
$$T_1(W)(x)\stackrel{\text{def}}{=}T(W)(x),W𝒫.$$
(4.5)
Here $`T`$ on the right hand side is the linear representation defined in the last chapter. The operator valued distributions obtained by the time ordering are called time ordered products or $`T`$\- products. The time ordered product of the polynomials $`W_1,\mathrm{},W_n`$ is written as
$$T(W_1,\mathrm{},W_n)(x_1,\mathrm{},x_n).$$
(4.6)
The definition above implies that the tensor product of two time ordered products with $`m`$ and $`n`$ arguments is a well defined operator valued distribution in $`_{m+n}(𝒟)`$. Arguments that are multiples of the identity can be removed according to
$$T(W_1,\mathrm{},W_n,a1\mathrm{l})(x_1,\mathrm{},x_n,y)=aT(W_1,\mathrm{},W_n)(x_1,\mathrm{},x_n)a\text{ }\mathrm{C}.$$
(4.7)
P2 (Graded symmetry): Time ordered products are totally graded symmetric under permutations of their indices. That means
$$\begin{array}{cc}& T(W_{\pi (1)},\mathrm{},W_{\pi (n)})(x_{\pi (1)},\mathrm{},x_{\pi (n)})\hfill \\ & =(1)^{f(\pi )+g(\pi )}T(W_1,\mathrm{},W_n)(x_1,\mathrm{},x_n)\pi 𝒫_{\underset{¯}{n}},\hfill \end{array}$$
(4.8)
where the integers $`f(\pi )`$, $`g(\pi )`$ were defined in eqn. (4.1).
P3 (Causality): Time ordered products are causal, that means they fulfill eqn. (4.1) for non coinciding points. Even more, outside the total diagonal $`\mathrm{Diag}_n`$ (see below) the time ordered product of $`n`$ arguments is completely determined by those that have fewer arguments. The total diagonal $`\mathrm{Diag}_nM^n`$ is the set where all points coincide:
$$\mathrm{Diag}_n=\{(x_1,\mathrm{},x_n)M^n:x_1=\mathrm{}=x_n\}.$$
(4.9)
If not all points $`x_i`$ coincide there exists a spacelike surface $`\mathrm{\Sigma }M`$ that separates the points $`X=\{x_1,\mathrm{},x_n\}`$ into a future subset $`Z`$ and a past subset $`Z^c=XZ`$ such that
$$\mathrm{\Sigma }X=\mathrm{},Z(\mathrm{\Sigma }+\overline{V}_+),Y(\mathrm{\Sigma }+\overline{V}_{}).$$
(4.10)
This situation will be denoted as $`ZZ^c`$. Furthermore we introduce the abbreviation
$$T\left(W_Z\right)(x_Z)\stackrel{\text{def}}{=}T(W_1,\mathrm{},W_k)(x_1,\mathrm{}x_k)\text{if }Z=\{x_1,\mathrm{}x_k\}.$$
(4.11)
Causality means that the time ordered product $`T\left(W_X\right)(x_X)`$ is required to satisfy causal factorization:
$$T\left(W_X\right)(x_X)=T\left(W_Z\right)(x_Z)T\left(W_{Z^c}\right)(x_{Z^c})\text{if }ZZ^c.$$
(4.12)
It provides a recursive definition of the time ordered products up to the diagonal $`\mathrm{Diag}_n`$. There the separation into future and past subsets is impossible and therefore no causal factorization exists. Validity of causal factorization for every number of arguments implies that spacelike separated time ordered products (anti-) commute<sup>15</sup><sup>15</sup>15The notation $`ZZ^c`$ means that $`Z`$ and $`Z^c`$ are spacelike separated, i.e. $`ZZ^c`$ and $`Z^cZ`$.:
$$[T\left(W_Z\right)(x_Z),T\left(W_{Z^c}\right)(x_{Z^c})]_{}=0\text{if }ZZ^c.$$
(4.13)
P4 (Translational invariance): Time ordered products are translationally invariant, that means that for every $`aM`$ the following equation holds:
$$\begin{array}{cc}\hfill \left(\mathrm{Ad}U(p)\right)& T(W_1,\mathrm{},W_n)(x_1,\mathrm{},x_n)=\hfill \\ & =T(W_1,\mathrm{},W_n)(x_1a,\mathrm{},x_na)p=(a,1\mathrm{l})𝒫_+^{}.\hfill \end{array}$$
(4.14)
Here $`U`$ is the representation of the Poincaré group in the Fock space introduced in the last chapter.
### 4.2. Inductive construction of time ordered products
In this section the inductive construction of the time ordered products is outlined. It goes back to Epstein and Glaser \[EG73\]. We use a formulation of their procedure proposed by Stora \[Sto93\] and recently elaborated by Brunetti and Fredenhagen in \[BF99\]. This section will not contain the proofs of the theorems. For them we refer to the latter article.
Formally the time ordering is also defined for a single argument by the linear representation $`T`$. The latter is uniquely defined for all $`W𝒫`$. This will serve as a starting point for the induction. Obviously the representation satisfies properties P1 \- P4.
We suppose that all $`T`$-products for up to $`n1`$ arguments are already constructed and satisfy properties P1 \- P4. Due to property P3 the time ordered products for $`n`$ arguments are therefore completely determined on $`M^n\mathrm{Diag}_n`$, i.e. for all test functions in $`𝒟(M^n\mathrm{Diag}_n)`$. To construct the distributions off the diagonal we introduce at first a partition of $`M^n\mathrm{Diag}_n`$ into the spaces
$$\begin{array}{cc}\hfill \mathrm{}_Z\stackrel{\text{def}}{=}\{(x_1,\mathrm{},x_n)M^n:x_i(x_j+\overline{V}_{}),iZ,jZ^c\}& \\ \hfill \text{for every }Z\mathrm{},ZX& .\hfill \end{array}$$
(4.15)
It is easy to see (and has been proven in \[BF99, Lemma 4.1\]) that
$$\underset{\stackrel{Z\mathrm{},}{ZX}}{}\mathrm{}_Z=M^n\mathrm{Diag}_n.$$
(4.16)
Furthermore we define
$$T_Z(W_X)(x_X)\stackrel{\text{def}}{=}\{\begin{array}{cc}T\left(W_Z\right)(x_Z)T\left(W_{Z^c}\right)(x_{Z^c})\text{if }(x_1,\mathrm{},x_n)\mathrm{}_Z,\hfill & \\ 0\text{otherwise}.\hfill & \end{array}$$
(4.17)
$`T_Z(W_X)(x_X)`$ is a well defined operator valued distribution in $`_n(𝒟)`$. Finally we choose an arbitrary locally finite $`C^{\mathrm{}}`$-partition of unity for $`M^n\mathrm{Diag}_n`$,
$$\begin{array}{cc}\hfill \left\{f_Z\right\}:& \underset{Z}{}f_Z=1\text{ on }M^n\mathrm{Diag}_n,\hfill \\ & f_Z\mathrm{}_Z,f_ZC^{\mathrm{}}(M^n\mathrm{Diag}_n).\hfill \end{array}$$
(4.18)
The restriction of $`T\left(W_X\right)(x_X)`$ to $`M^n\mathrm{Diag}_n`$<sup>16</sup><sup>16</sup>16That means $`T^0\left(W_X\right)(x_X):𝒟(M^n\mathrm{Diag}_n)(𝒟)`$ can now be defined as
$$T^0\left(W_X\right)(x_X)\stackrel{\text{def}}{=}\underset{Z}{}f_Z(x_X)T_Z(W_X)(x_X).$$
(4.19)
This definition does not depend on the choice of $`\left\{f_Z\right\}`$ because we assumed that eqns. (4.12) and (4.13) hold for the $`T`$-products with fewer arguments. This makes the $`T^0`$-products well defined operator valued distributions on test functions in $`𝒟(M^n\mathrm{Diag}_n)`$ that satisfy the properties P1 \- P4. For the proofs see \[BF99\].
For the construction of the time ordered products with $`n`$ arguments the $`T^0`$-products must be extended to the diagonal. They are linear combinations of products of numerical distributions $`t^{\mathrm{\hspace{0.17em}0}}`$ with Wick products $`:W_1(x_1)\mathrm{}W_n(x_n):`$, where the $`W_i(x_i)`$ are Wick monomials in $`_1(𝒟)`$. It is not trivial that these products exist, because distributions are multiplied at the same spacetime point, but it was shown by Epstein and Glaser that translational invariance implies that this product is indeed well defined — this result is referred to as “Theorem 0” in \[EG73, p. 229\]. From the modern point of view the product exists because the wave front sets of the distributions do not linearly combine to zero in the cotangent spaces, see \[BF99\].
For the extension of the operator valued distributions to the diagonal it suffices to extend each numerical distribution $`t^{\mathrm{\hspace{0.17em}0}}`$ and to prove that the resulting product is well defined. The latter is no problem here because the “Theorem 0” applies also to the extended distributions.
Translational invariance (P4) implies that the numerical distributions $`t^{\mathrm{\hspace{0.17em}0}}`$ depend only on the relative coordinates $`(y_1,\mathrm{},y_{n1})\stackrel{\text{def}}{=}(x_1x_n,\mathrm{},x_{n1}x_n)`$ such that $`\mathrm{Diag}_n`$ is the origin in the space of the $`y`$’s. This allows us to give a further restriction to the extension of the numerical distributions to the diagonal: Each distribution $`t^{\mathrm{\hspace{0.17em}0}}`$ is regarded as a distribution in the space of relative coordinates. Then the scaling degree of the extended distribution $`t`$ must not exceed that of the original distribution $`t^{\mathrm{\hspace{0.17em}0}}`$ in relative coordinates. Brunetti and Fredenhagen \[BF99\] prove that such an extension always exists as a well defined distribution for test functions in $`𝒟(M^{n1})`$ — or in $`𝒟(M^n)`$ if one returns to the original coordinates. It is unique only if the original distribution has a scaling degree $`(t^{\mathrm{\hspace{0.17em}0}})`$ that satisfies the following inequality:
$$(t^{\mathrm{\hspace{0.17em}0}})<(n1)\times d$$
(4.20)
where $`n1`$ is the number of relative coordinates and $`d`$ the spacetime dimension. This can be seen as follows: The distribution $`t`$ is already determined up to the diagonal $`\mathrm{Diag}_n`$. In other words, two extensions may differ only by a delta distribution with support at the origin of the relative coordinates or by a derivative of it. If the scaling degree of $`t^{\mathrm{\hspace{0.17em}0}}`$ satisfies the inequality above, it is not possible to add a delta distribution or a derivative of it without violating the restriction on the scaling degree. Therefore the solution is unique then. In general the inequality does not hold and the extension is therefore ambiguous, corresponding to the freedom of finite renormalization in other renormalization procedures.
### 4.3. Antichronological products
In this section we define antichronological products. This definition can be given recursively as $`\overline{T}_0=1\mathrm{l}`$ and<sup>17</sup><sup>17</sup>17The notation is the same as in eqn. (4.11)
$$\begin{array}{cc}& \overline{T}\left(W_X\right)(x_X)\stackrel{\text{def}}{=}\hfill \\ & =\underset{YX,Y\mathrm{}}{}(1)^{\left|Y\right|}T\left(W_Y\right)(x_Y)\overline{T}\left(W_{Y^c}\right)(x_{Y^c})\hfill \\ & =\underset{YX,YX}{}(1)^{\left|Y^c\right|}\overline{T}\left(W_Y\right)(x_Y)T\left(W_{Y^c}\right)(x_{Y^c})\hfill \end{array}$$
(4.21)
for $`n1`$. Here possible signs that come from changes in the order of the arguments are neglected for simplicity. They can be easily recovered using P2, which holds for the $`\overline{T}`$-products, too (see below).
Iterating the recursive definition above one finds the following explicit expression for the $`\overline{T}`$-products:
$$\overline{T}\left(W_X\right)(x_X)=\underset{P}{}(1)^{\left|P\right|+\left|X\right|}\underset{pP}{}T\left(W_p\right)(x_p).$$
(4.22)
Here the sum runs over all partitions $`P`$ of $`X`$ into $`\left|P\right|`$ nonempty subsets. With this definition the antichronological products become for non coinciding points, i.e. $`x_ix_jij`$,
$$\begin{array}{cc}\hfill \overline{T}\left(W_X\right)(x_X)=& \underset{\pi 𝒫_{\underset{¯}{n}}}{}\theta (x_{\pi (1)}^0x_{\pi (2)}^0)\mathrm{}\hfill \\ & \mathrm{}\theta (x_{\pi (n1)}^0x_{\pi (n)}^0)T\left(W_{\pi (n)}\right)(x_{\pi (n)})\mathrm{}T\left(W_{\pi (1)}\right)(x_{\pi (1)}).\hfill \end{array}$$
(4.23)
The antichronological products satisfy properties P1, P2 and P4. Property P3 holds for them in the reverse order. That means that under the same conditions and with the same notation as in P3 antichronological products satisfy
$$\overline{T}\left(W_X\right)(x_X)=T\left(W_{Z^c}\right)(x_{Z^c})T\left(W_Z\right)(x_Z),ZZ^c$$
(P3’)
justifying their name since they are defined like the time ordered products but with the opposite order.
### 4.4. Normalization conditions
In this section we formulate the normalization conditions that restrict the ambiguity in the extension of the $`\overline{T}`$-products to the diagonal. They implement Poincaré covariance (N1) and unitarity (N2). They define the time ordered products up to a $`\mathrm{C}`$-number distribution (N3) and determine them uniquely if at least one argument is a generator from $`𝒢`$ (N4). Finally they determine Ward identities for the ghost current (N5) and the BRS current (N6). It is proven that the conditions (N1) - (N5) have common solutions. For condition (N6) this must be done for the individual models.
The first normalization condition establishes Poincaré covariance w.r.t. the representation $`U`$ of the Poincaré group $`𝒫_+^{}`$ introduced in chapter (3). It reads
$$\begin{array}{cc}& \left(\mathrm{Ad}U(p)\right)T(W_1,\mathrm{},W_n)(x_1,\mathrm{},x_n)=\hfill \\ & =T(_\mathrm{\Lambda }(W_1),\mathrm{},_\mathrm{\Lambda }(W_n))(\mathrm{\Lambda }^1x_1a,\mathrm{},\mathrm{\Lambda }^1x_na)\hfill \end{array}$$
(N1)
for every $`p=(a,\mathrm{\Lambda })𝒫_+^{}`$ and all monomials $`W_i𝒫`$. Here $`_\mathrm{\Lambda }`$ is the representation of the Lorentz group on $`𝒫`$ introduced in section (3.1). Property $`\mathrm{𝐏𝟒}`$ is in view of (N1) only the special case with $`p=(a,1\mathrm{l})`$.
Popineau and Stora \[PS82\] have proven that this condition has always a solution, but their article is unfortunately not published. So we refer the reader to Scharf \[Sch95, p. 282\] for the proof. Recently Prange, Bresser and Pinter, \[BPP99\] and \[Pra99\], have found even a general construction prescription for covariant normalizations.
The second normalization condition establishes pseudo-unitarity by means of
$$T(W_1,\mathrm{},W_n)^{}(x_1,\mathrm{},x_n)=\overline{T}(W_n^{},\mathrm{},W_1^{})(x_n,\mathrm{},x_1)W_i𝒫,$$
(N2)
where the -involution on the left hand side is the Krein adjoint on $`(𝒟)`$, while the -involution on the right hand side is the adjoint operation in $`𝒫`$ defined in section (3.1). Note that the order of the arguments is reversed. It can of be put into the original order by means of P2.
It was already shown by Epstein and Glaser \[EG73\] that eqn. (N2) can always be accomplished. Their argument and the compatibility of (N2) with (N1) can be easily understood: Suppose, (N2) holds for all integers $`m<n`$ simultaneously with eqn. (N1). Then for every normalization $`T^{}=T(W_1,\mathrm{},W_n)`$ that is compatible with eqn. (N1) the distribution $`T=\frac{1}{2}(T^{}+T_{}^{}{}_{}{}^{})`$ satisfies eqn. (N2) and will also be an extension of $`T^0`$ because (N2) holds for the $`T^0`$-products by induction. It will automatically be a solution of eqn. (N1) since the representation $`U`$ was chosen to be pseudo-unitary, i.e. $`U(p)^{}=U(p)^1`$.
To formulate the third normalization condition we remind the reader of the commutator function $`\mathrm{\Delta }_{jk}(x)`$, eqn. (3.60) in section \[3.3\]. The normalization condition reads:
$$\begin{array}{cc}& [T(W_1,\mathrm{},W_n)(x_1,\mathrm{},x_n),\phi _i(y)]_{}=\hfill \\ & =i\underset{k=1}{\overset{n}{}}\underset{j}{}\mathrm{\Delta }_{ij}(x_ky)T(W_1,\mathrm{},\frac{W_k}{\phi _j},\mathrm{},W_n)(x_1,\mathrm{},x_n),\hfill \end{array}$$
(N3)
for every $`W_i𝒫,\phi _i(y)=T\left(\phi _i\right)(y),\phi _i𝒢`$. The second sum runs over all generators in $`𝒢`$, not only the basic generators.
Since an element of $`(𝒟)`$ is a multiple of the identity if it (anti-) commutes with all the $`\phi _i(y)`$ — see eqn. (3.56) in section (3.2) —, this condition determines the time ordered products uniquely up to a $`\mathrm{C}`$-number, provided the time ordered products that involve the sub monomials are known. This can be explicitly seen in an equivalent equation, the causal Wick expansion
$$\begin{array}{cc}\hfill T(W_1,\mathrm{},W_n)(x_1,\mathrm{},x_n)=\underset{\gamma _1,\mathrm{},\gamma _n}{}\omega _0\left(T(W_1^{(\gamma _1)},\mathrm{},W_n^{(\gamma _n)})(x_1,\mathrm{},x_n)\right)& \\ \hfill \times \frac{:\phi ^{\gamma _1}(x_1)\mathrm{}\phi ^{\gamma _n}(x_n):}{\gamma _1!\mathrm{}\gamma _n!}& .\hfill \end{array}$$
(4.24)
Here the $`\gamma _i\mathrm{I}\mathrm{N}^r`$ are multi indices, vectors with one entry for each of the $`r`$ generators in $`𝒢`$, i.e.
$$\gamma _i=((\gamma _i)_1,\mathrm{},(\gamma _i)_r)\mathrm{I}\mathrm{N}^r$$
(4.25)
The $`W^{(\gamma _i)}`$ are derivatives,
$$W^{(\gamma _i)}\stackrel{\text{def}}{=}\frac{^{\left|\gamma _i\right|}W}{^{(\gamma _i)_1}\phi _1\mathrm{}^{(\gamma _i)_r}\phi _r},$$
(4.26)
where $`\left|\gamma _i\right|=_{k=1}^r(\gamma _i)_k`$. The $`\phi ^{\gamma _i}`$ are defined as
$$\phi ^{\gamma _i}(x)\stackrel{\text{def}}{=}T\left(\underset{k=1}{\overset{r}{}}\phi _k^{(\gamma _i)_k}\right)(x).$$
(4.27)
Finally
$$(\gamma _i)!\stackrel{\text{def}}{=}\underset{k=1}{\overset{r}{}}(\gamma _i)_k!.$$
(4.28)
It is shown in the appendix, section A.1, that the causal Wick expansion is indeed equivalent with (N3). Compatibility with eqn. (N1) is easily verified since (N3) respects the Poincaré transformation properties. With the same construction as after eqn. (N2) one can show that for every common solution of (N1) and (N3) a normalization can be constructed that is also a solution of (N2).
In particular in the formulation (4.24) of (N3) it is immediately clear that only $`\omega _0\left(T(W_1,\mathrm{},W_n)(x_1,\mathrm{},x_n)\right)`$ — the term with $`\gamma _1=\mathrm{}=\gamma _n=0`$ in (4.24) — is left open to be normalized, since all other terms are determined by the time ordered products for the sub monomials. These distributions correspond to the vacuum diagrams of the respective time ordered product in the Feynman graph picture. So condition (N3) has the consequence that only vacuum diagrams need to be (re-) normalized, a fact that is well known from other renormalization procedures.
The fourth normalization condition is a differential equation that uniquely determines time ordered products with at least one generator $`\phi _i𝒢`$ among its arguments. This assertion holds under the assumption that the time ordered products for fewer arguments are already known. The condition reads:
$$\begin{array}{cc}& \underset{j}{}D_{ij}^yT(W_1,\mathrm{},W_n,\phi _j)(x_1,\mathrm{},x_n,y)=\hfill \\ & =i\underset{k=1}{\overset{n}{}}T(W_1,\mathrm{},\frac{W_k}{\phi _i},\mathrm{},W_n)(x_1,\mathrm{},x_n)\delta (x_ky),\hfill \end{array}$$
(N4)
where $`W_i𝒫,\phi _j𝒢`$. It is proven in the appendix, section A.2, that condition (N4) has common solutions with condition (N3). Compatibility with condition (N1) is again immediate since (N4) is Poincaré covariant. A solution of (N1), (N3) and (N4) that satisfies also (N2) can be found by the same procedure as above.
In the chapter concerning the interacting theory we will see that eqn. (N4) already implies the interacting field equations.
Condition (N4) possesses an alternative formulation, like (N3). Its integrated version reads
$$\begin{array}{cc}& T(W_1,\mathrm{},W_n,\phi _i)(x_1,\mathrm{},x_n,y)=\hfill \\ & =i\underset{k=1}{\overset{n}{}}\underset{j}{}\mathrm{\Delta }_{ij}^F(yx_k)T(W_1,\mathrm{},\frac{W_k}{\phi _j},\mathrm{},W_n)(x_1,\mathrm{},x_n)\hfill \\ & +\underset{\gamma _1\mathrm{}\gamma _n}{}\omega _0\left(T(W_1^{(\gamma _1)},\mathrm{},W_n^{(\gamma _n)})(x_1,\mathrm{},x_n)\right)\frac{:\phi ^{\gamma _1}(x_1)\mathrm{}\phi ^{\gamma _n}(x_n)\phi _i(y):}{\gamma _1!\mathrm{}\gamma _n!}.\hfill \end{array}$$
(4.29)
The sum over $`j`$ runs again over all generators including the higher ones.
This formulation shows explicitely that with eqn. (N4) the time ordered products with at least one generator among its arguments are already determined. In appendix (A.2) the equivalence of the two formulations is proven.
Eqn. (N4) uniquely fixes the Feynman propagators for derivated fields. These in turn determine all tree level diagrams. Comparing (N4) with results from other renormalization procedures shows an important difference between the causal approach and other approaches: The definition of the propagators for the derivated fields differ between the causal approach and other approaches. Therefore also the Green’s functions at tree level are different. The difference between the conventional propagators and our prescription is labelled by the normalization constants $`C_{\phi _i,k}`$. Only if all these constants are set to zero the difference disappears. But we saw already that the propagators are then no longer invertible. For example, in the conventional renormalization procedures we have
$$\omega _0\left(T(_x^\mu A_\nu (x),_y^\nu A_\rho (y))\right)=i_\rho ^x_x^\mu D^F(xy),$$
(4.30)
while the corresponding propagator in our causal theory reads
$$\omega _0\left(T((A_\nu )^{1,\mu },(A_\rho )^{1,\nu })(x,y)\right)=i_\rho ^x_x^\mu D^F(xy)iC_{A,1}\delta _\rho ^\mu \delta (xy).$$
(4.31)
Now we come to the Ward identities for the ghost current. This is a normalization condition for time ordered products that contain a ghost current $`k^\mu `$ — see section (3.5) — as an argument. It reads
$$\begin{array}{cc}& _\mu ^yT(W_1,\mathrm{},W_n,k^\mu )(x_1,\mathrm{},x_n,y)=\hfill \\ & =\underset{k=1}{\overset{n}{}}g(W_k)\delta (yx_k)T(W_1,\mathrm{},W_n)(x_1,\mathrm{},x_n)\hfill \end{array}$$
(N5)
for all monomials $`W_i𝒫`$. It holds if none of the arguments contains a generator $`(u^a)^{(\alpha )}`$ or $`(\stackrel{~}{u}^a)^{(\alpha )}`$ with $`\left|\alpha \right|2`$.
The proof that this normalization condition has common solutions with the conditions (N1) - (N4) is given in appendix (B.1). For a technical reason that will be explained there this normalization condition can only be proven for arguments $`W_i`$ that do not contain $`k^\mu `$ as a sub monomial — in particular $`k^\mu `$ itself is excluded. In the examples where (N5) is applied in the following chapters this limitation will not be relevant.
We state here one particular fact that will come out in the proof: There exists exactly one choice for the normalization constant $`C_{u,1}`$ such that condition (N5) has common solutions with (N1) - (N4). This choice is $`C_{u,1}=1`$.
Following Dütsch and Fredenhagen \[DF99\] who made the calculation for the Ward identities for the electric current (see below) we prove in appendix (B.1) that there exists an integrated version of (N5), namely
$$\begin{array}{cc}& s_cT(W_1,\mathrm{},W_n)(x_1,\mathrm{},x_n)=\hfill \\ & =\left(\underset{k=1}{\overset{n}{}}g(W_k)\right)T(W_1,\mathrm{},W_n)(x_1,\mathrm{},x_n).\hfill \end{array}$$
(4.32)
So as a consequence of (N5) the ghost number of a time ordered product is simply the sum of the ghost numbers of its arguments.
Eqn. (N5) and (4.32) are equivalent in the following sense: If (N5) holds then (4.32) is automatically valid, too. If (4.32) holds, then a normalization can be found that is compatible with (N5). For details see appendix (B.1).
Dütsch and Fredenhagen \[DF99\] proved an analogous Ward identity for the electric current $`j_{\mathrm{el}}^\mu =\overline{\psi }\gamma ^\mu \psi `$. Here $`\psi `$ and $`\overline{\psi }`$ are the electron and the positron field, respectively. Their Ward identity reads in our language
$$\begin{array}{cc}& _\mu ^yT(W_1,\mathrm{},W_n,j_{\mathrm{el}}^\mu )(x_1,\mathrm{},x_n,y)=\hfill \\ & =i\left(\underset{k=1}{\overset{n}{}}f(W_k)\delta (yx_k)\right)T(W_1,\mathrm{},W_n)(x_1,\mathrm{},x_n).\hfill \end{array}$$
(N5’)
It holds if the monomials $`W_i`$ do not contain generators $`\psi ^{(\alpha )}`$ or $`\overline{\psi }^{(\alpha )}`$ with $`\left|\alpha \right|1`$. The existence of common solutions of (N5’) with the other normalization conditions can be proven along the same lines as for the ghost current Ward identities, provided that either none of the arguments $`W_i`$ contains $`j_{\mathrm{el}}^\mu `$ as a sub monomial or that all the $`W_i`$ are the QED Lagrangian $`_{QED}=A_\mu j_{\mathrm{el}}^\mu `$ or sub monomials of it.
To formulate the Ward identity for the BRS current we anticipate here a condition for the Lagrangian that will be illuminated more closely in section (5.4). In QED and Yang-Mills theory there exist so called $`Q(n)`$-vertices for the Lagrangians. These are polynomials $`_1^\mu ,_2^{\mu \rho },\mathrm{}𝒫`$ totally antisymmetric in their Lorentz-indices for which the following identities hold:
$$s_cT\left(_i^{\mu _1,\mathrm{},\mu _i}\right)(x)=i_\rho ^xT\left(_{i+1}^{\mu _1,\mathrm{},\mu _i,\rho }\right)(x).$$
(4.33)
We admit only polynomials $``$ as Lagrangians if there exist such $`Q(n)`$-vertices and in addition so called $`R(n)`$-vertices $`_1,_2,\mathrm{}`$ that are polynomials in $`𝒫`$ which satisfy the following condition: There exists a normalization of $`T(_i^{\mu _1,\mathrm{},\mu _i},j^\mu )`$ that is compatible with the normalization conditions (N1) - (N4) and for which the equation
$$\begin{array}{cc}\hfill _\mu ^yT(_i^{\mu _1,\mathrm{},\mu _i},j^\mu )(x,y)=& i_\nu ^x\left(\delta (xy)_{i+1}^{\mu _1,\mathrm{},\mu _i,\nu }(x)\right)\hfill \\ & +i\left(_\nu ^x\delta (xy)\right)_{i+1}^{\mu _1,\mathrm{},\mu _i,\nu }(x)\hfill \end{array}$$
(4.34)
holds. The series of equations terminates at a certain point i.e. there exists an $`m\mathrm{I}\mathrm{N}`$ with $`_m=0,_m=0`$. This is the condition (C4) in section (5.4). The $`R(n)`$-vertices are totally antisymmetric in their Lorentz indices, too.
With the notion of $`Q(n)`$-vertices and the $`R(n)`$-vertices we can give the next normalization condition, the Ward identities for the BRS current:
$$\begin{array}{cc}& _\mu ^yT(_{i_1},\mathrm{},_{i_n},j^\mu )(x_1,\mathrm{},x_n,y)=\hfill \\ & =i\underset{k=1}{\overset{n}{}}_\nu ^k\left(\delta (yx_k)T(_{i_1},\mathrm{},_{i_k+1}^\nu ,\mathrm{},_{i_n})(x_1,\mathrm{},x_n)\right)\hfill \\ & +i\underset{k=1}{\overset{n}{}}\left(_\nu ^k\delta (yx_k)\right)T(_{i_1},\mathrm{},_{i_k+1}^\nu ,\mathrm{},_{i_n})(x_1,\mathrm{},x_n)\hfill \end{array}$$
(N6)
where $`i\mathrm{I}\mathrm{N}`$ and we define $`_0=`$.
The same calculation leading to eqn. (4.32) can also be applied to condition (N6) and gives the generalized operator gauge invariance
$$\begin{array}{cc}& s_0T(_{i_1},\mathrm{},_{i_n})(x_1,\mathrm{},x_n)=\hfill \\ & =i\underset{k=1}{\overset{n}{}}_\nu ^kT(_{i_1},\mathrm{},_{i_k+1}^\nu ,\mathrm{},_{i_n})(x_1,\mathrm{},x_n).\hfill \end{array}$$
(4.35)
Dütsch, Hurth, Krahe and Scharf, \[DHKS94a\] \- \[DHS95b\], found that for eqn. (4.35) to hold in Yang-Mills theory for two arguments the normalization constant $`C_{A,1}`$ in eqn. (3.91) must be $`C_{A,1}=\frac{1}{2}`$.
Unfortunately there exists no general proof that condition (N6) can always be accomplished or that it is compatible with (N3) and (N4)<sup>18</sup><sup>18</sup>18At a first sight it may seem that (N4) has nothing to do with (N6) since there is no generator in the time ordered products whose normalization (N6) determines. The point is that compatibility with (N3) requires a set of relations among which are also some that involve time ordered products that contain a generator. Then (N4) could fix their normalization in a way that compatibility between (N3) and (N6) is inhibited. In this sense we think that (N4) and (N6) shall be compatible.. But we can show that the generalized operator gauge invariance together with (N5) is already sufficient for (N6). For the construction of solutions of (N6) under the assumption that generalized operator gauge invariance holds see appendix (B.2).
The proof that either the eqn. (N6) or eqn. (4.35) have common solutions with the other normalization conditions must be done in individual models.
As far as we know QED is the only example where this is done — for the proof see section (7.1). The existence of solutions for eqn. (4.35) is in QED a direct consequence of the existence of solutions for the electric current Ward identities (N5’).
In Yang-Mills theories the solutions for eqn. (4.35) can be explicitely given in first order, see section (7.2). A detailed study of (4.35) with $`i_1=\mathrm{}=i_n=0`$ for Yang-Mills theory without matter fields can be found in \[DHKS94a\] \- \[DHS95b\] — this equation is called operator gauge invariance. This result has been generalized to Yang-Mills theory with matter fields by Dütsch \[Düt96\]. They come to the result that operator gauge invariance holds in that theory provided a weak assumption concerning the infrared behaviour is satisfied. Loosely speaking the infrared behaviour must not be too bad. It is usually assumed that this assumption is satisfied, otherwise not even off shell Green’s functions would exist.
It should be possible to prove the generalized version of operator gauge invariance under the same assumption and along the same lines as in their calculation, but this has not been done up to now — and it is probably a long winded work, the original calculation took a series of four articles.
Another promising strategy to prove generalized operator gauge invariance is to translate the results of algebraic renormalization \[PS95\] to causal perturbation theory. The descent equations can be viewed as the generalized operator gauge invariance version of that framework. It has been proven in \[PS95\] that they can be accomplished for Yang-Mills theories. Unlike our causal approach algebraic renormalization is a loop expansion, i.e. an expansion in the parameter $`\mathrm{}`$ and not in the coupling constant. Furthermore it is a functional approach, in contrast to causal perturbation theory which is an operator approach. So in order to make the results cited above available to the causal theory some translational work has to be performed. This has not been done up to now.
For the equations (N6) with $`i_n5`$ the compatibility of normalization conditions is easy to prove: These $`T`$-products comply automatically with (N3) and (N4) since their ghost number, which is the minimal number of field operators in the Wick products in the causal Wick expansion, exceeds the spacetime dimension, so we are in the situation of the inequality (4.20) and therefore the extension is unique and complies with (N3), (N4) and (N5).
We have stated altogether six normalization conditions for the time ordered products (where for the last it remains open whether it can always be accomplished). One could ask whether these conditions suffice to make the extension of the $`T^0`$-products to the diagonal unique. Unfortunately this is not true. There remains a certain ambiguity, even though calculations in first order show that the normalization conditions restrict the freedom of the extensions severely — in fact there are many examples where the above conditions suffice to make the extension unique.
The decisive point is that the normalization conditions suffice to prove a lot of relations in the interacting theory like field equations, nilpotency of the interacting BRS charge and others, notwithstanding the remaining ambiguity.
Another interesting feature of these normalization conditions is that a subsequent enlargement of the algebra $`𝒫`$ — by the introduction of new basic fields or by inclusion of generators for higher derivatives of the basic fields than before — does not change the normalization of the time ordered products with arguments in the original, smaller algebra. Moreover these normalizations do not depend on the model with regard to which they are considered. For example there are certain time ordered products that occur both in QED and in Yang-Mills theory, but due to our construction their normalization is the same in both cases, provided the normalization constants $`C_{\phi _i,k}`$ are chosen equal. This is of course a consequence of the fact that the normalizations are completely independent of the Lagrangian. The latter is in this context a polynomial in $`𝒫`$ not outstanding from the others. So the idea behind the whole construction is to determine all time ordered products a priori, store them in a big library and fetch them if they are needed for a certain calculation. The remaining ambiguity of the time ordered products is certainly a handicap. Ambiguous time ordered products should be laid down in this library with an endorsement that they are ambiguous and what the allowed normalizations are.
## 5. Local causal perturbation theory
This chapter is devoted to the formulation of local causal perturbation theory. It will establish the connection of time ordered products with interacting quantum field theory. In the framework of causal perturbation theory the S-matrix and the interacting field operators are defined in terms of time ordered products, see below.
As usual the interaction will be defined by the S-matrix. But as we investigate local theories, there will be no interpretation of the S-matrix available as an operator mapping in-states onto out-states. These asymptotic states are a global concept that looses its meaning in a local framework. Nevertheless the S-matrix is the central object of the interacting theory. It determines the theory since the local interacting field operators are defined in terms of it.
In the causal approach infrared divergences are completely independent of the ultraviolet divergences — in particular there cannot be a cancellation of infrared with ultraviolet divergences. We circumvent the problem of infrared divergences by considering only local theories. By a local theory we mean the following situation: We choose an open, bounded domain $`𝒪M`$ in Minkowski space — usually such that it is causally complete — in which the interacting fields are localized and consider the field algebra generated by these fields.
The crucial observation that makes it possible to abandon the adiabatic limit and therefore to avoid infrared divergences is due to Brunetti and Fredenhagen \[BF97\]. They found that a modification of the interaction outside the domain $`𝒪`$ induces only a unitary transformation of the field algebra. Since this does not touch the physical content of the theory, it is in particular possible to switch off the interaction outside $`𝒪`$. With the coupling being a test function, infrared divergences cannot occur.
The chapter gives a short presentation of causal perturbation theory in the formulation of Epstein and Glaser \[EG73\]. For the reader interested in details of the causal approach we refer to the textbook of Scharf \[Sch95\]. We use here the notation of Epstein and Glaser which is different from that in the book of Scharf.
At first we construct the S-matrix by means of time ordered products. In the second section we define interacting field operators in terms of retarded products. Advanced and causal products are defined in the third section. The model we consider is determined by an interaction Lagrangian. It is a polynomial in $`𝒫`$, but not every polynomial in $`𝒫`$ can serve as a Lagrangian. We postulate in the fourth section five conditions such a polynomial must satisfy in order to define a possible Lagrangian.
### 5.1. The S-matrix
The S-matrix is defined as a formal power series in terms of time ordered products of the Lagrangian as
$$S(g)\stackrel{\text{def}}{=}\underset{n=0}{\overset{\mathrm{}}{}}\frac{i^n}{n!}d^4x_1\mathrm{}d^4x_ng(x_1)\mathrm{}g(x_n)T(,\mathrm{},)(x_1,\mathrm{},x_n)$$
(5.1)
Here $`g`$ is the coupling “constant”, i.e. in our approach a real test function in $`𝒟(M)`$. The notation $`(g)`$ in the argument of $`S`$ is of cause only symbolic, since the product of a test function in $`𝒟(M)`$ with a symbol in $`𝒫`$ is not defined. It means that the polynomials are the arguments of the time ordering which are smeared out with the test functions. For sums the symbolic notation means e.g.
$$S(g_1𝒲_1+g_2𝒲_2)\stackrel{\text{def}}{=}1\mathrm{l}+id^4x_1\left[g_1(x_1)T\left(𝒲_1\right)(x_1)+g_2(x_1)T\left(𝒲_2\right)(x_1)\right]+\mathrm{}.$$
(5.2)
The S-matrix is an element in $`\stackrel{~}{\text{ }\mathrm{C}}𝒟`$, i.e. the set of formal power series whose elements are endomorphisms on $`𝒟`$. This is true because $`T(,\mathrm{},)(x_1,\mathrm{},x_n)_n(𝒟)`$ and $`g(x_1)\mathrm{}g(x_n)𝒟(M^n)`$.
The S-matrix is also the generating functional of the time ordered products, i.e. the time ordered products can be recovered from the S-matrix by means of
$$T(W_1,\mathrm{},W_n)(x_1,\mathrm{},x_n)=\frac{\delta ^n}{i^n\delta g_1(x_1)\mathrm{}\delta g_n(x_n)}S\left(\underset{k=1}{\overset{n}{}}g_kW_k\right)|_{g_1=\mathrm{}g_n=0}.$$
(5.3)
The inverse S-matrix $`S^1(g)`$ is also a formal power series. From eqn. (4.21) we conclude
$$S^1(g)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{(i)^n}{n!}d^4x_1\mathrm{}d^4x_ng(x_1)\mathrm{}g(x_n)\overline{T}(,\mathrm{})(x_1,\mathrm{},x_n).$$
(5.4)
The S-matrix is also pseudo unitary, $`S(g)^{}=S^1(g)`$, by means of normalization condition (N2).
### 5.2. Interacting fields and retarded products
The interacting fields are constructed according to Bogoliubov as operator valued distributions by
$$\left(W_i\right)_{\mathrm{int}}^g\left(y\right)\stackrel{\text{def}}{=}S(g)^1\frac{\delta }{i\delta h(y)}S(g+hW_i)|_{h=0}.$$
(5.5)
Here $`h`$ is a test function in $`𝒟(M)`$. The corresponding localized field operators are
$$\left(W_i\right)_{\mathrm{int}}^g\left(f\right)\stackrel{\text{def}}{=}d^4yf(y)\left(W_i\right)_{\mathrm{int}}^g\left(y\right)$$
(5.6)
where $`f`$ is a test function with support in the domain $`𝒪`$ as. The algebra of field operators that are localized in $`𝒪`$ is denoted as $`\stackrel{~}{}(𝒪)`$.
Inserting the definition of the S-matrix the distributional field operators can be written as
$$\begin{array}{cc}\hfill \left(W_i\right)_{\mathrm{int}}^g\left(y\right)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{i^n}{n!}d^4x_1\mathrm{}d^4x_ng(x_1)\mathrm{}g(x_n)& \\ \hfill \times R(,\mathrm{},;W_i)(x_1,\mathrm{},x_n;y)& .\hfill \end{array}$$
(5.7)
This expression contains the so called retarded or $`R`$-products whose definition in terms of $`T`$\- and $`\overline{T}`$-products reads
$$\begin{array}{cc}& R(W_1,\mathrm{},W_n;W_i)(x_1,\mathrm{},x_n;y)\hfill \\ & \stackrel{\text{def}}{=}\underset{YX}{}(1)^{\left|Y\right|}\overline{T}\left(W_Y\right)(x_Y)T(W_{Y^c},W_i)(x_{Y^c},y).\hfill \end{array}$$
(5.8)
Here $`X=\{x_1,\mathrm{},x_n\}`$. For the notation we refer to eqn. (4.11).
According to eqn. (4.21) the retarded products can be alternatively expressed as
$$\begin{array}{cc}& R(W_1,\mathrm{},W_n;W_i)(x_1,\mathrm{},x_n;y)\hfill \\ & =\underset{YX}{}(1)^{\left|Y\right|}\overline{T}(W_Y,W_i)(x_Y,y)T\left(W_{Y^c}\right)(x_{Y^c}).\hfill \end{array}$$
(5.9)
Causality (4.13) implies that the retarded products have retarded support (justifying their name), i.e.
$$\begin{array}{cc}& R(W_X;W_i)(x_X,y)\hfill \\ & \{(x_1,\mathrm{},x_n,y)M^{n+1}:x_i(y+\overline{V}_{})x_iX\}.\hfill \end{array}$$
(5.10)
The interacting fields in $`\stackrel{~}{}(𝒪)`$ therefore depend only on the interaction in the past of $`𝒪`$. From the definition of the interacting field distributions Dütsch and Fredenhagen derive in \[DF99\] the commutator relation:
$$\begin{array}{cc}\hfill [\left(W^1\right)_{\mathrm{int}}^g\left(x\right),\left(W^2\right)_{\mathrm{int}}^g\left(y\right)]_{}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{i^n}{n!}d^4x_1\mathrm{}d^4x_ng(x_1)\mathrm{}g(x_n)\times & \\ \hfill \{R(,\mathrm{},,W^1;W^2)(x_1,\mathrm{},x_n,x;y)& \\ \hfill R(,\mathrm{},,W^2;W^1)(x_1,\mathrm{},x_n,y;x)\}.& \end{array}$$
(5.11)
### 5.3. The advanced and the causal product
The advanced product is defined as
$$\begin{array}{cc}& A(W_1,\mathrm{},W_n;W_i)(x_1,\mathrm{},x_n;y)\hfill \\ & \stackrel{\text{def}}{=}\underset{YX}{}(1)^{\left|Y\right|}T\left(W_{Y^c}\right)(x_{Y^c})\overline{T}(W_Y,W_i)(x_Y,y)\hfill \end{array}$$
(5.12)
or, with the alternative expression analogous to eqn. (5.9),
$$\begin{array}{cc}& A(W_1,\mathrm{},W_n;W_i)(x_1,\mathrm{},x_n;y)\hfill \\ & =\underset{YX}{}(1)^{\left|Y\right|}T(W_{Y^c},W_i)(x_{Y^c},y)\overline{T}\left(W_Y\right)(x_Y).\hfill \end{array}$$
(5.13)
They have advanced support,
$$\begin{array}{cc}& A(W_X;W_i)(x_X,y)\hfill \\ & \{(x_1,\mathrm{},x_n,y)M^{n+1}:x_i(y+\overline{V}_+)x_iX\}.\hfill \end{array}$$
(5.14)
The interacting fields can also be defined in terms of advanced products instead of retarded products without changing the local field algebra if we define
$$\left(W_i\right)_{\mathrm{int}}^g\left(f\right)=d^4yf(y)\frac{\delta }{i\delta h(y)}S(g+hW_i)|_{h=0}\times S(g)^1.$$
(5.15)
This would only result in a unitary transformation on $`\stackrel{~}{}(𝒪)`$ with $`S(g)`$ as the unitary operator.
Finally we define the causal product as
$$\begin{array}{cc}& D(W_1,\mathrm{},W_n;W_i)(x_1,\mathrm{},x_n;y)\hfill \\ & \stackrel{\text{def}}{=}R(W_1,\mathrm{},W_n;W_i)(x_1,\mathrm{},x_n;y)A(W_1,\mathrm{},W_n;W_i)(x_1,\mathrm{},x_n;y)\hfill \end{array}$$
(5.16)
which has obviously causal support:
$$\begin{array}{cc}& D(W_X;W_i)(x_X,y)\hfill \\ & \{(x_1,\mathrm{},x_n,y)M^{n+1}:x_i(y+\overline{V})x_iX\}.\hfill \end{array}$$
(5.17)
### 5.4. Conditions on the interaction Lagrangian
Up to now the Lagrangian density $``$ that defines the model via the S-matrix could have been an arbitrary polynomial in $`𝒫`$. There is a number of restrictions that such a polynomial must satisfy before it can define a reasonable physical model. In this section we will collect these restrictions.
At first, it must be Lorentz invariant:
$$\left(\mathrm{Ad}U(p)\right)T\left(\right)(x)=T\left(\right)(\mathrm{\Lambda }^1x)p=(0,\mathrm{\Lambda })𝒫_+^{}.$$
(C1)
The second condition it must satisfy is pseudo-unitarity:
$$\left(T\left(\right)\right)^{}(x)=T\left(\right)(x).$$
(C2)
Furthermore it must have vanishing ghost number,
$$s_cT\left(\right)(x)=0.$$
(C3)
A Lagrangian with non vanishing ghost number would define a strange theory. The individual orders in perturbation theory of an interacting field would have a ghost number increasing (or decreasing) with the order. Such a theory would be super-renormalizable, provided it is power counting renormalizable, see below.
Since the S-matrix should also be BRS-invariant, one could also expect an equation like
$$s_0T\left(\right)(x)=0$$
(5.18)
to hold. Unfortunately it is in general — and specifically in QED and Yang-Mills-theory — impossible to find a Lagrangian for which eqn. (5.18) holds. So we must weaken the condition a little. Therefore we demand that there exist polynomials $`_n^{\mu _1\mathrm{}\mu _n}`$ in $`𝒫`$, the so called $`Q(n)`$-vertices, such that the following equations hold:
$$s_0T\left(_n^{\mu _1,\mathrm{},\mu _n}\right)(x)=i_\rho ^xT\left(_{n+1}^{\mu _1,\mathrm{},\mu _n,\rho }\right)(x).$$
(C4)
The $`Q(n)`$-vertices must be totally antisymmetric in their Lorentz indices. The other index indicates the ghost number
$$g\left(_n\right)=n_n.$$
(5.19)
In (C4) there will be only a finite number of nontrivial equations, i.e. there exists an $`m\mathrm{I}\mathrm{N}`$ such that $`_m=0`$. The $`Q(n)`$-vertices have always the same canonical dimension as the original vertex $``$, and they also contain the same number of generators. Therefore $`_5=0`$ for power counting renormalizable theories (see below) since $`_5`$ must have ghost number five and it is impossible to construct a polynomial with ghost number five and a canonical dimension not exceeding four. If the original vertex contains only three generators as it is usually the case then already $`_4=0`$. In Yang-Mills theory — with or without matter — even $`_3=0`$ and in QED $`_2=0`$. These results can be derived by explicit calculation.
The last condition on the Lagrangian we want to impose is power counting renormalizability. Perturbation theories can be divided into three groups according to the canonical dimension of their Lagrangian: Those with a canonical dimension less than the spacetime dimension are super renormalizable, that means the number of free normalization parameters decreases with the order and finally vanishes, so the theory is completely determined by a finite number of such parameters. Power counting renormalizable theories are those where the canonical dimension equals the space time dimension. For those theories there exists for all orders in perturbation theory a common upper bound for the number of free parameters in the extension. Non renormalizable theories have Lagrangians whose canonical dimension exceeds the spacetime dimension, and this leads to a number of free normalization parameters that may increase with the order. Although the predictive power of such theories — perturbative gravitation is an example of those — is rather poor, it is nevertheless possible to deal with them in the framework of causal perturbation theory.
For our considerations non renormalizable Lagrangians play no role and therefore we exclude them explicitely. As we always work in four spacetime dimensions, the condition for renormalizability reads
$$\mathrm{deg}4,$$
(C5)
where $`\mathrm{deg}`$ means the canonical dimension.
## 6. The interacting theory
In this chapter we come back to the program for the construction of interacting gauge theories outlined in chapter (2). We formulated at the end of section (2.3) four requirements for an interacting gauge theory. With the construction of local interacting field theories in the last chapter and the normalization conditions in chapter (4) we are now able to determine under which conditions these requirements can be accomplished. The first condition — the condition that suitable ghost and BRS charges can be found in the free model — must be verified for the individual model. This has been done for the free models underlying QED and Yang-Mills theory in section (3.5). In this chapter we will see that the other three conditions hold if all normalization conditions (N1) - (N6) are satisfied and if the conditions (C1) - (C5) are valid for the Lagrangian $``$ which defines the model. We assume throughout this chapter that these preconditions hold.
In the first section we collect a number of properties all interacting fields share from their very definition. Among them are e.g. covariance and locality. In addition we derive a relation between the interacting field operators for the higher generators and those for the basic generators.
In the second section we formulate field equations for the interacting field operators. These equations are determined by normalization condition (N4).
In the third section we come to interacting operators that are of particular importance in gauge theories. In this section we define the interacting ghost current, the interacting ghost charge and the ghost number of interacting fields. We prove that the interacting ghost current is conserved and that the higher order contributions of the ghost charge vanish. As a consequence every interacting field has the same ghost number as the corresponding free field.
In the fourth section we define the most essential operators in an interacting gauge theory: the interacting BRS current, the interacting BRS charge and the interacting BRS transformation. We find that the interacting BRS current is conserved only where the test function $`g`$ that defines the coupling is constant. The BRS charge is constructed only for spacetimes that are compactified in spacelike directions. Otherwise its definition would not be well posed. We prove also that with our definitions the BRS algebra holds. This means in particular that the interacting BRS charge is nilpotent.
In the last section we examine the relation between the quantum field theory defined above and its corresponding classical theory and formulate a correspondence law for these theories.
### 6.1. General properties of interacting fields
We begin our considerations with
C-numbers: From the definition of the retarded products, eqn. (5.8), we can find that they vanish if at least one of their arguments is a multiple of the identity — provided the total number of arguments is at least two, see \[DF99\]. This implies immediately for interacting fields that are generated by $`\mathrm{C}`$-numbers that they possess no higher order terms:
$$\left(\alpha 1\mathrm{l}\right)_{\mathrm{int}}^g\left(x\right)=\alpha 1\mathrm{l},\alpha \text{ }\mathrm{C}.$$
(6.1)
Lorentz covariance: The fact that the Lagrangian is a Lorentz scalar implies, together with condition (N1), the Lorentz transformation properties of the interacting field operators:
$$\left(\mathrm{Ad}U(p)\right)\left(W_i\right)_{\mathrm{int}}^g\left(x\right)=\left(_\mathrm{\Lambda }\left(W_i\right)\right)_{\mathrm{int}}^{g^p}\left(xa\right),p=(a,\mathrm{\Lambda })𝔓_+^{}$$
(6.2)
where $``$ is the representation of the Lorentz group (or its covering group) defined in section (3.1) and $`g^p=g(\mathrm{\Lambda }^1xa)`$.
Pseudo-hermiticity: Due to the conditions (C2) and (N2) the Krein adjoint of the interacting fields is given by
$$\left(\left(W_i\right)_{\mathrm{int}}^g\left(x\right)\right)^{}=\left(W_i^{}\right)_{\mathrm{int}}^g\left(x\right)W_i𝒫.$$
(6.3)
The -involution on the right hand side is the one introduced in section (3.1). Locality: A very important property of interacting fields is their locality. This means that two interacting field operators (anti-) commute with each other if they are localized in spacelike separated regions. This can immediately be derived from eqn. (5.11):
$$[\left(W^1\right)_{\mathrm{int}}^g\left(x\right),\left(W^2\right)_{\mathrm{int}}^g\left(y\right)]_{}=0\text{ if }xy.$$
(6.4)
Primary interacting fields: Due to normalization condition (N4) the interacting fields for the higher generators may be expressed by those for the basic generators as:
$$\begin{array}{cc}\hfill \left((\phi _i)^{(n,\nu _1\mathrm{}\nu _n)}\right)_{\mathrm{int}}^g\left(x\right)=& _x^{\nu _1}\mathrm{}_x^{\nu _n}\left((\phi _i)^{(0)}\right)_{\mathrm{int}}^g\left(x\right)\hfill \\ & +C_{\phi _i,n}g(x)\left(\frac{}{\stackrel{~}{\phi }_i^{(n,\nu _1\mathrm{}\nu _n)}}\right)_{\mathrm{int}}^g\left(x\right),\hfill \end{array}$$
(6.5)
where $`\stackrel{~}{\phi }_i`$ is the field conjugated to $`\phi _i`$.
### 6.2. The interacting field equations
Now we state field equations for the interacting field theory. They are again already determined by condition (N4) and read
$$\underset{j}{}D_{ij}^x\left(\phi _j\right)_{\mathrm{int}}^g\left(x\right)=g(x)\left(\frac{}{\phi _i}\right)_{\mathrm{int}}^g\left(x\right).$$
(6.6)
Inserting here the definition of $`D^x`$ — eqn. (3.86) and the following ones — we find that this implies in particular
$$\begin{array}{c}\hfill K^{\phi _i,x}\left((\phi _i)^{(0)}\right)_{\mathrm{int}}^g\left(x\right)=\underset{n=0}{\overset{\mathrm{}}{}}(1)^n_x^{\nu _1}\mathrm{}_x^{\nu _n}\left(g(x)\left(\frac{}{(\stackrel{~}{\phi }_i)^{(n,\nu _1\mathrm{}\nu _n)}}\right)_{\mathrm{int}}^g\left(x\right)\right),\end{array}$$
(6.7)
where $`K^{\phi _i,x}`$ was defined in eqn. (3.88). These are exactly the field equations that are derived as the Euler-Lagrange equations for a classical field theory with a Lagrangian $`_0+`$, where $``$ is the interaction Lagrangian and $`_0`$ is the free Lagrangian that implies the free field equations
$$K^{\phi _i,x}\phi _i(x)=0,\phi _i(x)\text{ a classical field.}$$
(6.8)
But there is one important difference between the field equations in the classical theory and those in the quantum theory. While the classical field equations govern the dynamics of the system, this in not true for the quantum field equations. The reason is that the classical theory has fewer independent variables. The field equations determine the time evolution of the basic fields on the left hand side. Therefore the time evolution of the entire classical theory is determined by the field equations, since all variables are basic fields or products thereof. This is not true in the quantum theory, because the interacting fields for composed elements in the algebra $`𝒫`$ are not products of those for the generators<sup>19</sup><sup>19</sup>19A product of distributional field operators is not defined a priori. It can be examined in the framework of operator product expansions \[Wil69, Wil71, Zim73\], but we will not discuss this here.. Therefore the time evolution of the interacting fields for composed elements of $`𝒫`$ is left open by the equations above.
The quantum field equations are completely independent of the normalization constants $`C_{\phi _i,k}`$ in eqn. (3.78). They are also independent of the normalization of time ordered products, provided condition (N4) applies.
### 6.3. The interacting ghost current and the ghost charge
The interacting ghost current is defined as the interacting field operator that is generated by the free ghost current $`k^\mu `$, see section (3.5):
$$\stackrel{~}{k}^\mu (x)\stackrel{\text{def}}{=}\left(k^\mu \right)_{\mathrm{int}}^g\left(x\right).$$
(6.9)
This current is conserved as is easily derived by means of (N5) and (C3):
$$_\mu ^x\stackrel{~}{k}^\mu (x)=0.$$
(6.10)
From eqn. (6.3) and the fact that the free ghost current is anti-pseudo-hermitian we find that the interacting ghost current is anti-pseudo-hermitian, too:
$$\left(\stackrel{~}{k}^\mu (x)\right)^{}=\stackrel{~}{k}^\mu (x).$$
(6.11)
The interacting ghost charge is defined as
$$\stackrel{~}{Q}_c\stackrel{\text{def}}{=}\underset{\lambda 0}{lim}d^4yh_\lambda (y)\stackrel{~}{k}^0(y),$$
(6.12)
where $`h_\lambda (x^0,𝐱)=\lambda h^t(\lambda x^0)b(\lambda 𝐱)`$, see eqn. (3.103). Here the coordinate frame is chosen such that the origin $`0`$ is in the domain $`𝒪`$ where the fields are localized. We restrict the admissible spatial test functions $`b`$: At first the temporal test function $`h^t`$ is selected such that $`0h^t`$ and the following equation holds:
$$\left((g)\left[𝒪+\overline{V}_+\right]\right)\left(h^t\times \mathrm{I}\mathrm{R}^3\right)\left((g)\left[𝒪+\overline{V}_{}\right]\right).$$
(6.13)
Then only test functions $`b`$ are admitted with the following properties: $`b(𝐱)=1`$ for all $`𝐱\mathrm{I}\mathrm{R}^3`$ for which an $`x^0h^t`$ exists such that
$$(x^0,𝐱)\left(g+\overline{V}_+\right).$$
(6.14)
The question arises whether the limit in the definition of $`\stackrel{~}{Q}_c`$ exists. We will show that this is indeed true.
The zeroth order of the interacting ghost current is simply the free ghost current. For the free current we know already that the limit exists, so we confine our attention to the higher orders.
We will prove that the higher orders of the ghost charge do not depend on $`\lambda `$. For this purpose we calculate for the $`n^{\mathrm{th}}`$ order of the ghost charge, $`n1`$ and $`\lambda 1`$:
$$Q_{c,\lambda }^nQ_{c,1}^n=d^4x\left(h_\lambda (x)h_1(x)\right)\stackrel{~}{k}^{0,n}(x).$$
(6.15)
Here $`\stackrel{~}{k}^{\mu ,n}`$ is the $`n^{\mathrm{th}}`$ order of the ghost current. We have for all $`n1`$ that $`\stackrel{~}{k}^{\mu ,n}\left(g+\overline{V}_+\right)`$ due to the support properties of the retarded products.
With our conventions for the test functions we can substitute in eqn. (6.15) on the right hand side $`h^t(x^0)b(\lambda x)`$ for $`h_1(x)=h^t(x^0)b(x)`$ because $`h^t(x^0)\left(b(\lambda x)b(x)\right)`$ vanishes on the support of $`\stackrel{~}{k}^{\mu ,n}`$, $`n1`$. Then eqn. (6.15) becomes
$$Q_{c,\lambda }^nQ_{c,1}^n=d^4x\left(\lambda h^t(\lambda x^0)h^t(x^0)\right)b(\lambda 𝐱)\stackrel{~}{k}^{0,n}(x).$$
(6.16)
There exists a test function $`H_\lambda 𝒟(\mathrm{I}\mathrm{R})`$ such that
$$_0^xH_\lambda (x^0)=\left(\lambda h^t(\lambda x^0)h^t(x^0)\right).$$
(6.17)
Inserting this into (6.15) we get
$$\begin{array}{cc}\hfill Q_{c,\lambda }^nQ_{c,1}^n& =d^4x\left(_0^xH_\lambda (x^0)\right)b(\lambda 𝐱)\stackrel{~}{k}^{0,n}(x)\hfill \\ & =d^4xH_\lambda (x^0)\left(_i^xb(\lambda 𝐱)\right)\stackrel{~}{k}^{i,n}(x),\hfill \end{array}$$
(6.18)
where we have partially integrated and used the fact that $`\stackrel{~}{k}^\mu `$ is conserved. By construction we have
$$\left(H_\lambda (x^0)\left(_i^xb(\lambda 𝐱)\right)\right)\left(g+\overline{V}_+\right)=\mathrm{}.$$
(6.19)
Comparing this with the support of $`\stackrel{~}{k}^{\mu ,n}`$, we see that the integral vanishes. Therefore the higher orders of $`\stackrel{~}{Q}_c`$ do not depend on $`\lambda `$. Even more, because of current conservation, eqn. (6.10), one can choose $`h^t`$ such that the support of $`h_1`$ is entirely in the past of $`g`$. Then the higher order terms vanish due to the support properties of the retarded products, so the interacting ghost charge coincides with the free ghost charge or, strictly speaking since $`\stackrel{~}{Q}_c`$ is a formal power series,
$$\stackrel{~}{Q}_c=(Q_c,0,0,\mathrm{}).$$
(6.20)
Since the ghost current is anti-pseudo-hermitian, the ghost charge is it, too:
$$\stackrel{~}{Q}_c^{}=\stackrel{~}{Q}_c.$$
(6.21)
The interacting ghost number of a localized field operator is measured by the following derivation:
$$\stackrel{~}{s}_c\left(\left(W_i\right)_{\mathrm{int}}^g\left(x\right)\right)\stackrel{\text{def}}{=}[\stackrel{~}{Q}_c,\left(W_i\right)_{\mathrm{int}}^g\left(x\right)]_{}.$$
(6.22)
As the interacting ghost charge coincides with the free one, we have
$$\stackrel{~}{s}_c\left(\left(W_i\right)_{\mathrm{int}}^g\left(x\right)\right)=s_c\left(\left(W_i\right)_{\mathrm{int}}^g\left(x\right)\right).$$
(6.23)
This implies immediately, due to (N5) and (C3), that the interacting field operators have the same ghost number as the corresponding free fields:
$$\stackrel{~}{s}_c\left(\left(W_i\right)_{\mathrm{int}}^g\left(x\right)\right)=g(W_i)\left(W_i\right)_{\mathrm{int}}^g\left(x\right),g(W_i)\mathrm{𝖹𝖹}.$$
(6.24)
### 6.4. The interacting BRS current, BRS charge and BRS transformation
The natural choice for the BRS current,
$$\stackrel{~}{ȷ}_B^\mu (x)=\left(j_B^\mu \right)_{\mathrm{int}}^g\left(x\right),$$
(6.25)
is not conserved in general, so this cannot be the correct interacting BRS current. The situation is even worse: Explicit calculations in first order QED and Yang-Mills theory shows that there exists no normalization of the time ordered products such that this current is conserved even in first order, irrespective of our normalization conditions. The best one can achieve is that the current is conserved where the coupling is constant, and even this seemingly liberal condition fixes the normalization in first order uniquely.
A direct calculation reveals that this normalization is not compatible with the normalization conditions (N3) and (N4). But there is an expression for the interacting BRS current that is compatible with the normalization conditions in first order and that is conserved in the sense above, not only for Yang-Mills theories but for every theory. Adopting this expression as the definition of the interacting BRS current we have
$$\stackrel{~}{ȷ}_B^\mu (x)\stackrel{\text{def}}{=}\left(j_B^\mu \right)_{\mathrm{int}}^g\left(x\right)g(x)\left(_1^\mu \right)_{\mathrm{int}}^g\left(x\right),$$
(6.26)
where $`_1^\mu `$ is the $`R(1)`$-vertex, see condition (C4). Normalization condition (N6) implies that this current is indeed conserved where the coupling $`g`$ is constant:
$$_\mu ^x\stackrel{~}{ȷ}_B^\mu (x)=(_\nu g)(x)\left(_1^\nu \right)_{\mathrm{int}}^g\left(x\right).$$
(6.27)
The fact that the interacting BRS current is not everywhere conserved is a severe drawback, since it complicates the definition of the BRS charge, see below. So the question arises whether a more clever choice for the BRS current could have yielded one that is everywhere conserved. But this turns out to be impossible in general. Concretely, in QED as well as in Yang-Mills theory the explicit calculation shows already in first order that no such choice exists. So in general this result cannot be improved.
By the same reasoning as for eqn. (6.11) one derives that $`\stackrel{~}{ȷ}_B^\mu `$ is pseudo-hermitian
$$\left(\stackrel{~}{ȷ}^\mu (x)\right)^{}=\stackrel{~}{ȷ}^\mu (x).$$
(6.28)
Now we come to the definition of the BRS charge. As already mentioned this definition is more difficult than that of the ghost charge was, since the BRS current is not everywhere conserved. The problem can be seen as follows: The natural choice for the BRS charge would be
$$\stackrel{~}{Q}_B=\underset{\lambda 0}{lim}d^4yh_\lambda (y)\stackrel{~}{ȷ}_B^0(y),$$
(6.29)
with $`h_\lambda `$ like above. Unfortunately this expression would depend on the choice of $`h_\lambda `$, unlike for $`\stackrel{~}{Q}_c`$, and the higher orders would depend on $`\lambda `$, both because the BRS current is not conserved. If the higher orders depend on $`\lambda `$ the limit is no longer under control.
In order to make $`\stackrel{~}{Q}_B`$ independent of $`h_\lambda `$, the support of $`h_\lambda `$ must be for every $`\lambda `$ in a region where $`g`$ is constant. This would mean that $`g`$ is everywhere constant, i.e. the adiabatic limit must be performed, and this limit does not exist in general.
Another possibility would be not to perform the limit and choose e.g. $`h_1`$ as a test function in the definition of the BRS charge. In this case the BRS charge would clearly become well defined, but it would also be a local operator, and such an operator could not annihilate states with finite energy due to the theorem of Reeh and Schlieder. It is unlikely that the cohomology defined with it has good properties, and therefore we exclude this possibility.
The way out of this seemingly pitfall was found by Dütsch and Fredenhagen \[DF99\]. In order to allow functions that are constant in spacelike directions as test functions, they embed the double cone $`𝒪`$ isometrically into the cylinder $`\mathrm{I}\mathrm{R}\times C_L`$ with $`\mathrm{I}\mathrm{R}`$ the time axis and $`C_L`$ a cube of length $`L`$ sufficiently big to contain $`𝒪`$. This spatial compactification does not change the properties of the local algebra $`\stackrel{~}{}(𝒪)`$. This is why the quantization of free fields in a box, mentioned in chapter (3), is important for us. For the details of the construction we refer to \[DF99\].
In the compactified space $`h`$ and $`g`$ can be chosen to be test functions such that $`g`$ is constant on $`h`$ with the same value as on $`𝒪`$.
With these test functions we are able to give a definition of the interacting BRS current in a spatially compactified spacetime. At first, we choose the test function $`h`$ to be
$$h(x)=h^t(x_0),h^t\text{ like in (}\text{3.103}\text{)},h𝒟(\mathrm{I}\mathrm{R}\times C_L),$$
(6.30)
and the coupling $`g`$ such that
$$g|_h=g|_𝒪=\text{ constant},g𝒟(\mathrm{I}\mathrm{R}\times C_L).$$
(6.31)
With $`g`$ and $`h`$ now both being a test function — on $`\mathrm{I}\mathrm{R}\times C_L`$ — the BRS charge can be defined as
$$\stackrel{~}{Q}_B\stackrel{\text{def}}{=}_{\mathrm{I}\mathrm{R}\times C_L}d^4yh(y)\stackrel{~}{ȷ}_B^{\mathrm{\hspace{0.17em}0}}(y).$$
(6.32)
It is easy to see by an analogous reasoning as for the ghost charge that this BRS charge is independent of $`h`$.
The zeroth order of this BRS charge agrees with the free BRS charge in the $`\mathrm{I}\mathrm{R}\times C_L`$ spacetime, $`\stackrel{~}{Q}_{B,0}=Q_B`$, if $`h_\lambda `$ is replaced there by $`h`$. Of course the limit is then not performed because it would be void. Unlike the interacting ghost charge the interacting BRS charge has also non vanishing higher order contributions. The reasoning which showed that the higher contributions of $`\stackrel{~}{Q}_c`$ vanish cannot be applied here, since $`h`$ may not be (not even partly) in the past of $`g`$ from its very definition.
Like the BRS current the BRS charge is pseudo-hermitian:
$$\stackrel{~}{Q}_B^{}=\stackrel{~}{Q}_B.$$
(6.33)
Since $`s_cQ_B=Q_B`$ in the free theory, we find with eqn. (4.32)
$$[\stackrel{~}{Q}_c,\stackrel{~}{Q}_B]_{}=\stackrel{~}{s}_c\left(\stackrel{~}{Q}_B\right)=\stackrel{~}{Q}_B.$$
(6.34)
So the first part of the BRS algebra holds. The most important property of the BRS charge is its nilpotency, the second part of the BRS algebra. This will be proven next.
To this end we write at first $`\stackrel{~}{Q}_B`$ in a different form that is more adequate for the proof. We use for the interacting field $`\left(j^\mu \right)_{\mathrm{int}}^g\left(x\right)`$ in the definition of $`\stackrel{~}{ȷ}^\mu `$ the equation (B.25) from the appendix and the identity $`\stackrel{~}{Q}_c=Q_c`$, eqn. (6.20). With it the interacting BRS charge can be written as
$$\begin{array}{cc}\hfill \stackrel{~}{Q}_B=Q_B+\underset{n=0}{\overset{\mathrm{}}{}}\frac{i^n}{n!}& d^4yd^4zd^4x_1\mathrm{}d^4x_nh(y)(_\nu g)(z)\times \hfill \\ & \times g(x_1)\mathrm{}g(x_n)R(,\mathrm{},,_1^\nu ;k^0)(x_1,\mathrm{}x_n,z;y).\hfill \end{array}$$
(6.35)
The first term in the sum on the right hand side is the free BRS charge. Since its properties are already known, we confine our attention to the higher order terms. From the form of these terms one can see immediately that the interacting BRS charge becomes the free one in the adiabatic limit, if this limit exists. We are here particularly interested in theories where the adiabatic limit does not necessarily exist.
We introduce a test function $`H𝒟(\mathrm{I}\mathrm{R}\times C_L)`$ with the property
$$_\mu ^yH(y)=\delta _\mu ^0h(y),H𝒟(\mathrm{I}\mathrm{R}\times C_L),$$
(6.36)
such that $`H(x)=1`$ for all $`x`$ in the past of $`g`$ and $`H(x)=0`$ for all $`x`$ in the future of $`g`$. Inserting this into the expression above, we find by partial integration and with the help of eqn. (N5) the following alternative formulation for the $`n^{\mathrm{th}}`$ order of $`\stackrel{~}{Q}_B`$, $`n1`$:
$$\begin{array}{cc}\hfill \stackrel{~}{Q}_B^{(n)}=& \frac{i^{n1}}{(n1)!}d^4zd^4x_1\mathrm{}d^4x_{n1}H(z)(_\nu g)(z)\times \hfill \\ & \times g(x_1)\mathrm{}g(x_{n1})R(,\mathrm{},;_1^\nu )(x_1,\mathrm{},x_{n1};z).\hfill \end{array}$$
(6.37)
The $`n^{\text{th}}`$ order of $`(\stackrel{~}{Q}_B)^2`$ decomposes according to
$$\left((\stackrel{~}{Q}_B)^2\right)^{(n)}=\underset{k=0}{\overset{n}{}}(\stackrel{~}{Q}_B)^{(k)}(\stackrel{~}{Q}_B)^{(nk)}=s_0\left((\stackrel{~}{Q}_B)^{(n)}\right)+\underset{k=1}{\overset{n1}{}}(\stackrel{~}{Q}_B)^{(k)}(\stackrel{~}{Q}_B)^{(nk)}.$$
(6.38)
At first we will calculate $`s_0\left((\stackrel{~}{Q}_B)^{(n)}\right)`$. With the help of the generalized operator gauge invariance, eqn. (4.35), with $`i_1=1`$ and $`i_k=0`$ otherwise, we get
$$\begin{array}{cc}& s_0(\stackrel{~}{Q}_B)^{(n)}=\hfill \\ & =\frac{i^{n2}}{(n2)!}d^4x_1\mathrm{}d^4x_{n2}d^4yd^4zg(x_1)\mathrm{}g(x_{n2})\hfill \\ & \times (_\rho g)(y)H(y)(_\mu g)(z)H(z)R(,\mathrm{},,_1^\rho ;_1^\mu )(x_1,\mathrm{},x_{n2},y;z)\hfill \\ & +\frac{i^n}{(n1)!}d^4x_1\mathrm{}d^4x_{n1}d^4zg(x_1)\mathrm{}g(x_{n1})\hfill \\ & \times \left(_\rho ^z\left[(_\mu g)(z)H(z)\right]\right)R(,\mathrm{},;_2^{\mu \rho })(x_1,\mathrm{},x_{n1};z).\hfill \end{array}$$
(6.39)
Here an additional factor $`H(y)`$ has been inserted in the first integral. This factor does not change the result due to the retarded support of the distribution.
Let us at first consider the second integral on the right hand side. Calculating the derivative of the square bracket in the last line, we get $`(_\mu _\rho g)(z)H(z)+(_\mu g)(z)(_\rho )H(z)`$. The second term vanishes since the supports of $`g`$ and $`H`$ are disjoint. The first term is symmetric in $`\mu `$ and $`\rho `$ while the retarded product is antisymmetric in these indices, due to the antisymmetry of the $`Q(2)`$-vertex. Therefore the entire second integral vanishes.
Now we come to the first integral. Here the test functions are also symmetric under permutation of $`(z,\mu )`$ and $`(y,\rho )`$. If we look at the definition of the retarded products, eqn. (5.8), we see that there are terms where both $`_1^\mu `$ and $`_1^\rho `$ appear as arguments in the same time ordered product or antichronological product. These contributions vanish, because the distributions are antisymmetric in $`(z,\mu )`$ and $`(y,\rho )`$ due to graded symmetry $`(\mathrm{𝐏𝟐})`$. The only contributions that remain lead to our final expression for $`s_0(\stackrel{~}{Q}_B)^{(n)}`$:
$$\begin{array}{cc}& s_0(\stackrel{~}{Q}_B)^{(n)}=\hfill \\ & =\frac{i^{n2}}{(n2)!}d^4x_1\mathrm{}d^4x_{n2}d^4yd^4zg(x_1)\mathrm{}g(x_{n2})(_\rho g)(y)H(y)\hfill \\ & \times (_\mu g)(z)H(z)\underset{YX}{}(1)^{\left|Y\right|}\overline{T}(,\mathrm{},,_1^\rho )(x_Y,y)T(,\mathrm{},,_1^\mu )(x_{Y^c},z)\hfill \end{array}$$
(6.40)
with $`X=\{x_1,\mathrm{},x_{n2}\}`$.
To calculate $`_{k=1}^{n1}(\stackrel{~}{Q}_B)^{(k)}(\stackrel{~}{Q}_B)^{(nk)}`$ we make use of the two ways to express $`R`$-products in terms of $`T`$\- and $`\overline{T}`$-products, that means we use eqn. (6.37) for the individual orders of the BRS charge, inserting eqn. (5.9) for the retarded products on the left hand side and eqn. (5.8) for those on the right hand side. Then we get after a little combinatorial analysis
$$\begin{array}{cc}& \underset{k=1}{\overset{n1}{}}(\stackrel{~}{Q}_B)^{(k)}(\stackrel{~}{Q}_B)^{(nk)}=\hfill \\ & =\frac{i^{n2}}{(n2)!}d^4x_1\mathrm{}d^4x_{n2}d^4yd^4zg(x_1)\mathrm{}g(x_{n2})(_\rho g)(y)H(y)(_\mu g)(z)\hfill \\ & \times H(z)[\underset{Y,Z,U,V}{}(1)^{\left|Z\right|+\left|V\right|}\left(\overline{T}(,\mathrm{},,_1^\mu )(x_Z,z)T(,\mathrm{},)(x_Y)\right)\hfill \\ & \times \left(\overline{T}(,\mathrm{},)(x_V)T(,\mathrm{},,_1^\rho )(x_U,y)\right)],\hfill \end{array}$$
(6.41)
where the sum in the square brackets runs over all disjoint partitions of $`X`$ into four subsets $`U,V,Y,Z`$. These subsets may be empty. This set of partitions can be divided into two subsets, namely the set of those partitions where $`Y`$ and $`V`$ are empty and its complement. This complement can in turn be divided in subsets with $`U`$ and $`Z`$ fixed, yielding terms proportional to
$$\underset{WXUZ}{}(1)^{\left|W\right|}T(,\mathrm{},)(x_W)\overline{T}(,\mathrm{},)(x_{XUZW}).$$
(6.42)
This expression vanishes due to eqn. (4.21) because $`XUZ\mathrm{}`$ according to our assumption. So there remains only a contribution from the partitions with $`Y=V=\mathrm{}`$, and since $`T_0=\overline{T}_0=1\mathrm{l}`$, there remains only
$$\begin{array}{cc}& \underset{k=1}{\overset{n1}{}}(\stackrel{~}{Q}_B)^{(k)}(\stackrel{~}{Q}_B)^{(nk)}=\hfill \\ & =\frac{i^{n2}}{(n2)!}d^4x_1\mathrm{}d^4x_{n2}d^4yd^4zg(x_1)\mathrm{}g(x_{n2})(_\rho g)(y)H(y)\hfill \\ & \times (_\mu g)(z)H(z)\underset{YX}{}(1)^{\left|Y\right|}\overline{T}(,\mathrm{},,_1^\rho )(x_Y,y)T(,\mathrm{},,_1^\mu )(x_{Y^c},z).\hfill \end{array}$$
(6.43)
Obviously this is just the negative of eqn. (6.40), yielding
$$\left((\stackrel{~}{Q}_B)^2\right)^{(n)}=s_0\left((\stackrel{~}{Q}_B)^{(n)}\right)+\underset{k=1}{\overset{n1}{}}(\stackrel{~}{Q}_B)^{(k)}(\stackrel{~}{Q}_B)^{(nk)}\stackrel{!}{=}0.$$
(6.44)
Reviewing our preconditions, we have proven that with our definition the BRS charge is nilpotent — and therefore the complete BRS algebra holds —, provided our normalization condition (N6) is valid.
At the end of this section we come to the interacting BRS transformation $`\stackrel{~}{s}`$. It could be defined as
$$\stackrel{~}{s}\left(\left(W\right)_{\mathrm{int}}^g\left(x\right)\right)=[\stackrel{~}{Q}_B,\left(W\right)_{\mathrm{int}}^g\left(x\right)]_{}.$$
(6.45)
But it turns out to be more clever to permute commutation and integration, and we define
$$\stackrel{~}{s}\left(\left(W\right)_{\mathrm{int}}^g\left(x\right)\right)\stackrel{\text{def}}{=}_{\mathrm{I}\mathrm{R}\times C_L}d^4yh(y)[\stackrel{~}{ȷ}_B^0(y),\left(W\right)_{\mathrm{int}}^g\left(x\right)]_{}$$
(6.46)
with $`h`$ and $`g`$ as in the definition of $`\stackrel{~}{Q}_B`$. The advantage of this definition is that it remains well defined for $`\stackrel{~}{s}`$ acting on local fields in $`\stackrel{~}{}(𝒪)`$ even if the spacetime is not compactified and $`h`$ has compact support only in timelike directions being constant in spacelike directions. This is well defined because locality, eqn. (6.4), holds — both $`\stackrel{~}{ȷ}_B^0`$ and $`\left(W\right)_{\mathrm{int}}^g\left(f\right)`$ are local fields. Therefore the commutator has causal support, so the integrand vanishes in the causal complement of $`𝒪`$.
$$\stackrel{~}{s}\left(\left(W\right)_{\mathrm{int}}^g\left(f\right)\right)\stackrel{\text{def}}{=}d^4yh(y)[\stackrel{~}{ȷ}_B^0(y),\left(W\right)_{\mathrm{int}}^g\left(f\right)]_{}$$
(6.47)
is a well defined expression in Minkowski space, if $`h`$ is chosen such that
$$\begin{array}{cc}& h(x)=h^t(x_0),h^t𝒟(\mathrm{I}\mathrm{R})\text{ as in eqn. (}\text{3.103}\text{)},\hfill \\ & g\text{ is constant on }\left(𝒪+\overline{V}_{}\right)\left(h+\overline{V}_+\right).\hfill \end{array}$$
(6.48)
This expression is independent of $`h`$. The BRS transformation is nilpotent. This can be seen by direct computation — the calculation is then completely analogous to that for $`(\stackrel{~}{Q}_B)^2=0`$ in the compactified spacetime. A different way to prove that $`\stackrel{~}{s}`$ is nilpotent is to consider $`\stackrel{~}{s}`$ in a compactified spacetime — where $`\stackrel{~}{s}^2=0`$ follows directly from $`(\stackrel{~}{Q}_B)^2=0`$. Then let the compactification length $`L`$ tend to infinity. The resultant space will be the Minkowski space and $`\stackrel{~}{s}^2=0`$ still holds since the algebra does not depend on the compactification length. Therefore
$$\stackrel{~}{s}^2A=0A\stackrel{~}{}(𝒪).$$
(6.49)
It is important to note that this reasoning holds only for local operators. In particular the argument of Nakanishi and Ojima \[NO90\] that a nilpotent BRS transformation defines a nilpotent BRS charge can not be applied here. Their argument is as follows:
$$\stackrel{~}{Q}_B\stackrel{\text{def}}{=}\stackrel{~}{s}\stackrel{~}{Q}_c\text{and}0=\stackrel{~}{s}^2(\stackrel{~}{Q}_c)=\stackrel{~}{s}(\stackrel{~}{Q}_B)=2\stackrel{~}{Q}_B^2,$$
(6.50)
but since $`\stackrel{~}{Q}_c`$ is not a local operator it is not in the domain of $`\stackrel{~}{s}`$ in the framework of ordinary spacetime.
So we arrive at the following result: For all investigations concerning the state space it is necessary to compactify spacetime, since we need the BRS charge to define the physical state space, and this is only defined in a compactified spacetime. But for investigations concerning only the algebra of local observables there is no need for a compactification because the definition of observables requires only the BRS transformation, not the BRS charge, and the former can also be defined in an ordinary spacetime.
Summarizing the results of this and the preceeding chapter we see that all the preconditions that we postulated at the end of section (2.3) are satisfied. The only restriction is that the BRS current is conserved only locally, but this is sufficient for the construction of the local interacting gauge theory.
This result was derived under the assumption that the normalization conditions (N1) - (N6) and the conditions on the Lagrangian are satisfied. We proved in chapter (4) that the first five normalization conditions have simultaneous solutions. So the essential point is whether condition (N6) can be satisfied for a model. If this is the case, the construction of the physical state space (in the spatially compactified spacetime) and of the local observable algebra can be performed.
### 6.5. The correspondence between quantum and classical theory
We have seen in section (4.4) that in our approach the propagators for the higher generators are different from the corresponding propagators for the derivated fields in other renormalization procedures. The propagators determine the tree diagrams, and these in turn are known to determine the classical limit of the theory. Therefore the question arises whether the classical limit of our theory is different from what one would expect from other approaches. We will see that this is indeed the case.
The classical fields are functions on a manifold, in this case the Minkowski space. Unlike the distributional field operators they may be multiplied at the same spacetime point. We take advantage of this property and define a representation $`C`$ of the algebra $`𝒫`$ by classical fields. Unlike the representation $`T`$ of $`𝒫`$ in section (3.3) this is not only a linear representation but also an algebra homomorphism. We define
$$\begin{array}{cc}\hfill C:𝒫C^{\mathrm{}}(M),& C(aA)=aC(A)a\text{ }\mathrm{C},A𝒫,\hfill \\ & C\left(\underset{i}{}\phi _i\right)(x)=\underset{i}{}C\left(\phi _i\right)(x),\phi _i𝒢.\hfill \end{array}$$
(6.51)
The representatives of the basic generators are the basic classical fields, i.e. we suppose that there exists for each $`\phi _i𝒢_b`$ a classical field $`\phi _i^{\mathrm{cl}}(x)`$ such that
$$C(\phi _i)(x)=\phi _i^{\mathrm{cl}}(x).$$
(6.52)
The question arises how the higher generators may be represented. The first attempt is to define their representatives as the derivatives of those for the basic generators, e.g.
$$C\left((\phi _i)^{(1,\mu )}\right)(x)=_x^\mu C\left((\phi _i)^{(0)}\right)(x).$$
(6.53)
But this definition is not consistent. This can be seen by comparing this equation with eqn. (6.5),
$$\begin{array}{cc}\hfill \left((\phi _i)^{(n,\nu _1\mathrm{}\nu _n)}\right)_{\mathrm{int}}^g\left(x\right)=& _x^{\nu _1}\mathrm{}_x^{\nu _n}\left((\phi _i)^{(0)}\right)_{\mathrm{int}}^g\left(x\right)\hfill \\ & +C_{\phi _i,n}g(x)\left(\frac{}{\stackrel{~}{\phi }_i^{(n,\nu _1\mathrm{}\nu _n)}}\right)_{\mathrm{int}}^g\left(x\right).\hfill \end{array}$$
(6.54)
If we adopted the definition above, the left hand side and the right hand side of this equation would be equal on the quantum level, but they would have different classical limits, and this cannot be true. We see that the correct prescription for the classical limit of the higher generators is
$$\begin{array}{cc}\hfill C\left((\phi _i)^{(n,\nu _1\mathrm{}\nu _n)}\right)(x)=& _x^{\nu _1}\mathrm{}_x^{\nu _n}C\left((\phi _i)^{(0)}\right)(x)\hfill \\ & +gC_{\phi _i,n}C\left(\frac{}{\stackrel{~}{\phi }_i^{(n,\nu _1\mathrm{}\nu _n)}}\right)(x).\hfill \end{array}$$
(6.55)
Here we have set the coupling $`g`$ constant, since in a classical theory there is no need for the interaction to be switched off. $`\stackrel{~}{\phi }_i`$ is the generator of the field conjugate to $`\phi _i`$.
The fields that correspond to the higher generators are labelled by the normalization constants $`C_{\phi _i,n}`$. This is what we expected when we pointed out the importance of the propagators for the classical limit, because these propagators are also labelled by the normalization constants.
With the representation $`C`$ now defined we can formulate the correspondence law. It states that the distributional interacting field operators become products of classical fields in the classical limit according to
$$\begin{array}{cc}& \left(W\right)_{\mathrm{int}}^g\left(x\right)C\left(W\right)(x)W𝒫.\hfill \\ & g(x)g=\text{ constant.}\hfill \end{array}$$
(6.56)
## 7. Two particular theories
In this chapter we will examine the consequences of our general results derived in the preceeding chapter for two well known models: Quantum electrodynamics and Yang-Mills theory.
### 7.1. Quantum electrodynamics
The fields involved in QED are vector bosons $`A_\mu `$ — the photons —, ghosts and anti-ghosts $`u,\stackrel{~}{u}`$ and charged spinors $`\psi ,\overline{\psi }`$ — the electrons and positrons.
For QED there exists a way to determine the physical state vector space without the BRS formalism — the Gupta-Bleuler procedure. Furthermore the ghosts do not couple to the other fields. Therefore it is not necessary to include the ghosts in the model. Nevertheless we do so because we investigate QED also as a preparation for Yang-Mills theory where the ghosts are indispensable.
The corresponding free theory for QED has been treated in section (3.5).
Therefore we start directly with the interaction. The interaction Lagrangian for QED reads
$$_{QED}=A_\mu j_{\mathrm{el}}^\mu 𝒫.$$
(7.1)
Here $`A_\mu `$ is the basic generator corresponding to the photon field, and the electric current $`j_{\mathrm{el}}^\mu `$ is defined as
$$j_{\mathrm{el}}^\mu \stackrel{\text{def}}{=}\overline{\psi }\gamma ^\mu \psi 𝒫$$
(7.2)
with the basic generators $`\psi ,\overline{\psi }`$ corresponding to the electron and the positron field.
It can be easily verified that this Lagrangian satisfies our requirements (C1) - (C3) and (C5). The canonical dimension of the spinors is 3/2 and that of the photons is 1, summing up to a total canonical dimension of 4, so the model is renormalizable. We will show that also condition (C4) is accomplished.
In addition we examine an important relation that we were not able to prove in the general framework: The normalization condition (N6). We prove that the other normalization conditions, in particular the Ward identities for the electric current, eqn. (N5’), already imply (N6) in QED. The proof will be given below.
To begin with we determine the $`Q(n)`$-vertices of QED from its Lagrangian. For condition (C4) to hold we must find $`Q(n)`$-vertices that satisfy the following equations:
$$s_0T\left(\right)(x)=i_\nu ^xT\left(_1^\nu \right)(x),s_0T\left(_1^\nu \right)(x)=i_\rho ^xT\left(_2^{\rho \nu }\right)(x),\mathrm{}$$
(7.3)
Observing the free BRS transformations introduced in section (3.5), we find that these $`Q(n)`$-vertices exist indeed:
$$_1^\nu =uj_{\mathrm{el}}^\nu ,_i=0i2.$$
(7.4)
To prove that (N6) is valid we must therefore calculate the following expression
$$_\mu ^yT(_1^{\nu _1},\mathrm{},_1^{\nu _k},,\mathrm{},,j_B^\mu )(x_1,\mathrm{},x_n,y)$$
(7.5)
with $`=_{QED}`$, $`_1`$ like above and $`j_B^\mu `$ as defined in section (3.5). If we insert this time ordered product into the causal Wick expansion, eqn, (4.24), we find that it can be written as
$$\begin{array}{cc}& T(_1^{\nu _1},\mathrm{},_1^{\nu _k},,\mathrm{},,(A_\rho )^{(1,\rho )})(x_1,\mathrm{},x_n,y)_y^\mu u(y)\hfill \\ & T(_1^{\nu _1},\mathrm{},_1^{\nu _k},,\mathrm{},,(A_\rho )^{(2,\rho \mu )})(x_1,\mathrm{},x_n,y)u(y).\hfill \end{array}$$
(7.6)
Since neither $``$ nor $`_1`$ contain a higher generator, conditions (N4) reveals that this expression is equal to
$$\begin{array}{cc}& \left(_y^\rho T(_1^{\nu _1},\mathrm{},_1^{\nu _k},,\mathrm{},,(A_\rho )^{(0)})(x_1,\mathrm{},x_n,y)\right)_y^\mu u(y)\hfill \\ & \left(_y^\rho _y^\mu T(_1^{\nu _1},\mathrm{},_1^{\nu _k},,\mathrm{},,(A_\rho )^{(0)})(x_1,\mathrm{},x_n,y)\right)u(y).\hfill \end{array}$$
(7.7)
Inserting the derivative and taking into account the field equations of $`u(y)`$, we find
$$\begin{array}{cc}& _\mu ^yT(_1^{\nu _1},\mathrm{},_1^{\nu _k},,\mathrm{},,j_B^\mu )(x_1,\mathrm{},x_n,y)\hfill \\ & =\left(_y^\mu \mathrm{}^yT(_1^{\nu _1},\mathrm{},_1^{\nu _k},,\mathrm{},,(A_\mu )^{(0)})(x_1,\mathrm{},x_n,y)\right)u(y).\hfill \end{array}$$
(7.8)
With condition (N4) this expression can be rewritten as
$$\begin{array}{cc}& i(\underset{m=k+1}{\overset{n}{}}\left(_\mu ^y\delta (yx_m)\right)\hfill \\ & \times T(_1^{\nu _1},\mathrm{},_1^{\nu _k},,\mathrm{},j_{\mathrm{el}}^\mu ,\mathrm{},)(x_1,\mathrm{},x_n))u(y).\hfill \end{array}$$
(7.9)
Here the vertex $`j_{\mathrm{el}}^\mu `$ is at the $`m^{\mathrm{th}}`$ position. Pulling the derivative out of the bracket we finally arrive at
$$\begin{array}{cc}& i\underset{m=k+1}{\overset{n}{}}_\mu ^m\left(\delta (yx_m)T(_1^{\nu _1},\mathrm{},_1^{\nu _k},,\mathrm{},_1^\mu ,\mathrm{},)(x_1,\mathrm{},x_n)\right)\hfill \\ & i\left(\underset{m=k+1}{\overset{n}{}}\delta (yx_m)_\mu ^mT(_1^{\nu _1},\mathrm{},_1^{\nu _k},,\mathrm{},j_{\mathrm{el}}^\mu ,\mathrm{},)(x_1,\mathrm{},x_n)\right)u(y).\hfill \end{array}$$
(7.10)
The vertices $`_1^\mu `$ and $`j_{\mathrm{el}}^\mu `$ are again in the $`m^{\mathrm{th}}`$ position. Comparing the last line with the Ward identities for the electric current, eqn. (N5’), we see that this term vanishes since $`f()=f(_1)=0`$. The remaining expression is exactly what condition (N6) predicts, provided that all the $`R(n)`$-vertices vanish, $`_1=_2=\mathrm{}=0`$. Condition (N6) was derived using the other normalization conditions, so it must be compatible with all these conditions.
Now we come to the definition of interacting fields. Since the Lagrangian contains no higher generators, the following relations hold due to eqn. (6.5)
$$\left((\phi _i)^{(n,\nu _1\mathrm{}\nu _n)}\right)_{\mathrm{int}}^{g_{QED}}\left(x\right)=_x^{\nu _1}\mathrm{}_x^{\nu _n}\left((\phi _i)^{(0)}\right)_{\mathrm{int}}^{g_{QED}}\left(x\right).$$
(7.11)
We define $`F^{\mu \nu }`$ as
$$F^{\mu \nu }\stackrel{\text{def}}{=}(A^\nu )^{(1,\mu )}(A^\mu )^{(1,\nu )}𝒫.$$
(7.12)
The easiest examples of interacting fields are the ghosts and the anti-ghosts. They do not appear in the Lagrangian $`_{QED}`$ and therefore do not interact. The causal Wick expansion, eqn. (4.24), implies together with the definition of the retarded products, eqn. (5.8),
$$\left(u\right)_{\mathrm{int}}^{g_{QED}}\left(x\right)=u(x)\text{and}\left(\stackrel{~}{u}\right)_{\mathrm{int}}^{g_{QED}}\left(x\right)=\stackrel{~}{u}(x).$$
(7.13)
Due to relation (7.11) we can establish the usual relation for the interacting photon field and the field strength tensor in QED:
$$\left(F^{\mu \nu }\right)_{\mathrm{int}}^{g_{QED}}\left(x\right)=_x^\mu \left(A_\nu \right)_{\mathrm{int}}^{g_{QED}}\left(x\right)_x^\nu \left(A_\mu \right)_{\mathrm{int}}^{g_{QED}}\left(x\right).$$
(7.14)
The field equations for QED are also the usual ones:
$$\begin{array}{cc}& \mathrm{}^x\left(A^\mu \right)_{\mathrm{int}}^{g_{QED}}\left(x\right)=g(x)\left(j_{\mathrm{el}}^\mu \right)_{\mathrm{int}}^{g_{QED}}\left(x\right)\hfill \\ \hfill \text{and}& (i/m)\left(\psi \right)_{\mathrm{int}}^{g_{QED}}\left(x\right)=g(x)\left(\gamma ^\mu A_\mu \psi \right)_{\mathrm{int}}^{g_{QED}}\left(x\right).\hfill \end{array}$$
(7.15)
Furthermore we find that the interacting ghost current and BRS current have a particularly easy form because the ghosts and anti-ghosts do not interact:
$$\begin{array}{cc}& \left(k^\mu \right)_{\mathrm{int}}^{g_{QED}}\left(x\right)=k^\mu (x)\hfill \\ \hfill \text{and}& \left(j_B^\mu \right)_{\mathrm{int}}^{g_{QED}}\left(x\right)=\left(_x^\rho \left(A_\rho \right)_{\mathrm{int}}^{g_{QED}}\left(x\right)\right)_x^\mu u(x)\hfill \\ & \left(_x^\rho _x^\mu \left(A_\rho \right)_{\mathrm{int}}^{g_{QED}}\left(x\right)\right)u(x).\hfill \end{array}$$
(7.16)
Dütsch and Fredenhagen \[DF99\] find the following commutator relations
$$\begin{array}{cc}& [_x^\mu \left(A_\mu \right)_{\mathrm{int}}^{g_{QED}}\left(x\right),\left(A_\nu \right)_{\mathrm{int}}^{g_{QED}}\left(y\right)]_{}=i^\nu D(xy)\hfill \\ \hfill \text{and}& [_x^\mu \left(A_\mu \right)_{\mathrm{int}}^{g_{QED}}\left(x\right),\left(\psi \right)_{\mathrm{int}}^{g_{QED}}\left(y\right)]_{}=g(x)D(xy)\left(\psi \right)_{\mathrm{int}}^g\left(y\right)\hfill \end{array}$$
(7.17)
if $`x,y𝒪`$. We can use the equation for the interacting BRS current to find the explicit form of the interacting BRS transformations, for example
$$\begin{array}{c}\hfill \begin{array}{cc}& \stackrel{~}{s}\left(\left(A_\mu \right)_{\mathrm{int}}^{g_{QED}}\left(x\right)\right)=i_\mu ^xu(x)\hfill \\ & \stackrel{~}{s}\left(_x^\mu \left(A_\mu \right)_{\mathrm{int}}^{g_{QED}}\left(x\right)\right)=0\hfill \\ & \stackrel{~}{s}\left(\left(\psi \right)_{\mathrm{int}}^{g_{QED}}\left(x\right)\right)=g(x)\left(\psi \right)_{\mathrm{int}}^{g_{QED}}\left(x\right)u(x)\hfill \\ & \stackrel{~}{s}\left(\left(\overline{\psi }\right)_{\mathrm{int}}^{g_{QED}}\left(x\right)\right)=g(x)\left(\overline{\psi }\right)_{\mathrm{int}}^{g_{QED}}\left(x\right)u(x)\hfill \end{array}\begin{array}{cc}& \stackrel{~}{s}\left(u(x)\right)=0\hfill \\ & \stackrel{~}{s}\left(\stackrel{~}{u}(x)\right)=i_x^\rho \left(A_\rho \right)_{\mathrm{int}}^{g_{QED}}\left(x\right)\hfill \\ & \stackrel{~}{s}\left(\left(F^{\mu \nu }\right)_{\mathrm{int}}^{g_{QED}}\left(x\right)\right)=0\hfill \\ & \stackrel{~}{s}\left(\left(j_{\mathrm{el}}^\mu \right)_{\mathrm{int}}^{g_{QED}}\left(x\right)\right)=0,\hfill \end{array}\end{array}$$
(7.18)
for $`x𝒪`$. The interacting electric current and the interacting field strength tensor are the only nontrivial observable quantities of those. The other two quantities with vanishing BRS transformation are not observable. The ghost $`u(x)`$ has non vanishing ghost number, and $`_x^\mu \left(A_\mu \right)_{\mathrm{int}}^{g_{QED}}\left(x\right)`$ is a coboundary and therefore equivalent to zero.
### 7.2. Yang-Mills-theory
The basic fields in Yang-Mills theory<sup>20</sup><sup>20</sup>20We consider here only pure, massless Yang-Mills theory, for simplicity are Lie algebra valued vector bosons $`A_\mu =A_\mu ^a\tau _a`$, ghosts $`u=u^a\tau _a`$ and anti-ghosts $`\stackrel{~}{u}=\stackrel{~}{u}^a\tau _a`$. The $`\tau _a`$ form a basis of the Lie algebra. Their Lie-bracket gives $`[\tau _a,\tau _b]=f_{ab}^c\tau _c`$. The $`f_{ab}^c`$ are the structure constants of the Lie-algebra. They satisfy the Jacobi-identity
$$f_{ab}^ef_{ec}^d+f_{bc}^ef_{ea}^d+f_{ca}^ef_{eb}^d=0$$
(7.19)
and are assumed to be totally antisymmetric.
The free field operators that belong to different components $`A_\mu ^a,u^a,\stackrel{~}{u}^a`$ of the fields $`A_\mu ,u`$ and $`\stackrel{~}{u}`$ have trivial commutation relations among each other, e.g.
$$\{u^a(x),\stackrel{~}{u}^b(y)\}_+=i\delta ^{ab}D(xy).$$
(7.20)
Therefore the free model underlying Yang-Mills theory is simply a $`p`$-fold copy of free QED if $`p`$ is the dimension of the Lie algebra. The underlying free model was considered in section (3.5)
The Lagrangian of Yang-Mills theory in causal perturbation theory is
$$_{YM}=\frac{1}{2}f_{ab}^cA_\mu ^aA_\nu ^bF_c^{\nu \mu }f_{ab}^cA_\mu ^bu^a^\mu \stackrel{~}{u}_c.$$
(7.21)
Here $`F_c^{\mu \nu }\stackrel{\text{def}}{=}(A_c^\nu )^{(1,\mu )}(A_c^\mu )^{(1,\nu )}`$. Note that there is no four-gluon-vertex present. It is created in second order perturbation theory due to $`C_{A,1}=\frac{1}{2}`$, see \[DHKS94a\] \- \[DHS95b\] for further details.
For the interacting fields we get
$$\begin{array}{cc}& \left(A_\mu ^a\right)_{\mathrm{int}}^{g_{YM}}\left(x\right)=\left(\left(A_\mu ^a\right)_{\mathrm{int}}^{g_{YM}}\left(x\right)\right)^{}\stackrel{~}{\text{ }\mathrm{C}}_1(𝒟),\hfill \\ & \left(u^a\right)_{\mathrm{int}}^{g_{YM}}\left(x\right)=\left(\left(u^a\right)_{\mathrm{int}}^{g_{YM}}\left(x\right)\right)^{}\stackrel{~}{\text{ }\mathrm{C}}_1(𝒟),\hfill \\ & \left(\stackrel{~}{u}^a\right)_{\mathrm{int}}^{g_{YM}}\left(x\right)=\left(\left(\stackrel{~}{u}^a\right)_{\mathrm{int}}^{g_{YM}}\left(x\right)\right)^{}\stackrel{~}{\text{ }\mathrm{C}}_1(𝒟).\hfill \end{array}$$
(7.22)
From eqn. (6.5) we get for the higher generators
$$\begin{array}{cc}& \left((A_\mu ^a)^{(1,\nu )}\right)_{\mathrm{int}}^{g_{YM}}\left(x\right)=_x^\nu \left(A_\mu ^a\right)_{\mathrm{int}}^{g_{YM}}\left(x\right)\frac{1}{2}g(x)\left(f_{bc}^aA_\mu ^bA_\nu ^c\right)_{\mathrm{int}}^{g_{YM}}\left(x\right),\hfill \\ & \left((u^a)^{(1,\nu )}\right)_{\mathrm{int}}^{g_{YM}}\left(x\right)=_x^\nu \left(u^a\right)_{\mathrm{int}}^{g_{YM}}\left(x\right)+g(x)\left(f_{bc}^aA^{b,\nu }u^c\right)_{\mathrm{int}}^{g_{YM}}\left(x\right).\hfill \\ & \left((\stackrel{~}{u}^a)^{(1,\nu )}\right)_{\mathrm{int}}^{g_{YM}}\left(x\right)=_x^\nu \left(\stackrel{~}{u}^a\right)_{\mathrm{int}}^{g_{YM}}\left(x\right).\hfill \end{array}$$
(7.23)
The first equation implies in particular
$$\left(F_{\mu \nu }^a\right)_{\mathrm{int}}^{g_{YM}}\left(x\right)=_\mu ^x\left(A_\nu ^a\right)_{\mathrm{int}}^{g_{YM}}\left(x\right)_\nu ^x\left(A_\mu ^a\right)_{\mathrm{int}}^{g_{YM}}\left(x\right)+g(x)\left(f_{bc}^aA_\mu ^bA_\nu ^c\right)_{\mathrm{int}}^{g_{YM}}\left(x\right)$$
(7.24)
and
$$\left((A_\mu ^a)^{(1,\mu )}\right)_{\mathrm{int}}^{g_{YM}}\left(x\right)=_x^\mu \left(A_\mu ^a\right)_{\mathrm{int}}^{g_{YM}}\left(x\right).$$
(7.25)
The first relation reproduces the usual relation between the interacting vector boson field and the field strength tensor in Yang-Mills theories. From the Lagrangian we can also derive the field equations using eqn. (6.7):
$$\begin{array}{cc}& \begin{array}{cc}\hfill \mathrm{}^x\left(A_\mu ^a\right)_{\mathrm{int}}^{g_{YM}}\left(x\right)=& _x^\nu \left[g(x)\left(f_{bc}^aA_\mu ^bA_\nu ^c\right)_{\mathrm{int}}^{g_{YM}}\left(x\right)\right]\hfill \\ & g(x)\left(f_{bc}^aA^{\nu ,b}F_{\nu \mu }^c\right)_{\mathrm{int}}^{g_{YM}}\left(x\right)+g(x)\left(f_{bc}^au^b(\stackrel{~}{u}^c)_{(1,\mu )}\right)_{\mathrm{int}}^{g_{YM}}\left(x\right),\hfill \end{array}\hfill \\ & \mathrm{}^x\left(u^a\right)_{\mathrm{int}}^{g_{YM}}\left(x\right)=_x^\mu \left[g(x)\left(f_{bc}^aA^{\mu ,b}u^c\right)_{\mathrm{int}}^{g_{YM}}\left(x\right)\right],\hfill \\ & \mathrm{}^x\left(\stackrel{~}{u}^a\right)_{\mathrm{int}}^{g_{YM}}\left(x\right)=g(x)\left(f_{bc}^aA_\mu ^b(\stackrel{~}{u}^c)^{(1,\mu )}\right)_{\mathrm{int}}^{g_{YM}}\left(x\right).\hfill \end{array}$$
(7.26)
The Lagrangian (7.21) obviously satisfies the conditions (C1), (C2),(C3) and (C5). It is also possible to find $`Q(n)`$-vertices and $`R(n)`$-vertices such that condition (C4) is valid. These vertices are
$$\begin{array}{cc}& _1^\mu =f_{ab}^cu^aA_\nu ^bF_c^{\nu \mu }\frac{1}{2}f_{ab}^cu^au^b(\stackrel{~}{u}_c)^{(1,\mu )},\hfill \\ & _2^{\mu \rho }=\frac{1}{2}f_{ab}^cu^au^bF_c^{\mu \rho },\hfill \\ & _i^{\mathrm{}}=_i^{\mathrm{}}=0i>2.\hfill \end{array}\begin{array}{cc}& _1^\mu =\frac{1}{2}f_{ab}^cu^au^b(\stackrel{~}{u}_c)^{(1,\mu )},\hfill \\ & _2^{\mu \rho }=\frac{1}{2}f_{ab}^cu^au^bF_c^{\mu \rho },\hfill \end{array}$$
(7.27)
The vertices have been chosen such that condition (N6) is compatible with all other normalization conditions in first order.
We remind the reader that the existence of solutions for condition (N6) has not been proven for an arbitrary number of arguments in Yang-Mills theory. Scharf and collaborators, \[DHKS94a\] \- \[DHS95b\], have proven operator gauge invariance, i.e. eqn. (4.35) for $`i_1=\mathrm{}=i_n=0`$, see section (4.4) for further details.
We already mentioned that the free model underlying Yang-Mills theory is a $`p`$-fold copy of free QED if $`p`$ is the dimension of the Lie-group. The question arises whether there are other interactions besides Yang-Mills theory with the same free model. Stora \[Sto97\] found out that the number of possible Lagrangians for such a model is severely restricted by the conditions (C1) - (C5). Lagrangians $`T\left(\right)(x)`$ may differ from the Yang-Mills Lagrangian only by a coboundary $`s_0T\left(K\right)(x)`$ or a derivative $`_\mu T\left(K^\mu \right)(x)`$ where $`K`$ is a scalar polynomial with ghost number $`1`$ and $`K^\mu `$ is a vector polynomial with ghost number zero. By Yang-Mills Lagrangian we mean here an expression like (7.21) with arbitrary constants $`f_{bc}^a`$ that are totally antisymmetric in their indices and satisfy the Jacobi identity (7.19). In particular the Lie-group structure needs not to be put in. The Jacobi identity for the constants $`f_{bc}^a`$ is a consequance of operator gauge invariance in second order, and operator gauge invariance in first order implies that they are totally antisymmetric. Since Stora’s paper is not published, we refer the reader to the articles of Aste and Scharf \[AS98\] and Grigore \[Gri98\].
It is usually argued that the addition of such coboundary or derivative terms does not change the model because coboundaries are equivalent to zero in cohomology and derivatives should not give a contribution in the adiabatic limit. Dütsch \[Düt96\] has proven that this is correct also in higher orders for theories where the adiabatic limit can be performed, e.g. in massive theories. But for models where this limit does not exist the question is still open. Concerning the coboundary terms we remark that it is not clear whether a coboundary in the free theory, $`T\left(A\right)=s_0T\left(B\right)`$ for some $`B𝒫`$, gives a coboundary in the interacting theory, such that $`\left(A\right)_{\mathrm{int}}^g\left(x\right)=\stackrel{~}{s}\left(C\right)_{\mathrm{int}}^g\left(x\right)`$ for some $`C𝒫`$. Direct calculations in first order indicate that this is indeed true for suitable normalizations, but as long as this question is not clarified coboundary terms in the Lagrangian must not be neglected. The same is true for derivated terms in these theories.
In the rest of the section we want to compare our results with those of Nakanishi and Ojima \[NO90\]. Their results have been derived in the context of quantum field theory, but they are also classical in the following sense: They use field equations derived as Euler-Lagrange equations from a classical action, and they deliberately neglect the distributional character of field operators and form products of field operators at the same spacetime point. Therefore it is possible to compare their results with the classical limit of our results. At first we note that Nakanishi and Ojima use a different convention for the anti-ghosts. Their ghosts $`C`$ and anti-ghosts $`\overline{C}`$ correspond to ours in the following way:
$$C^au^a,\overline{C}^ai\stackrel{~}{u}^a.$$
(7.28)
For the comparison we will always translate their results into our language. To make the notation shorter we introduce the covariant derivative of a field with Lie-algebra index, $`X(x)=X^a(x)\tau _a`$, as
$$(D_\mu X)^a(x)\stackrel{\text{def}}{=}_\mu ^xX^a(x)+f_{bc}^a\left(A_\mu ^b\right)_{\mathrm{cl}}\left(x\right)X^c(x).$$
(7.29)
We have the classical fields
$$C(A_\mu ^a)(x)=\left(A_\mu ^a\right)_{\mathrm{cl}}\left(x\right),C(u^a)(x)=\left(u^a\right)_{\mathrm{cl}}\left(x\right),C(\stackrel{~}{u}^a)(x)=\left(\stackrel{~}{u}^a\right)_{\mathrm{cl}}\left(x\right)$$
(7.30)
and for the higher generators the representation $`C`$ gives
$$\begin{array}{cc}& C\left((A_\mu ^a)^{(1,\nu )}\right)(x)=_x^\nu \left(A_\mu ^a\right)_{\mathrm{cl}}\left(x\right)\frac{1}{2}g(x)f_{bc}^a\left(A_\mu ^b\right)_{\mathrm{cl}}\left(x\right)\left(A_\nu ^c\right)_{\mathrm{cl}}\left(x\right),\hfill \\ & C\left((u^a)^{(1,\nu )}\right)(x)=_x^\nu \left(u^a\right)_{\mathrm{cl}}\left(x\right)+g(x)f_{bc}^a\left(A^{b,\nu }\right)_{\mathrm{cl}}\left(x\right)\left(u^c\right)_{\mathrm{cl}}\left(x\right).\hfill \\ & C\left((\stackrel{~}{u}^a)^{(1,\nu )}\right)(x)=_x^\nu \left(\stackrel{~}{u}^a\right)_{\mathrm{cl}}\left(x\right).\hfill \end{array}$$
(7.31)
The field equations (7.26) become in the classical limit
$$\begin{array}{cc}& \begin{array}{cc}\hfill (D^\mu F_{\mu \nu }^{\mathrm{cl}})^a(x)& =_\nu ^x_x^\mu \left(A_\mu ^a\right)_{\mathrm{cl}}\left(x\right)\hfill \\ & +gf_{bc}^a\left(_\nu ^x\left(\stackrel{~}{u}^b\right)_{\mathrm{cl}}\left(x\right)\right)\left(u^c\right)_{\mathrm{cl}}\left(x\right),\hfill \end{array}\hfill \\ & _x^\mu (D_\mu \left(u\right)_{\mathrm{cl}})^a(x)=0,\hfill \\ & (D_\mu ^\mu \left(u\right)_{\mathrm{cl}})^a(x)=0.\hfill \end{array}$$
(7.32)
Here $`F_{\mu \nu }^{a,\mathrm{cl}}`$ is the classical field strength tensor,
$$F_{\mu \nu }^{a,\mathrm{cl}}=_\mu ^x\left(A_\nu ^a\right)_{\mathrm{cl}}\left(x\right)_\nu ^x\left(A_\mu ^a\right)_{\mathrm{cl}}\left(x\right)+gf_{bc}^a\left(A_\mu ^b\right)_{\mathrm{cl}}\left(x\right)\left(A_\mu ^c\right)_{\mathrm{cl}}\left(x\right).$$
(7.33)
The field equations are exactly the same as those of Nakanishi and Ojima. For the ghost current we get
$$C\left(k^\mu \right)(x)=i\underset{a}{}\left(\left(u^a\right)_{\mathrm{cl}}\left(x\right)_x^\mu \left(\stackrel{~}{u}^a\right)_{\mathrm{cl}}\left(x\right)(D^\mu \left(u\right)_{\mathrm{cl}})^a(x)\left(\stackrel{~}{u}^a\right)_{\mathrm{cl}}\left(x\right)\right).$$
(7.34)
This is $`i`$ times the result of Nakanishi and Ojima. The factor $`i`$ comes from a different definition of the ghost current. They require that the ghost current and -charge be pseudo-hermitian, so that the eigenvalues of the ghost charge are in $`i\mathrm{𝖹𝖹}`$.
For the classical BRS current $`j_B^\mu (x)`$ we have according to definition (6.26)
$$j_B^\mu (x)=\left(j_B^\mu \right)_{\mathrm{cl}}\left(x\right)g\left(_1^\mu \right)_{\mathrm{cl}}\left(x\right).$$
(7.35)
This reads in terms of the basic fields
$$\begin{array}{cc}\hfill j_B^\mu (x)& =\underset{a}{}\left((D^\mu \left(u\right)_{\mathrm{cl}})^a(x)_x^\nu \left(A_\nu ^a\right)_{\mathrm{cl}}\left(x\right)\left(u^a\right)_{\mathrm{cl}}\left(x\right)_x^\mu _x^\nu \left(A_\nu ^a\right)_{\mathrm{cl}}\left(x\right)\right)\hfill \\ & \frac{1}{2}f_{ab}^c\left(u^a\right)_{\mathrm{cl}}\left(x\right)\left(u^b\right)_{\mathrm{cl}}\left(x\right)_x^\mu \left(\stackrel{~}{u}_c\right)_{\mathrm{cl}}\left(x\right)\hfill \end{array}$$
(7.36)
This is again — up to a minus sign which is pure convention — the same result as Nakanishi and Ojima. Therefore we realize a complete agreement between the results of Nakanishi and Ojima and ours, apart from different conventions. This supports both our results at the quantum level and also the correspondence law. The same relations at the quantum level would have given different results if we had adopted the correspondence law (6.53), for example.
## 8. Conclusions and Outlook
We presented a universal construction of local quantum gauge theories. It gives an algebra of local observables that has a Hilbert space representation. For this construction to work two preconditions must hold: The underlying free theory must be positive in the sense discussed in chapter (3), and the time ordered products of free field operators must satisfy the conditions (N1) - (N6). The second precondition can only be violated with respect to condition (N6), all other conditions can always be accomplished. If all the normalization conditions hold, a locally conserved BRS current and with it a nilpotent BRS transformation on the algebra of local fields can be defined. The algebra of local observables is then defined as the cohomology of the algebra of local fields w.r.t. the BRS transformation. If the underlying free model is positive, the Hilbert space representation can be constructed. Therefore spacetime must be compactified spatially in order to allow a nilpotent BRS charge to be defined. This compactification does not change the algebra. It is an open question whether two representations that are constructed with a different compactification length are equivalent or not.
The most crucial point for each model that is investigated in this framework is whether normalization condition (N6) can be accomplished together with the other normalization conditions. We have proven that this holds fod quantum electrodynamics, but for Yang-Mills theory the question is still open. We think that methods of algebraic renormalization can help to find a solution. To reduce the problem to an algebraic one it could be helpful to define a BRS transformation $`s`$ on the algebra $`𝒫`$, such that $`T(sA)=s_0T(A)A𝒫`$. This requires the introduction of an additional auxiliary field, the scalar Nakanishi-Lautrup field $`B𝒫`$ with the properties $`s\stackrel{~}{u}=iB,sB=0`$ and $`T(B)(x)=^\mu A_\mu (x)`$. With this definition the BRS transformation $`s`$ on $`𝒫`$ can be chosen to be nilpotent. These notions could make it possible to translate the language of algebraic renormalization into ours. Normalization condition (N6) takes on the form of the descent equations in algebraic renormalization. Since they are proven in Yang-Mills theory, this could also lead to a proof of (N6) for Yang-Mills theory.
The renormalization scheme underlying our construction is the one of Epstein and Glaser. It is formulated, unlike the other renormalization schemes, in configuration space. Therefore it is suitable for quantum field theories on curved spacetimes. Brunetti and Fredenhagen \[BF99\] have shown that the time ordered products can also be defined in globally hyperbolic spacetimes. To generalize our normalization conditions to these spacetimes, the propagators and differential operators introduced in chapter (3) must be substituted by suitably generalized ones. With the normalization conditions all relations derived from them carry over to curved spacetimes, in particular the field equations, the conservation of ghost and BRS current and the nilpotency of the BRS charge and the BRS transformation. So it is possible to define an algebra of local observables even in globally hyperbolic spacetimes, provided these spacetimes allow propagators and their corresponding differential operators to be defined.
Acknowledgement: I thank Prof. Klaus Fredenhagen for entrusting me with this interesting topic, for guidance and constant support during the course of this work. Detailed discussions with him were always a great motivation and help. He also managed to create an inspiring atmosphere in the research group.
I am also greatful to Michael Dütsch. I benefited a lot from his great experience withcausal perturbation theory, many valuable discussions and fruitful collaboration.
Thanks are also due to the other colleagues in our research group for many interesting discussions on various topics.
Financial support given by the DFG as part of the ‘Graduiertenkolleg für theoretische Elementarteilchenphysik’ is greatfully acknowledged.
## Appendix A Proof of (N3) and (N4)
### A.1. Proof of (N3)
The essential point in the proof that solutions for condition (N3) exist is to show that eqn. (N3) is equivalent to the causal Wick expansion (4.24). This suffices for a proof because it was already shown in \[BF99\] that (4.24) has solutions, see below.
The proof that both conditions are equivalent for a certain $`T(W_1,\mathrm{},W_n)`$ proceeds inductively. The induction hypothesis is that eqn. (N3) and eqn. (4.24) hold and are equivalent for the following time ordered products: all time ordered products that contain fewer arguments than $`n`$ and all that contain a combination of sub monomials of the $`W_i`$, if at least one of these sub monomials is a proper one.
At first we prove that eqn. (4.24) implies eqn. (N3). With eqn. (4.24) the time ordered product on the left hand side of (N3) can be written as
$$\begin{array}{cc}\hfill T(W_1,\mathrm{},W_n)(x_1,\mathrm{},x_n)=\underset{\gamma _1,\mathrm{},\gamma _n}{}\omega _0\left(T(W_1^{(\gamma _1)},\mathrm{},W_n^{(\gamma _n)})(x_1,\mathrm{},x_n)\right)& \\ \hfill \times \frac{:\phi ^{\gamma _1}(x_1)\mathrm{}\phi ^{\gamma _n}(x_n):}{\gamma _1!\mathrm{}\gamma _n!}& ,\hfill \end{array}$$
(A.1)
for the notation see the formulas following (4.24). To calculate the (anti-) commutator with the $`\phi _i(z)`$ in eqn. (N3), we note that the (anti-) commutator of the Wick product with the $`\phi _i(z)`$ gives
$$\begin{array}{cc}& [\frac{:\phi ^{\gamma _1}(x_1)\mathrm{}\phi ^{\gamma _n}(x_n):}{\gamma _1!\mathrm{}\gamma _n!},\phi _i(z)]_{}=\hfill \\ & =i\underset{k=1}{\overset{n}{}}\underset{j}{}\mathrm{\Delta }_{ij}(zx_k)\frac{:\phi ^{\gamma _1}(x_1)\mathrm{}\phi ^{\gamma _ke_j}(x_k)\mathrm{}\phi ^{\gamma _n}(x_n):}{\gamma _1!\mathrm{}(\gamma _ke_j)!\mathrm{}\gamma _n!}\hfill \end{array}$$
(A.2)
if the $`\gamma _k0`$, otherwise the respective term vanishes. Here $`e_j`$ is the unit vector with an entry $`1`$ at the $`j^{\mathrm{th}}`$ position and the other entries zero. Therefore we get for the complete commutator
$$\begin{array}{cc}\hfill i\underset{k=1}{\overset{n}{}}\underset{j}{}\mathrm{\Delta }_{ij}(zx_k)\underset{\stackrel{\gamma _1,\mathrm{},\gamma _n}{\gamma _k0}}{}\omega _0\left(T(W_1^{(\gamma _1)},\mathrm{},W_n^{(\gamma _n)})(x_1,\mathrm{},x_n)\right)& \\ \hfill \times \left[\frac{:\phi ^{\gamma _1}(x_1)\mathrm{}\phi ^{\gamma _ke_j}(x_k)\mathrm{}\phi ^{\gamma _n}(x_n):}{\gamma _1!\mathrm{}(\gamma _ke_j)!\mathrm{}\gamma _n!}\right].& \end{array}$$
(A.3)
This becomes after a shifting of indices
$$\begin{array}{cc}& \begin{array}{cc}& i\underset{k=1}{\overset{n}{}}\underset{j}{}\mathrm{\Delta }_{ij}(zx_k)\times \hfill \\ & \begin{array}{cc}\hfill \times \underset{\gamma _1,\mathrm{},\gamma _n}{}\omega _0\left(T(W_1^{(\gamma _1)},\mathrm{},W_k^{(\gamma _1+e_j)},\mathrm{},W_n^{(\gamma _n)})(x_1,\mathrm{},x_n)\right)\times & \\ \hfill \times \left[\frac{:\phi ^{\gamma _1}(x_1)\mathrm{}\phi ^{\gamma _n}(x_n):}{\gamma _1!\mathrm{}\gamma _n!}\right]& \end{array}\hfill \end{array}\hfill \\ & =i\underset{k=1}{\overset{n}{}}\underset{j}{}\mathrm{\Delta }_{ij}(zx_k)T(W_1,\mathrm{},W_k^{(e_j)},\mathrm{},W_n)(x_1,\mathrm{},x_n).\hfill \end{array}$$
(A.4)
The last identity is valid because eqn. (A.1) holds for the sub monomials according to our induction hypothesis. This proves that (N3) is a consequence of (A.1).
To complete the proof of equivalence we recall that we already saw that eqn. (N3) determines the time ordered product up to a $`\mathrm{C}`$-number distribution. To be precise, eqn. (N3) determines completely
$$T(W_1,\mathrm{},W_n)\omega _0\left(T(W_1,\mathrm{},W_n)\right)1\mathrm{l}$$
(A.5)
and leaves
$$\omega _0\left(T(W_1,\mathrm{},W_n)\right)$$
(A.6)
open. This is exactly the same with (A.1). The Wick products are determined anyway and the numerical distributions are determined by the $`T`$-products for the sub monomials if at least one $`\gamma _i0`$. Since both equations determine the same part of the distribution and leave the same part open, and moreover one of them is a consequence of the other, they must be equivalent.
The question arises whether the expression on the right hand side of eqn. (A.1) is well defined, because there appear products of distribution. The answer is the same as in section (4.2): Epstein and Glaser’s “Theorem 0” guarantees that the product is well defined.
### A.2. Proof of (N4)
Like for (N3) we do not prove the existence of solutions for (N4) itself but for its integrated version
$$\begin{array}{cc}& T(W_1,\mathrm{},W_n,\phi _i)(x_1,\mathrm{},x_n,y)=\hfill \\ & =i\underset{k=1}{\overset{n}{}}\underset{j}{}\mathrm{\Delta }_{ij}^F(yx_k)T(W_1,\mathrm{},\frac{W_k}{\phi _j},\mathrm{},W_n)(x_1,\mathrm{},x_n)\hfill \\ & +\underset{\gamma _1\mathrm{}\gamma _n}{}\omega _0\left(T(W_1^{(\gamma _1)},\mathrm{},W_n^{(\gamma _n)})(x_1,\mathrm{},x_n)\right)\frac{:\phi ^{\gamma _1}(x_1)\mathrm{}\phi ^{\gamma _n}(x_n)\phi _i(y):}{\gamma _1!\mathrm{}\gamma _n!}.\hfill \end{array}$$
(A.7)
At the end of the section we will prove that the two conditions are equivalent.
The right hand side of eqn. (A.7) is obviously well defined, because the first sum is a tensor product of distributions which is always well defined — the argument $`y`$ does not appear in the time ordered product — while the second sum is simply part of (A.1) which was already proven to be well defined.
The question is whether this expression has the correct causal factorization outside the diagonal. To show this we proceed again inductively, the induction hypothesis is that eqn. (A.7) is valid for all time ordered products of sub monomials of the $`W_i`$.
At first we compare the expression with (A.1), which reveals in the present case
$$\begin{array}{cc}& T(W_1,\mathrm{},W_n,\phi _i)(x_1,\mathrm{},x_n,y)=\hfill \\ & =\underset{\gamma _1\mathrm{}\gamma _n}{}\omega _0\left(T(W_1^{(\gamma _1)},\mathrm{},W_n^{(\gamma _n)},\phi _i)(x_1,\mathrm{},x_n,y)\right)\frac{:\phi ^{\gamma _1}(x_1)\mathrm{}\phi ^{\gamma _n}(x_n):}{\gamma _1!\mathrm{}\gamma _n!}\hfill \\ & +\underset{\gamma _1\mathrm{}\gamma _n}{}\omega _0\left(T(W_1^{(\gamma _1)},\mathrm{},W_n^{(\gamma _n)})(x_1,\mathrm{},x_n)\right)\frac{:\phi ^{\gamma _1}(x_1)\mathrm{}\phi ^{\gamma _n}(x_n)\phi _i(y):}{\gamma _1!\mathrm{}\gamma _n!}.\hfill \end{array}$$
(A.8)
So the second sum in eqn. (A.7) is already present and we must only show that
$$i\underset{k=1}{\overset{n}{}}\underset{j}{}\mathrm{\Delta }_{ij}^F(yx_k)T(W_1,\mathrm{},\frac{W_k}{\phi _j},\mathrm{},W_n)(x_1,\mathrm{},x_n)$$
(A.9)
is a possible extension of
$$\underset{\gamma _1\mathrm{}\gamma _n}{}\omega _0\left(T^0(W_1^{(\gamma _1)},\mathrm{},W_n^{(\gamma _n)},\phi _i)(x_1,\mathrm{},x_n,y)\right)\frac{:\phi ^{\gamma _1}(x_1)\mathrm{}\phi ^{\gamma _n}(x_n):}{\gamma _1!\mathrm{}\gamma _n!}$$
(A.10)
to the diagonal. Inserting eqn. (N3) into expression (A.9) gives
$$\begin{array}{cc}& i\underset{k=1}{\overset{n}{}}\underset{j}{}\mathrm{\Delta }_{ij}^F(yx_k)\times \hfill \\ & \begin{array}{cc}\hfill [\underset{\gamma _1\mathrm{}\gamma _n}{}\omega _0\left(T(W_1^{(\gamma _1)},\mathrm{},W_k^{(\gamma _k+e_j)},\mathrm{},W_n^{(\gamma _n)})(x_1,\mathrm{},x_n)\right)& \\ \hfill \times \frac{:\phi ^{\gamma _1}(x_1)\mathrm{}\phi ^{\gamma _n}(x_n):}{\gamma _1!\mathrm{}\gamma _n!}]& .\hfill \end{array}\hfill \end{array}$$
(A.11)
The latter is equal to expression (A.10) if
$$\begin{array}{cc}\hfill \begin{array}{cc}& i\underset{k=1}{\overset{n}{}}\underset{j}{}\mathrm{\Delta }_{ij}^F(yx_k)\times \hfill \\ & \left[\underset{\gamma _1\mathrm{}\gamma _n}{}\omega _0\left(T(W_1^{(\gamma _1)},\mathrm{},W_k^{(\gamma _k+e_j)},\mathrm{},W_n^{(\gamma _n)})(x_1,\mathrm{},x_n)\right)\right]\hfill \end{array}& \\ \hfill =\omega _0\left(T^0(W_1^{(\gamma _1)},\mathrm{},W_n^{(\gamma _n)},\phi _i)(x_1,\mathrm{},x_n,y)\right)& \end{array}$$
(A.12)
for all $`\gamma _1,\mathrm{},\gamma _n`$ and outside the diagonal. This equation is obviously true if at least one $`\gamma _i0`$ since eqn. (A.7) is valid for the sub monomials of the $`W_i`$ according to our induction hypothesis. So eqn. (A.7) can be accomplished if
$$i\underset{k=1}{\overset{n}{}}\underset{j}{}\mathrm{\Delta }_{ij}^F(yx_k)\omega _0\left(T(W_1,\mathrm{},W_k^{(e_j)},\mathrm{},W_n)(x_1,\mathrm{},x_n)\right)$$
(A.13)
is a possible extension of
$$\omega _0\left(T^0(W_1,\mathrm{},W_n,\phi _i)(x_1,\mathrm{},x_n,y)\right).$$
(A.14)
To see this we smear out both expressions with a test function $`\eta `$ that vanishes with all its derivatives on the diagonal $`\mathrm{\Delta }_{n+1}`$. Let the test function $`\eta `$ be fix and recall the definition of the partition of unity in eqn. (4.18). Then we can define for each subset $`Z\{x_1,\mathrm{},x_n,y\}`$ another test function $`\eta _Z𝒟(M^{n+1})`$
$$\eta _Z\stackrel{\text{def}}{=}\{\begin{array}{cc}f_Z\eta \hfill & \text{outside }\mathrm{\Delta }_{n+1}\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}$$
(A.15)
such that
$$\eta _Z\mathrm{}_Z\text{and}\underset{Z}{}\eta _Z=\eta .$$
(A.16)
Then the following equation is valid owing to causal factorization:
$$\begin{array}{cc}& d^4yd^4x_1\mathrm{}d^4x_n\eta (x_1,\mathrm{},x_n,y)T(W_1,\mathrm{},W_n,\varphi _i)(x_1,\mathrm{},x_n,y)=\hfill \\ & =\underset{ZX}{}d^4yd^4x_1\mathrm{}d^4x_n\eta _Z(x_1,\mathrm{},x_n,y)T\left(W_Z\right)(x_Z)T(W_{Z^c},\phi _i)(x_{Z^c},y)\hfill \\ & +\underset{ZX}{}d^4yd^4x_1\mathrm{}d^4x_n\eta _Z(x_1,\mathrm{},x_n,y)T(W_Z,\phi _i)(x_Z,y)T\left(W_{Z^c}\right)(x_{Z^c})\hfill \end{array}$$
(A.17)
because $`ZZ^c`$ on $`\eta _Z`$. Here $`X=\{x_1,\mathrm{},x_n\}`$.
Let us investigate $`T\left(W_Z\right)(x_Z)T(W_{Z^c},\phi _i)(x_{Z^c},y)`$ and assume for simplicity that $`Z=\{x_{k+1},\mathrm{},x_n\}`$ and $`Z^c=\{x_1,\mathrm{},x_k\}`$. Due to the validity of eqn. (A.7) in lower orders we have
$$\begin{array}{cc}& T\left(W_Z\right)(x_Z)T(W_{Z^c},\phi _i)(x_{Z^c},y)=\hfill \\ & \begin{array}{cc}\hfill =i\underset{m=1}{\overset{k}{}}\underset{j}{}\mathrm{\Delta }_{ij}^F(yx_k)T\left(W_Z\right)(x_Z)T(W_1,\mathrm{},W_m^{(e_j)},\mathrm{},W_k)(x_{Z^c})& \\ \hfill +\underset{\gamma _1\mathrm{}\gamma _k}{}\omega _0\left(T(W_1^{(\gamma _1)},\mathrm{},W_k^{(\gamma _k)})(x_{Z^c})\right)T\left(W_Z\right)(x_Z)\times & \\ \hfill \times \frac{:\phi ^{\gamma _1}(x_1)\mathrm{}\phi ^{\gamma _k}(x_k)\phi _i(y):}{\gamma _1!\mathrm{}\gamma _k!}& .\hfill \end{array}\hfill \end{array}$$
(A.18)
Since $`ZZ^c`$, the product in the first sum recombines to
$$\begin{array}{cc}\hfill T\left(W_Z\right)(x_Z)T(W_1,\mathrm{},W_m^{(e_j)},\mathrm{},W_k)(x_{Z^c})& \\ \hfill =T(W_1,\mathrm{},W_m^{(e_j)},\mathrm{},W_n)(x_1,\mathrm{},x_n)& .\hfill \end{array}$$
(A.19)
For the product in the second sum we get
$$\begin{array}{cc}& T\left(W_Z\right)(x_Z)\frac{:\phi ^{\gamma _1}(x_1)\mathrm{}\phi ^{\gamma _k}(x_k)\phi _i(y):}{\gamma _1!\mathrm{}\gamma _k!}=\hfill \\ & \begin{array}{cc}\hfill :\mathrm{}\phi _i(y):+\underset{l=k+1}{\overset{n}{}}\underset{j}{}T(W_{k+1},\mathrm{},W_l^{(e_j)},\mathrm{},W_n)(x_Z)\mathrm{\Delta }_{ij}^+(yx_l)& \\ \hfill \times \frac{:\phi ^{\gamma _1}(x_1)\mathrm{}\phi ^{\gamma _k}(x_k):}{\gamma _1!\mathrm{}\gamma _k!}& .\hfill \end{array}\hfill \end{array}$$
(A.20)
Inserting this into eqn. (A.18) and taking (A.1) into account, the second sum becomes
$$\begin{array}{cc}\hfill i\underset{l=1}{\overset{k}{}}\underset{j}{}\mathrm{\Delta }_{ij}^+(yx_k)T(W_{k+1},\mathrm{},W_l^{(e_j)},\mathrm{},W_n)(x_Z)T\left(W_{Z^c}\right)(x_{Z^c})& \\ \hfill +:\mathrm{}\phi _i(y):& .\hfill \end{array}$$
(A.21)
From the definition of the Feynman propagator we find
$$\mathrm{\Delta }_{ij}^+(yx_l)=\mathrm{\Delta }_{ij}^F(yx_l)\mathrm{\Delta }_{ij}^A(yx_l)=\mathrm{\Delta }_{ij}^F(yx_l)$$
(A.22)
since $`(yx_l)\overline{V}_{}`$. Recombining the terms we finally arrive at
$$\begin{array}{cc}& \omega _0\left(T\left(W_Z\right)(x_Z)T(W_{Z^c},\phi _i)(x_{Z^c},y)\right)=\hfill \\ & =i\underset{m=1}{\overset{k}{}}\underset{j}{}\mathrm{\Delta }_{ij}^F(yx_k)\omega _0\left(T(W_1,\mathrm{},W_m^{(e_j)},\mathrm{},W_n)(x_1,\mathrm{},x_n)\right)\hfill \\ & +i\underset{l=k+1}{\overset{n}{}}\underset{j}{}\mathrm{\Delta }_{ij}^F(yx_l)\omega _0\left(T(W_1,\mathrm{},W_l^{(e_j)},\mathrm{},W_n)(x_1,\mathrm{},x_n)\right)\hfill \\ & =i\underset{m=1}{\overset{n}{}}\underset{j}{}\mathrm{\Delta }_{ij}^F(yx_k)\omega _0\left(T(W_1,\mathrm{},W_m^{(e_j)},\mathrm{},W_n)(x_1,\mathrm{},x_n)\right).\hfill \end{array}$$
(A.23)
With the same argument we can see that
$$\begin{array}{cc}& \omega _0\left(T(W_Z,\phi _i)(x_Z,y)T\left(W_{Z^c}\right)(x_{Z^c})\right)=\hfill \\ & =i\underset{m=1}{\overset{n}{}}\underset{j}{}\mathrm{\Delta }_{ij}^F(yx_k)\omega _0\left(T(W_1,\mathrm{},W_m^{(e_j)},\mathrm{},W_n)(x_1,\mathrm{},x_n)\right).\hfill \end{array}$$
(A.24)
Taking the vacuum expectation value of eqn. (A.17) and inserting the expressions above, we finally get
$$\begin{array}{cc}& d^4yd^4x_1\mathrm{}d^4x_n\eta (x_1,\mathrm{},x_n,y)\omega _0\left(T(W_1,\mathrm{},W_n,\varphi _i)(x_1,\mathrm{},x_n,y)\right)\hfill \\ & =id^4yd^4x_1\mathrm{}d^4x_n\eta (x_1,\mathrm{},x_n,y)\times \hfill \\ & \times \underset{m=1}{\overset{n}{}}\underset{j}{}\mathrm{\Delta }_{ij}^F(yx_k)\omega _0\left(T(W_1,\mathrm{},W_m^{(e_j)},\mathrm{},W_n)(x_1,\mathrm{},x_n)\right).\hfill \end{array}$$
(A.25)
So we have proven that expression (A.13) is a possible extension of (A.14), and this implies that eqn. (A.7) has the correct causal factorization. From the construction it is clear that (A.7) is compatible with (4.24) and thus with (N3). It is obvious that it respects the Poincaré transformation properties and is therefore compatible with (N1). The same calculation as for the compatibility of eqns. (N1) and (N2) reveals that it is also compatible with (N2).
Finally we have to prove that (N4) and (A.7) are equivalent. Eqn. (A.7) implies (N4) immediately: Application of the operator $`D^y`$, eqn. (3.86), from the left on eqn. (A.7) gives the desired result.
On the other hand a solution of (N4) is unique. This can best be seen for the corresponding equation for the retarded products,
$$\begin{array}{cc}& \underset{j}{}D_{ij}^yR(W_1,\mathrm{},W_n;\phi _j)(x_1,\mathrm{},x_n;y)=\hfill \\ & =i\underset{k=1}{\overset{n}{}}R(W_1,\mathrm{},\stackrel{ˇ}{k},\mathrm{},W_n;\frac{W_k}{\phi _i})(x_1,\mathrm{},\stackrel{ˇ}{k},\mathrm{},x_n;x_k)\delta (x_ky).\hfill \end{array}$$
(N4)
The difference of two solutions of this differential equation is a solution of the homogeneous differential equation. Due to the support properties of the retarded products there exists a Cauchy surface in the $`y`$-space such that all Cauchy data are zero. Therefore zero is a solution of that equation, and $`D^y`$ is an operator with a unique solution for the Cauchy problem, see page 3.97. So the retarded products are uniquely determined and with them the time ordered products. This completes the proof that (N4) and (A.7) are equivalent.
## Appendix B Proofs concerning the Ward identities
This appendix contains in its first section the proof that the ghost number Ward identities have common solutions with the other normalization conditions and that the ghost number Ward identities imply eqn. (4.32). In the second section we prove that the validity of the generalized operator gauge invariance already implies that there exists a solution of condition (N6).
### B.1. Proof of the ghost number Ward identities
We begin with the proof that the equation
$$\begin{array}{cc}& s_cT\left(W_1\mathrm{}W_n\right)(x_1,\mathrm{},x_n)=\hfill \\ & =\left(\underset{k=1}{\overset{n}{}}g(W_k)\right)T\left(W_1\mathrm{}W_n\right)(x_1,\mathrm{},x_n).\hfill \end{array}$$
(B.1)
is a direct consequence of condition (N5),
$$\begin{array}{cc}& _\mu ^yT(W_1,\mathrm{},W_n,k^\mu )(x_1,\mathrm{},x_n,y)=\hfill \\ & =\underset{k=1}{\overset{n}{}}g(W_k)\delta (yx_k)T(W_1,\mathrm{},W_n)(x_1,\mathrm{},x_n).\hfill \end{array}$$
(B.2)
Suppose, $`𝒪`$ is an open, bounded and causally complete region in spacetime such that all points $`x_1,\mathrm{},x_n`$ in eqn. (N5) lie in $`𝒪`$ — obviously for every set of points such a region can be found. Then we choose a test function $`f𝒟(M)`$ such that $`f(x)=1x𝒪^{}`$ with $`𝒪^{}`$ another open, bounded and causally complete region such that $`\overline{𝒪}𝒪^{}`$. Then we can find a Lorentz frame where a $`C^{\mathrm{}}`$-function $`H(y)`$ exists with the following properties:
$$\begin{array}{cc}& HC^{\mathrm{}}(M),H^tC^{\mathrm{}}(\mathrm{I}\mathrm{R}):H(y)=H^t(y^0),\hfill \\ & H^t(y^0)=1y^0<ϵ,H^t(y^0)=0y^0>ϵ,ϵ\mathrm{I}\mathrm{R},0<ϵ1,\hfill \\ & (H_\mu f)(\overline{V}_++𝒪)=\mathrm{},((1H)_\mu f)(\overline{V}_{}+𝒪)=\mathrm{}.\hfill \end{array}$$
(B.3)
The following calculations will be done in that Lorentz frame. Smearing out the left hand side of eqn. (N5) with $`f`$ gives
$$\begin{array}{cc}& d^4yf(y)_\mu ^yT(W_1,\mathrm{},W_n,k^\mu )(x_1,\mathrm{},x_n,y)=\hfill \\ & =d^4y(_\mu f)(y)H(y)T(W_1,\mathrm{},W_n,k^\mu )(x_1,\mathrm{},x_n,y)\hfill \\ & d^4y(_\mu f)(y)(1H(y))T(W_1,\mathrm{},W_n,k^\mu )(x_1,\mathrm{},x_n,y).\hfill \end{array}$$
(B.4)
According to our assumptions about the supports of the test functions $`(_\mu f)H(y)`$ and $`(_\mu f)(1H(y))`$ we have in the first integral on the right hand side $`yx_ii`$ and in the second integral on the right hand side $`x_iyi`$. Owing to causal factorization the time ordered product $`T(W_1,\mathrm{},W_n,k^\mu )(x_1,\mathrm{},x_n,y)`$ decomposes in the first integral according to $`T\left(k^\mu \right)(y)T(W_1,\mathrm{},W_n)(x_1,\mathrm{},x_n)`$ and in the second one according to $`T(W_1,\mathrm{},W_n)(x_1,\mathrm{},x_n)T\left(k^\mu \right)(y)`$. Therefore the integral can be written as
$$\begin{array}{cc}& d^4yf(y)_\mu ^yT(W_1,\mathrm{},W_n,k^\mu )(x_1,\mathrm{},x_n,y)=\hfill \\ & =d^4y(_\mu f)(y)H(y)[T\left(k^\mu \right)(y),T(W_1,\mathrm{},W_n)(x_1,\mathrm{},x_n)]_{}\hfill \\ & d^4y(_\mu f)(y)T(W_1,\mathrm{},W_n)(x_1,\mathrm{},x_n)T\left(k^\mu \right)(y).\hfill \end{array}$$
(B.5)
Then the second integral vanishes since $`k^\mu `$ is a conserved current. Partial integration in the first integral reveals according to the properties of $`H`$ and $`f`$
$$\begin{array}{cc}& d^4yf(y)_\mu ^yT(W_1,\mathrm{},W_n,k^\mu )(x_1,\mathrm{},x_n,y)=\hfill \\ & =[Q_c,T(W_1,\mathrm{},W_n)(x_1,\mathrm{},x_n)]_{}=s_cT(W_1,\mathrm{},W_n)(x_1,\mathrm{},x_n).\hfill \end{array}$$
(B.6)
As the smearing of the right hand side of eqn. (N5) with $`f`$ is trivial since $`f=1x_k`$, we finally arrive at
$$\begin{array}{cc}& s_cT(W_1,\mathrm{},W_n)(x_1,\mathrm{},x_n)=\hfill \\ & =\left(\underset{k=1}{\overset{n}{}}g(W_k)\right)T(W_1,\mathrm{},W_n)(x_1,\mathrm{},x_n).\hfill \end{array}$$
(B.7)
The proof that the ghost number Ward identities have common solution with the other normalization conditions proceeds along the same lines as the proof of Dütsch and Fredenhagen \[DF99\] for the electric current. An important difference between the proofs is that for their proof it suffices to have eqn. (N4) for the basic generators, while it is here important to have it also for the higher generators since the ghost current $`k^\mu `$ contains also higher generators.
The proof is subdivided into two parts. At first we prove that it is possible to normalize $`T\left(W_1\mathrm{}W_n\right)`$ such that it satisfies eqn. (B.1). Then we prove the same statement for condition (N5). This seems to be a detour because we just saw that (B.1) is a consequence of (N5), but (B.1) will be needed in the proof of (N5).
Like all these proofs this one goes by induction, so we put forward the induction hypothesis that both (N5) and (B.1) hold for fewer arguments than $`n`$ and for the sub monomials of the $`W_i`$, provided that at least one sub monomial is a proper one. Then the causal Wick expansion — eqn. (4.24) — tells us that eqn. (B.1) can only be violated by an unsuitable normalization of $`\omega _0\left(T(W_1,\mathrm{},W_n)\right)`$. Applying $`\omega _0`$ to eqn. (B.1) and taking $`\omega _0s_c=0`$ into account, we see that either $`\left(_{k=1}^ng(W_k)\right)=0`$ or $`\omega _0\left(T(W_1,\mathrm{},W_n)\right)=0`$. In the first case eqn. (B.1) is true for an arbitrary normalization of $`\omega _0\left(T(W_1,\mathrm{},W_n)\right)`$. In the second case validity of (B.1) in lower orders guarantees that $`\omega _0\left(T(W_1,\mathrm{},W_n)\right)`$ vanishes outside the diagonal but not necessarily on the entire $`M^n`$. Nevertheless it is always possible to extend a distribution that vanishes outside the diagonal by a distribution that vanishes everywhere, and such an extension is obviously compatible with all other normalization conditions. So it is always possible to find a normalization of $`T(W_1,\mathrm{},W_n)`$ that is a solution of all normalization conditions including (B.1).
Now we come to the second part, the proof that normalizations can be found for which the ghost number Ward identities (N5)
$$\begin{array}{cc}& _\mu ^yT(W_1(x_1),\mathrm{},W_n(x_n),k^\mu (y))=\hfill \\ & =\underset{k=1}{\overset{n}{}}\delta (yx_k)g(W_k)T(W_1(x_1),\mathrm{},W_n(x_n))\hfill \end{array}$$
(B.8)
hold such that the normalization is also in accordance with (N1) - (N4), provided none of the $`W_i`$ is equal to $`k^\mu `$ or contains it as a sub monomial, and none of them contains generators $`(u^a)^{(\alpha )}`$ or $`(\stackrel{~}{u}^a)^{(\alpha )}`$ with $`\left|\alpha \right|2`$.
To this end we define a possible anomaly as
$$\begin{array}{cc}\hfill a(x_1,\mathrm{},x_n,y)& =_\mu ^yT(W_1,\mathrm{},W_n,k^\mu )(x_1,\mathrm{},x_n,y)\hfill \\ & \underset{k=1}{\overset{n}{}}\delta (yx_k)g(W_k)T(W_1,\mathrm{},W_n)(x_1,\mathrm{},x_n)\hfill \end{array}$$
(B.9)
and show that a normalization can be found — in agreement with eqns. (N1) - (N4) — such that the anomaly vanishes. Recalling our induction hypothesis we want to show this under the assumption that all these anomalies vanish for the time ordered products of fewer arguments than $`n`$ and in all equations that involve the sub monomials of the $`W_i`$. The proof will be divided into three steps.
Step 1: At first we commute the anomaly with the basic fields $`\phi _i(x)`$ in order to find that this commutator vanishes. Thereby we make repeated use of condition (N3) and the fact that according to our induction hypothesis eqn. (N5) is already established for the lower orders and for the sub monomials. The result of that calculation is
$$\begin{array}{cc}& [a(x_1,\mathrm{},x_n,y),\phi _i(z)]_{}=\hfill \\ & \begin{array}{cc}\hfill =ig(\phi _i)\underset{k=1}{\overset{n}{}}\underset{j}{}\mathrm{\Delta }_{ij}(x_kz)\delta (yx_k)T(W_1,\mathrm{},\frac{W_k}{\phi _j},\mathrm{},W_n)& \\ \hfill +i\underset{j}{}\left(_\mu ^y\mathrm{\Delta }_{ij}(yz)\right)T(W_1,\mathrm{},W_n,\frac{k^\mu }{\phi _j})& \\ \hfill +i\underset{j}{}\mathrm{\Delta }_{ij}(yz)_\mu ^yT(W_1,\mathrm{},W_n,\frac{k^\mu }{\phi _j}),& \end{array}\hfill \end{array}$$
(B.10)
where we have omitted the spacetime arguments of the time ordered products because the expressions would not fit into the line otherwise. We will do this throughout this proof. It should not cause confusion since it is already clear from the arguments of the time ordered products which the spacetime arguments are.
To show that the expression above vanishes we distinguish three cases:
Case 1: $`\phi _i(z)u^a(z),\stackrel{~}{u}^a(z)`$. In this case both $`\frac{k^\mu }{\phi _i}=0`$ and $`g(\phi _i)=0`$, so the commutator vanishes immediately.
Case 2: $`\phi _i(z)=u^a(z)`$. At first we note that $`g(u^a)=1`$. Furthermore we have $`k^\mu =i(\stackrel{~}{u}_a)^{(1,\mu )}u^ai\stackrel{~}{u}_a(u^a)^{(1,\mu )}`$, so we get in particular $`\frac{k^\mu }{\stackrel{~}{u}_a}=i(u^a)^{(1,\mu )}`$ and $`\frac{k^\mu }{(\stackrel{~}{u}_a)^{(1,\nu )}}=i\delta _\nu ^\mu u^a`$. Taking this and the definition of the commutator function $`\mathrm{\Delta }_{ij}`$ into account, we get for the last two lines in (B.10)
$$\begin{array}{cc}& \left(_\mu ^yD(yz)\right)T(W_1,\mathrm{},W_n,(u^a)^{(1,\mu )})\hfill \\ & +D(yz)_\mu ^yT(W_1,\mathrm{},W_n,(u^a)^{(1,\mu )})\hfill \\ & \left(_\mu ^yD(yz)\right)_y^\mu T(W_1,\mathrm{},W_n,u^a).\hfill \end{array}$$
(B.11)
According to (N4) the expression above transforms into
$$\begin{array}{cc}& \left(_y^\mu D(yz)\right)\left[iC_{u,1}\underset{k=1}{\overset{n}{}}\delta (yx_k)T(W_1,\mathrm{},\frac{W_k}{(\stackrel{~}{u}_a)^{(1,\mu )}},\mathrm{},W_n)\right]\hfill \\ & +D(yz)[+iC_{u,1}\underset{k=1}{\overset{n}{}}\delta (yx_k)T(W_1,\mathrm{},\frac{W_k}{\stackrel{~}{u}_a},\mathrm{},W_n)\hfill \\ & (1+C_{u,1})\mathrm{}T(W_1,\mathrm{},W_n,u^a)\hfill \\ & i\frac{C_{u,1}}{C_{u,2}}_y^\alpha _y^\beta \left[\underset{k=1}{\overset{n}{}}\delta (x_ky)T(W_1,\mathrm{},\frac{W_k}{(\stackrel{~}{u}^a)^{2,\alpha \beta }},\mathrm{},W_n)\right]+\mathrm{}]\hfill \end{array}$$
(B.12)
Since we required that the $`W_i`$ do not contain generators $`(u^a)^{(\alpha )}`$ with $`\left|\alpha \right|2`$, the last line and the following terms containing derivatives w.r.t. higher generators on the $`W_i`$, indicated by the dots, vanish. Then comparing the remaining expression with the first line in (B.10) reveals that these expressions cancel each other if and only if $`C_{u,1}=1`$. So the choice $`C_{u,1}=1`$ is a necessary (and, as it will turn out, sufficient) condition for eqn. (N5) to hold.
Case 3: $`\phi _i(z)=\stackrel{~}{u}^a(z)`$. The calculation for this case is completely analogous the the one before and reveals $`C_{u,1}=1`$ as a necessary condition for the commutator to vanish, too.
So with $`C_u=1`$ the commutator of the anomaly with every free field vanishes. Consequently $`a(x_1,\mathrm{},x_n,y)`$, smeared with an arbitrary test function, is a $`\mathrm{C}`$-number distribution.
Step 2: We already know that time ordered products with at least one generator as an argument are completely determined by the time ordered products in lower orders and those for the sub monomials. We will now examine whether this normalization is compatible with (N5).
Since the anomaly can at most be a $`\mathrm{C}`$-number distribution, it is sufficient to calculate its $`\mathrm{C}`$-number part $`\omega _0(a(x_1,\mathrm{},x_n,y))`$. So we want to prove that
$$\begin{array}{cc}& _\mu ^y\omega _0\left(T(W_1,\mathrm{},W_n,\phi _i,k^\mu )\right)=\hfill \\ & \left(\underset{k=1}{\overset{n}{}}\delta (yx_k)g(W_k)+\delta (yz)g(\phi _i)\right)\omega _0\left(T(W_1,\mathrm{},W_n,\phi _i)\right).\hfill \end{array}$$
(B.13)
With a repeated use of eqn. (4.29) this can be transformed into
$$\begin{array}{cc}& ig(\phi _i)\underset{k=1}{\overset{n}{}}\underset{j}{}\mathrm{\Delta }_{ij}^F(zx_k)\delta (yx_k)\omega _0\left(T(W_1,\mathrm{},\frac{W_k}{\phi _j},\mathrm{},W_n)\right)\hfill \\ & +i\underset{j}{}\left(_\mu ^y\mathrm{\Delta }_{ij}^F(zy)\right)\omega _0\left(T(W_1,\mathrm{},W_n,\frac{k^\mu }{\phi _j})\right)\hfill \\ & +i\underset{j}{}\mathrm{\Delta }_{ij}^F(zy)_\mu ^y\omega _0\left(T(W_1,\mathrm{},W_n,\frac{k^\mu }{\phi _j})\right)\hfill \\ & =g(\phi _i)\delta (yz)\left(i\underset{k=1}{\overset{n}{}}\underset{j}{}\mathrm{\Delta }_{ij}^F(zx_k)\omega _0\left(T(W_1,\mathrm{},\frac{W_k}{\phi _j},\mathrm{},W_n)\right)\right).\hfill \end{array}$$
(B.14)
Again we can distinguish different cases here.
In the first case, $`\phi _i(u^a)^{(\alpha )},(\stackrel{~}{u}^a)^{(\alpha )}`$, we have again both $`\frac{k^\mu }{\phi _i}=0`$ and $`g(\phi _i)=0`$, so the equation holds automatically. The cases $`\phi _i=(u^a)^{(\alpha )}`$ or $`\phi _i=(\stackrel{~}{u}^a)^{(\alpha )}`$ with $`\left|\alpha \right|2`$ cannot occur because they were explicitely excluded. So there remain four cases where we have to prove that the equation above is indeed valid: $`\phi _i=(u^a),(u^a)^{(1,\mu )},(\stackrel{~}{u}^a)`$ and $`(\stackrel{~}{u}^a)^{(1,\mu )}`$. For simplicity we will treat only $`\phi _i=(u^a)`$, the calculation for the other cases is analogous.
Remembering $`g(u^a)=1`$, $`\frac{k^\mu }{\stackrel{~}{u}_a}=i(u^a)^{(1,\mu )}`$ and $`\frac{k^\mu }{(\stackrel{~}{u}_a)^{(1,\nu )}}=i\delta _\nu ^\mu (u^a)`$ from the first step, we see that the sum of the second and third line on the left hand side of equation (B.14) give, where eqn. (N4) has been used,
$$\begin{array}{cc}& \left(_y^\mu D^F(zy)\right)\left[i\underset{k=1}{\overset{n}{}}\delta (yx_k)\omega _0\left(T(W_1,\mathrm{},\frac{W_k}{(\stackrel{~}{u}^a)^{(1,\mu )}},\mathrm{},W_n)\right)\right]\hfill \\ & +D^F(zy)\left[i\underset{k=1}{\overset{n}{}}\delta (yx_k)\omega _0\left(T(W_1,\mathrm{},\frac{W_k}{(\stackrel{~}{u}^a)},\mathrm{},W_n)\right)\right]\hfill \\ & \delta (zy)[i\underset{k=1}{\overset{n}{}}D^F(yx_k)T(W_1,\mathrm{},\frac{W_k}{(\stackrel{~}{u}^a)},\mathrm{},W_n)\hfill \\ & i\underset{k=1}{\overset{n}{}}\left(_y^\mu D^F(yx_k)\right)T(W_1,\mathrm{},\frac{W_k}{(\stackrel{~}{u}^a)^{(1,\mu )}},\mathrm{},W_n)].\hfill \end{array}$$
(B.15)
Comparing this with the other lines in eqn. (B.14), we see that the last two lines cancel the right hand side of that equation while the first two lines cancel the first line on the the left hand side. So the equation is indeed satisfied. As we already remarked, it can be proven by an analogous calculation that this is also true if $`\phi _i=(u^a)^{(1,\mu )},(\stackrel{~}{u}^a)`$ or $`(\stackrel{~}{u}^a)^{(1,\mu )}`$. With this we have proven that the time ordered products with at least one generator among their arguments satisfy condition (N5) automatically.
Step 3: We know up to now that eqn. (N5) can only be violated by $`T`$-products that have no generator among their arguments, and this violation can be at most a $`\mathrm{C}`$-number. In addition we know that the anomaly must be local because of causal factorization and validity of (N5) in lower orders, so it can be written as
$$a(x_1,\mathrm{},x_n,y)=\omega _0\left(a(x_1,\mathrm{},x_n,y)\right)=P()\delta (yx_1)\mathrm{}\delta (yx_n)$$
(B.16)
for some polynomial of spacetime derivatives $`P()`$. To show that such an anomaly can always be removed we notice that
$$0=d^4yf(y)a(x_1,\mathrm{},x_n,y)=d^4ya(x_1,\mathrm{},x_n,y)$$
(B.17)
where $`f`$ is a test function like in the proof of eqn. (4.32). The first identity is an immediate consequence of that equation. This is the point in the proof of (N5) where it is necessary to know in advance that (B.1) holds. The second identity is true since $`f=1`$ in a domain around each $`x_k`$.
Let us consider the Fourier transformation of the anomaly,
$$\begin{array}{cc}\hfill \widehat{a}(x_1,\mathrm{},x_n,y)& =(2\pi )^nd^4x_1\mathrm{}d^4x_na(x_1,\mathrm{},x_n,y)e^{i(p_1x_1+\mathrm{}+p_nx_n)}\hfill \\ & =(2\pi )^nP(ip_1,\mathrm{},ip_n)e^{i(p_1+\mathrm{}+p_n)y}.\hfill \end{array}$$
(B.18)
For the second identity we have adopted eqn. (B.16) for the anomaly. Inserting (B.18) back into eqn. (B.17), we find that the polynomial $`P(ip_1,\mathrm{},ip_n)`$ vanishes on the hyperplane $`p_1+\mathrm{}+p_n=0`$:
$$P(ip_1,\mathrm{},ip_n)\delta \left(p_1+\mathrm{}+p_n\right)=0.$$
(B.19)
Now we define $`\stackrel{~}{P}(q,p_1,\mathrm{},p_{n1})\stackrel{\text{def}}{=}P(ip_1,\mathrm{},ip_n)`$ with $`q\stackrel{\text{def}}{=}p_1+\mathrm{}+p_n`$ and consider its Taylor expansion around the origin:
$$\stackrel{~}{P}(q,p_1,\mathrm{},p_{n1})=\underset{k=1}{\overset{\mathrm{degree}\stackrel{~}{P}}{}}\underset{\left|\alpha \right|+\left|\beta \right|=k}{}\frac{q^\alpha p^\beta }{\alpha !\beta !}\left(\frac{^{\left|\alpha \right|}^{\left|\beta \right|}}{q^\alpha p^\beta }\stackrel{~}{P}\right)(0)$$
(B.20)
where $`p\stackrel{\text{def}}{=}(p_1,\mathrm{},p_{n1})`$. So the derivatives $`\frac{^{\left|\alpha \right|}}{q^\alpha }`$ describe a variation orthogonal to the hyperplane $`p_1+\mathrm{}+p_n=0`$, the derivatives $`\frac{^{\left|\beta \right|}}{p^\beta }`$ a variation within it. Since $`\stackrel{~}{P}`$ vanishes throughout the entire plane, terms with $`\left|\alpha \right|=0`$ must vanish. Therefore the Taylor expansion can be rewritten as
$$\stackrel{~}{P}(q,p_1,\mathrm{},p_{n1})=q\underset{k=0}{\overset{\mathrm{degree}\stackrel{~}{P}1}{}}\underset{\left|\alpha \right|+\left|\beta \right|=k}{}\frac{q^\alpha p^\beta }{\alpha !\beta !}\left(\frac{^{\left|\alpha \right|}^{\left|\beta \right|}}{q^\alpha p^\beta }\stackrel{~}{P}_1^\mu \right)(0)$$
(B.21)
with a new polynomial $`\stackrel{~}{P}_1^\mu `$. Reversing the Fourier transformation we find
$$P()=\left(\underset{i=1}{\overset{n}{}}_\mu ^i\right)P_1^\mu ()$$
(B.22)
where the polynomial $`P_1^\mu `$ is the Fourier transform of $`\stackrel{~}{P}_1^\mu `$. With this expression we can write the anomaly as
$$a(x_1,\mathrm{},x_n,y)=_\mu ^y\left(nP_1^\mu \delta (x_1y)\mathrm{}\delta (x_ny)\right).$$
(B.23)
So the anomaly can be removed by addition of $`nP_1^\mu \delta (x_1y)\mathrm{}\delta (x_ny)`$ to the the previous normalization of $`T(W_1,\mathrm{},W_n,k^\mu )(x_1,\mathrm{},x_n,y)`$. This is obviously a valid normalization and so the desired normalization has been found.
The question remains why we had excluded polynomials with $`k^\mu `$ as a sub polynomial. The reason is that we can assure that a normalization with the desired properties exists, but we cannot assure that time ordered products like
$$T(k^\mu ,k^\nu ,W_1,\mathrm{},W_n)(y,z,x_1,\mathrm{},x_n)$$
(B.24)
are symmetric under simultaneous exchange of $`\mu ,\nu `$ and $`y,z`$ as they must. There is indeed a counterexample for the Ward identities of the axial current $`j_A^\mu =:\overline{\psi }\gamma ^\mu \gamma ^5\psi :`$, where it is not possible to find a normalization of $`T(j_A^\mu ,j_A^\mu ,j_A^\mu )(x,y,z)`$ with the required symmetries. Excluding the respective polynomials from the allowed arguments makes sure that this situation does not occur.
From the proof above it is clear that the normalization we have found is compatible both with (N3) and (N4). But we can also immediately see that (N5) respects Poincaré transformation properties and therefore (N1). Taking the adjoint of (N5) finally reveals that it complies also with (N2) and therfore eventually with all other normalization conditions.
### B.2. Relation between (N6) and generalized operator gauge invariance
It is always possible to define an operator valued distribution $`T(_{i_1},\mathrm{},_{i_n},j^\mu )`$ by
$$\begin{array}{cc}& T(_{i_1},\mathrm{},_{i_n},j^\mu )(x_1,\mathrm{},x_n,y)\stackrel{\text{def}}{=}\hfill \\ & \stackrel{\text{def}}{=}s_0T(_{i_1},\mathrm{},_{i_n},k^\mu )(x_1,\mathrm{},x_n,y)\hfill \\ & +i\underset{m=1}{\overset{n}{}}_\nu ^mT(_{i_1},\mathrm{},_{i_m+1}^\nu ,\mathrm{},_{i_n},k^\mu )(x_1,\mathrm{},x_n,y)\hfill \\ & i\underset{m=1}{\overset{n}{}}\delta (x_my)T(_{i_1},\mathrm{},_{i_m+1}^\mu ,\mathrm{},_{i_n})(x_1,\mathrm{},x_n)\hfill \\ & +i\underset{m=1}{\overset{n}{}}\delta (x_my)i_mT(_{i_1},\mathrm{},_{i_m+1}^\mu ,\mathrm{},_{i_n})(x_1,\mathrm{},x_n).\hfill \end{array}$$
(B.25)
$`T(_{i_1},\mathrm{},_{i_n},j^\mu )`$ is at this point only a name for that distribution, we must still prove that it is indeed an extension of $`{}_{}{}^{0}T(_{i_1},\mathrm{},_{i_n},j^\mu )`$. Before we do that, we point out that it implicates (N6) almost immediately. Of course the time ordered products on the right hand side must satisfy eqn. (N5). Taking the derivative w.r.t. the $`y`$ coordinate, we find with (N5)
$$\begin{array}{cc}& _\mu ^yT(_{i_1},\mathrm{},_{i_n},j^\mu )(x_1,\mathrm{},x_n,y)=\hfill \\ & =\left(\underset{m=1}{\overset{n}{}}\delta (x_my)i_m\right)s_0T(_{i_1},\mathrm{},_{i_n})(x_1,\mathrm{},x_n)\hfill \\ & \begin{array}{cc}\hfill +i\underset{l=1}{\overset{n}{}}_\nu ^l[T(_{i_1},\mathrm{},_{i_l+1}^\nu ,\mathrm{},_{i_n})(x_1,\mathrm{},x_n)\times & \\ \hfill \times (\underset{m=1}{\overset{n}{}}\delta (x_my)i_m+\delta (x_ly))]& \end{array}\hfill \\ & +i\underset{m=1}{\overset{n}{}}\left(_\nu ^m\delta (x_my)\right)T(_{i_1},\mathrm{},_{i_m+1}^\nu ,\mathrm{},_{i_n})(x_1,\mathrm{},x_n)\hfill \\ & i\underset{m=1}{\overset{n}{}}\left(_\nu ^m\delta (x_my)i_m\right)T(_{i_1},\mathrm{},_{i_m+1}^\nu ,\mathrm{},_{i_n})(x_1,\mathrm{},x_n).\hfill \end{array}$$
(B.26)
Smearing out this equation with a test function $`f`$ like the one defined following eqn. (B.2) gives eqn. (B.1), the calculation is the same as at the beginning of the last section. Inserting this result into eqn. (B.26) we get immediately eqn. (N6).
Eqn. (B.25) is obviously a well posed definition since all operations involved in it are well defined — in-particular the time ordered product in the last line contains no vertex at $`y`$, so the product with the delta distribution is a tensor product.
The crucial question is whether the operator valued distribution has the correct causal factorization outside the diagonal $`\mathrm{Diag}_{n+1}`$. Only then it is really a time ordered product of its arguments as the notation suggests. Basically we must do the same construction as in the respective point for (N4), see section (A.2). We give here only a simplified version of this proof where the essential point may be more easily understood. The detailed version can easily be derived from this sketch.
Suppose the points $`x_1,\mathrm{},x_n`$ are in a relative position such that
$$\mathrm{}I=\{x_1,\mathrm{},x_k\}\{x_{k+1},\mathrm{},x_n,y\}.$$
(B.27)
This is the situation we encounter in eqn. (A.17) in the first sum — if $`I=\{x_1,\mathrm{},x_k\}\{x_{k+1},\mathrm{},x_n,y\}`$, corresponding to the second sum there, the argument works as well. Then
$$T(_{i_1},\mathrm{},_{i_n},j^\mu )=T(_{i_1},\mathrm{},_{i_k})T(_{i_{k+1}},\mathrm{},_{i_n},j^\mu )$$
(B.28)
where we omitted the spacetime indices for simplicity. Eqn. (B.25) is valid for $`T(_{i_{k+1}},\mathrm{},_{i_n},j^\mu )`$ since we assumed that eqn. (B.25) holds already for time ordered products with fewer arguments. Together with eqn. (B.28) this gives the following expression
$$\begin{array}{cc}& T(_{i_1},\mathrm{},_{i_n},j^\mu )=\hfill \\ & =s_0\left[T(_{i_1},\mathrm{},_{i_k})T(_{i_{k+1}},\mathrm{},_{i_n},k^\mu )\right]\hfill \\ & +\left[s_0T(_{i_1},\mathrm{},_{i_k})\right]T(_{i_{k+1}},\mathrm{},_{i_n},k^\mu )\hfill \\ & +\underset{l=k+1}{\overset{n}{}}_\nu ^l\left[T(_{i_1},\mathrm{},_{i_k})T(_{i_{k+1}},\mathrm{},_{i_l+1}^\nu ,\mathrm{},_{i_n},k^\mu )\right]\hfill \\ & i\underset{l=k+1}{\overset{n}{}}\delta (x_ly)\left[T(_{i_1},\mathrm{},_{i_k})T(_{i_{k+1}},\mathrm{},_{i_l+1}^\mu ,\mathrm{},_{i_n})\right]\hfill \end{array}$$
(B.29)
where we have omitted spacetime arguments for simplicity. For $`s_0T(_{i_1},\mathrm{},_{i_k})`$ we may use the generalized operator gauge invariance (4.35) in lower orders as long as $`kn`$. Unfortunately also the case $`k=n`$ occurs if all the $`x_i`$ coincide and only $`y`$ is separated from them. This is the only case where we must know in advance that (4.35) holds. If this would not be true then our definition (B.25) would be a well defined operator valued distribution, but no an extension of $`T^0(_{i_1},\mathrm{},_{i_n},j^\mu )`$ to the diagonal — this means that it could differ from the $`T^0`$-product even outside the diagonal. Hence we need to assume that (4.35) is valid also for $`k=n`$. Furthermore we may add in the last sum the terms with $`l=1,\mathrm{},k`$ since the delta distributions vanish because $`y`$ and the $`x_1,\mathrm{},x_k`$ may never coincide. Recombining the products of $`T`$-products into a single $`T`$-product according to eqn. (B.28) one gets immediately (B.25). As already remarked the calculation comes to the same result if $`\mathrm{}I=\{x_1,\mathrm{},x_k\}\{x_{k+1},\mathrm{},x_n,y\}`$. So (B.25) is a well defined operator valued distribution that agrees — as long as (4.35) is valid — with $`T^0(_{i_1},\mathrm{},_{i_n},j^\mu )`$ if smeared with a test function that vanishes with all its derivatives on the diagonal, so it is an extension of that $`T^0`$-product to the diagonal and therefore a possible normalization of $`T(_{i_1},\mathrm{},_{i_n},j^\mu )`$.
So we have just proven that the conditions (N6) and (4.35) are equivalent. |
warning/0001/astro-ph0001205.html | ar5iv | text | # Regular magnetic fields in the dwarf irregular galaxy NGC 4449
## 1 Introduction
The generation of large-scale galactic magnetic fields from small-scale field perturbations caused by turbulence (as postulated by the dynamo concept) requires a preferred sense of twisting of turbulent gas motions, called the $`\alpha `$-effect (Wielebinski & Krause wielebinski (1993)). In normal spiral galaxies it is determined by Coriolis forces caused by the disk rotation giving rise to strong dynamo action and to the observed spiral-like regular magnetic fields (Beck et al. 1996b ). To make the dynamo process work, either the differential rotational shear or the galaxy’s angular speed (in case of rigid rotation) must exceed certain threshold values (Ruzmaikin et al. ruzmaikin (1988)).
Dwarf irregulars are small, low-mass galaxies with a patchy distribution of star-forming regions. Though they exhibit a large variety of rotation curves (Hunter et al. 1998a ) many of them show slow rotation with much less rotational shear than in normal spirals. Some dwarf irregulars show complex velocity fields with chaotic motions comparable in speed to the overall rotation. Even if the dynamo could still work in such conditions, the generation time scales of the magnetic fields estimated from classical dynamo theory would be very long and strong large-scale magnetic fields are not expected. Their observational detection would mean that the dynamical role of global magnetic fields in gas dynamics and star formation in irregular galaxies has to be reconsidered.
Signatures of a global magnetic field were already detected in the Large Magellanic Cloud (LMC, Klein et al. klein93 (1993)). However, this galaxy still shows a significant degree of differential rotation (Luks & Rohlfs luks (1992)) so that, like in normal spirals, the standard dynamo process could be at work. In this paper we present a sensitive radio polarization study of the dwarf irregular galaxy NGC 4449 which exhibits only weak signs of global rotation (cf. also Sabbadin et al. sabbadin84 (1984), Hartmann et al. hartmann (1986)). The radial velocities in NGC 4449 relative to the systemic one reach $`\pm 20`$$`30`$ km/s. However, the analysis of the high-resolution HI data cube (kindly supplied by Dr D. Hunter) does not show the classical picture of a global rotation. Instead, NGC 4449 shows velocity jumps and gradients along both the major and minor axis with centroids not coincident with the optical centre. They are intermixed with chaotic velocity variations with an amplitude of about 10 – 15 km/s. These very complex and chaotic kinematics, partly due to the interactions with DDO125 (Hunter at al. 1998b ) and possibly also to a high star formation rate, make NGC 4449 an interesting target to investigate the magnetic field structure under conditions very difficult for the classical galactic dynamo. The basic parameters of NGC 4449 are summarized in Table 1.
A low-resolution detection of polarized emission (Klein at al. klein96 (1996)) showed that the magnetic field in NGC 4449 is running across its bright star-forming body, very different from that in normal galaxies. In this work we present a total power and polarization study of this galaxy with a resolution and sensitivity several times better than that in the work of Klein et al. (klein96 (1996)). The use of two frequencies (8.46 and 4.86 GHz) enables us to determine the distribution of Faraday rotation over the disk of NGC 4449, allowing to discriminate between the galaxy-scale uniform fields and those passively stretched and compressed in the gas flows powered by huge star-forming regions.
## 2 Observations and data reduction
The maps of total power and linearly polarized radio emission of NGC 4449 at 8.46 GHz and 4.86 GHz were obtained using the Very Large Array (VLA) of the National Radio Astronomy Observatory (NRAO) <sup>1</sup><sup>1</sup>1The National Radio Astronomy Observatory is a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc.. To attain highest sensitivity to smooth extended structures the most compact (D) configuration was used. The observations were carried out on 30 August and 1 – 2 September 1996 using 27 antennas at two independent IFs, each with a bandwidth of 50 MHz, separated by 50 MHz, with 11<sup>h</sup> integration time on NGC 4449 at 6.2 cm and 16<sup>h</sup> at 3.5 cm. The data were reduced using the standard AIPS software package. The flux density scale and the position angle of polarization was calibrated by observing the point source 3C 286. Instrumental polarization was corrected by observing 1216+487, which was also used for gain and phase calibration. The calibrated and edited visibility data were cleaned and self-calibrated (in phase only) using the AIPS package, yielding maps of Stokes parameters I, Q and U.
These IUQ data were combined with Effelsberg measurements using the program EFFMERG, a version of the SDE task IMERG (Cornwell et al. cornwell (1995)) modified by P. Hoernes (see Beck & Hoernes 1996a ). This program deconvolves both clean maps with their beams, and Fourier transforms them back into the UV plane. Then a combination of Effelsberg data for small spacings and VLA data for large spacings is performed with a linear interpolation in the overlap domain. The combined data are transformed back into the image plane with the synthesized VLA beam. To avoid ring-like distortions around strong unresolved sources, introduced by the combination technique, the bright unresolved sources were first subtracted from the maps and added again after the combination (see Beck et al. beck97 (1997)). For our 8.46 GHz map we used the single dish data at 10.55 GHz from Klein et al. (klein96 (1996)), scaled to our frequency. We used mean spectral index of $`0.7`$ (S$`{}_{\nu }{}^{}\nu ^\alpha `$). For 4.86 GHz we performed separate observations with the 100-m Effelsberg radio telescope at 4.85 GHz.
The Q and U maps were combined to get maps of the linearly polarized emission (corrected for the positive zero level offset) and of the position angle of polarization E-vectors. The final maps have the synthesized beam of $`12\mathrm{}`$ at 8.46 GHz and $`19\mathrm{}`$ at 4.86 GHz. To obtain the map of Faraday rotation the data at both frequencies were convolved to a common beam of $`19\mathrm{}`$.
## 3 Results
### 3.1 Total power and polarized emission at 8.46 GHz
The total power map of NGC 4449 at 8.46 GHz with apparent B-vectors of polarized intensity is shown in Fig. 1. As no correction for Faraday rotation was applied, the orientation of observed B-vectors may differ from the magnetic field directions by some $`3\mathrm{°}`$$`6\mathrm{°}`$ on average in the disk, the maximum difference reaching $`10\mathrm{°}`$ in small regions of high Faraday rotation measures (see Sect. 3.3). The map shows details of the radio structure in the inner disk. The total power emission shows strong peaks at the position of bright star-forming regions. In addition to that diffuse radio emission away from the optically bright star-forming body has been detected as well. This radio envelope extends along the galaxy’s minor axis up to $`2\mathrm{}`$ (corresponding to $`2.2`$ kpc) from the main plane. The extent of the radio envelope at 8.46 GHz is larger than that of the faint diffuse H$`\alpha `$ emission (see Fig. 1). In the southern disk the radio emission has an extension towards a nebulous object at RA$`{}_{2000}{}^{}=12^h28^m06\stackrel{s}{.}7`$, Dec$`{}_{2000}{}^{}=44\mathrm{°}03\mathrm{}39\mathrm{}`$ (probably a supernova remnant), forming a faint peak at its position.
The contour map of polarized brightness with apparent B-vectors proportional to the polarization degree is shown in Fig. 2. The extended radio emission is substantially polarized (locally up to 50%), with extended ($`1`$ kpc) domains of highly aligned B-vectors. The magnetic field structure in the inner disk looks unusual at first glance (Figs. 1 and 2). The projected magnetic vectors in NGC 4449 show two distinct kinds of structure. From the bright central star-forming region they are directed radially outwards, on each side forming a polarized “fan”. The B-vectors are parallel to the H$`\alpha `$ filaments discussed in detail by Sabbadin & Bianchini (sabbadin79 (1979)) and by Bomans et al. (bomans (1997)). In the galaxy’s outskirts, the magnetic vectors run along a polarized ridge encircling the galaxy on the northern, north-eastern and eastern side. Between this structure and the eastern “fan” an elongated unpolarized “valley” is due to a geometrical superposition of mutually perpendicular polarization directions in the “fan” and in the polarized ridge.
### 3.2 Total power and polarization at 4.86 GHz
Our maps at 4.86 GHz have a considerably worse resolution than those at 8.46 GHz, and the orientations of the B-vectors may be subject to stronger Faraday rotation (on average about $`10\mathrm{°}`$ but locally up to $`30^o`$, see Fig. 5). However, due to a higher signal-to-noise ratio at 4.86 GHz the radio emission is traced much further out (Fig. 3). At this frequency we can trace the radio envelope in the sky plane out to $`3\stackrel{}{.}2`$ (3.5 kpc) from the galaxy’s major axis. The nonthermal emission thus extends into the halo beyond one isophotal (at the level of 25$`{}_{}{}^{m}/(\mathit{}\mathrm{})`$) major axis radius, which is a rare (though not exceptional) phenomenon among spiral galaxies (e.g. Hummel et al. hummel (1991)).
The map at 4.86 GHz again shows the polarized “fans”, however the eastern one is less conspicuous at this frequency than at 8.46 GHz, which suggests stronger Faraday depolarization in this region. The polarized ridge in the northwestern portion of the galaxy, already visible in Fig. 2, turns out to be part of a larger polarized ring surrounding the galaxy from the northeast through north, east and south down to the southwest, with a well-organized, coherent pattern of magnetic vectors (Fig. 4). Another weak fragment of the polarized ring is visible west of the centre. The ring coincides well with a similar feature visible in HI (Hunter, priv. comm., Fig. 4), with one of the brightest polarization peaks lying close to the densest neutral gas clump. Along the ring the polarization B-vectors are not exactly tangential to the azimuthal directions or to the HI shell. They deviate systematically from the azimuthal directions by some $`20\mathrm{°}`$$`40\mathrm{°}`$. A detailed discussion of the magnetic field directions is presented in Sect. 4.
Due to the higher sensitivity at 4.86 GHz our map shows very well the unpolarized “valley” not only at the interface of the eastern “fan” and the ridge but also a similar feature in the NW disk. In both cases they result from a geometrical superposition of magnetic field directions in the “fans” and in the polarized ring, seen almost perpendicular to each other when projected to the sky plane.
### 3.3 Faraday rotation
The distribution of Faraday rotation measures between 8.46 and 4.86 GHz is shown in Fig. 5. The northern and eastern parts of the polarized ring, as well as the “magnetic fan” east of the central star-forming complex show coherently positive Faraday rotation measures (RM) over areas with sizes of about $`1.5\mathrm{}`$, with a mean value of about +50 rad/m<sup>2</sup>. The values of RM are rising locally up to +200 rad/m<sup>2</sup>. The western “fan” and the southern part of the polarized ring are dominated by negative RMs, on average of about $`50`$ rad/m<sup>2</sup> but also reaching $`150`$ rad/m<sup>2</sup> locally. The errors in these regions vary from $`\pm 10`$ to $`\pm 20`$ rad/m<sup>2</sup> in regions of low RMs, exceeding $`\pm 50`$ rad/m<sup>2</sup> in regions of high rotation measures. However, though in individual points the values of RM do not generally exceed the errors by more than 2 – 2.5$`\sigma `$ r.m.s. errors, coherent areas of the same sign of RM extend over many beam sizes. The statistical significance of our determinations of RM is discussed in detail in Sect. 4.
## 4 Discussion
### 4.1 Distribution of thermal emission
In order to determine the equipartition magnetic field strength in selected regions of NGC 4449 we need to estimate the distribution of thermal emission in the galaxy. The spectral index computed between the maps of NGC 4449 at 8.46 GHz and 4.86 GHz changes from about $`0.3`$ (S$`{}_{\nu }{}^{}\nu ^\alpha `$) in strongly star-forming regions to $`1.1`$ locally in the outer southern region. Somewhat smaller variations are found by Klein et al. (klein96 (1996)), probably because of the much lower resolution used by these authors. Although Klein et al. (klein96 (1996)) found variations of the nonthermal spectral index $`\alpha _{nt}`$ between $`0.5`$ in the central star-forming region to $`0.8`$ in the eastern part of the halo, for our purpose it was sufficient to assume $`\alpha _{nt}`$ constant over the whole galaxy. Possible uncertainties due to this assumption were included in the errors. To determine the nonthermal spectral index we compared the radial distribution of the thermal brightness S<sub>th</sub> at 8.46 GHz and that in the H$`\alpha `$ line (S), convolved to the beam of $`19\mathrm{}`$. We found that they are identical for $`\alpha _{nt}`$ = $`0.9`$. This value differs only by about 1.5$`\sigma `$ r.m.s. from the value obtained by Klein et al. (klein96 (1996)) from the radio spectrum, but agrees better with their estimate based on thermal flux obtained from the H$`\alpha `$ emission.
The distribution of thermal fraction f<sub>th</sub> at 8.46 GHz in NGC 4449 (Fig. 6) shows clear peaks at the positions of bright star-forming complexes, f<sub>th</sub> reaches 80% there. After subtraction of the thermal emission these regions are still considerably brighter than the diffuse emission from the surroundings by some 40%. Away from bright star-forming complexes the emission is largely nonthermal, the free-free emission amounts to not more than 10%.
As an additional test of our assumption of $`\alpha _{nt}`$ constant over the galaxy’s body we analyzed the point-to-point correlation between maps of the radio thermal flux at 8.46 GHz and that in the H$`\alpha `$ line convolved to $`19\mathrm{}`$. We checked that, using the Monte-Carlo simulations of two-dimensional arrays of points convolved to various beams, values in map points separated by 1.2 times the beam size are correlated only by some 10 – 12% and are for our purposes almost independent. Therefore we used points separated by $`22\stackrel{s}{.}8`$. In order to eliminate an artificial correlation caused by the radial decrease of all quantities, each map was divided by an axisymmetric model obtained by integrating the map in elliptical rings with the position angle and inclination taken from Tab. 1. A correlation slope significantly larger than 1 would mean that we have overestimated the thermal radio emission in strongly star-forming regions while in fact they have much flatter nonthermal spectra. Using the orthogonal regression we obtained S$`{}_{th}{}^{}`$ S$`{}_{}{}^{0.96\pm 0.09}{}_{H\alpha }{}^{}`$, thus close to a linear relationship, though few regions deviate strongly from the best-fit line. This means that $`\alpha _{nt}`$ shows some place-to place variations, but our assumption of $`\alpha _{nt}`$ = constant does not introduce large, systematic errors in determining the thermal fraction. A detailed multi-dimensional analysis of both radio emission components involving the H$`\alpha `$, CO, HI and X-ray data will be the subject of a separate study.
### 4.2 Magnetic field strengths
To determine the magnetic field strengths in NGC 4449 from the synchrotron emission we assumed the equipartition conditions between magnetic fields and cosmic rays to be valid everywhere in the galaxy. Furthermore we adopted a proton-to-electron ratio of energy densities of 100 and a lower energy cutoff of cosmic ray electrons of 300 MeV. We assumed a face-on thickness of the nonthermal disk of 2 kpc, resulting from a typical scale height of galactic radio disks of 1 kpc (Hummel et al. hummel (1991)), determined by the propagation range of cosmic ray electrons. With the inclination of NGC 4449 from Tab. 1 this implies a mean pathlength through the galaxy of 2.8 kpc. The errors of estimated magnetic field strengths include an uncertainty of these quantities of a factor two. The thermal fractions were taken from results described in Sect. 4.1.
Under these assumptions we determined the mean magnetic field strength for the whole galaxy and for selected regions; the results are summarized in Tab. 2. Regular magnetic fields derived from the polarized intensity were found to reach locally up to $`7\pm 2\mu `$G in the western magnetic “fan” and about $`8\pm 3\mu `$G in the radio-bright part of the polarized ring. The total magnetic field in these regions, determined from the total power emission reaches $`14\pm 4\mu `$G, comparable to that in the radio-brightest spiral galaxies (Beck et al. 1996b ). A slow rotation of NGC 4449 accompanied by chaotic gas motions apparently does not exclude the existence of strong, regular magnetic fields.
### 4.3 Magnetic field structure
#### 4.3.1 Magnetic field coherence
The presence of polarized emission alone does not provide a definite proof for dynamo-generated, spatially coherent magnetic fields. Substantial polarization may be also produced by random fields, made anisotropic by squeezing or stretching, e.g. by stellar winds or large-scale shocks from multiple supernova events, however, frequent field reversals along the line of sight would completely cancel the Faraday rotation. Non-zero rotation measures imply the magnetic fields in the observed galaxy coherent over scales much larger than the telescope beam.
Although the values of RM in individual points of our Faraday rotation map (Fig. 5) do not exceed the errors by much, we note that they deviate coherently from zero, forming large domains of constant RM sign (both positive and negative). These regions with mean RM of $`\pm 50`$ rad/m<sup>2</sup> are up to 20 times larger than the telescope beam area. A correction for the foreground rotation of $`35`$ rad/m<sup>2</sup> was estimated from background sources present in our map and checked with the galactic RM map by Simard-Normandin & Kronberg (simard (1980)). At the galactic latitude of NGC 4449 of $`72\mathrm{°}`$ the existence of foreground rotation structures changing sign over angular scales of $`2\mathrm{}`$$`4\mathrm{}`$ with an amplitude of 100 rad/m<sup>2</sup>, correlated with particular features in the galaxy’s polarized intensity, is unlikely. Thus, the observed Faraday effects almost certainly originate inside NGC 4449.
To check quantitatively the coherence of the non-zero Faraday rotation we computed values of RM and its error $`\sigma _{RM}`$ in a grid of points separated by $`22\stackrel{s}{.}8`$ (1.2 times the beam size). In case of a lack of systematic Faraday rotation such points would show only little correlation (see Sect. 4.1) and the variable defined as RM/$`\sigma _{RM}`$ is expected to fluctuate randomly from point to point with a zero mean and unity variance. However, we found that its mean value deviates from zero in the eastern polarized ridge by more than $`4.1\sigma _{mean}`$ as well as in the western “fan” by more than $`4.3\sigma _{mean}`$, $`\sigma _{mean}`$ being the r.m.s. error of mean RM in a given region. This implies that the probability of creating at random such large non-zero RM domains is less than $`10^5`$. In the eastern, weaker “fan”, the deviation amounts to only $`1.6\sigma _{mean}`$ (the probability of a random occurrence of non-zero RM of 10%), because of a worse signal-to-noise ratio. The results were found to be independent of the assumed foreground rotation. Thus we conclude that NGC 4449 contains genuine unidirectional fields, rather than stretched and compressed random magnetic field. The latter one would have different sky-projected components yielding a substantial polarization while the line-of-sight component would frequently change sign which would cancel any systematic Faraday rotation. We note that the growth of galaxy-scale coherent, unidirectional fields lies at the foundations of the dynamo process.
#### 4.3.2 Magnetic field geometry
Fig. 7 a and b presents the distribution of magnetic field orientations in the azimuth-ln(R) frame (R being the radial distance form galaxy’s optical centre), in which the logarithmic spiral appears as a set of straight lines inclined by the spiral’s pitch angle. At 8.46 GHz (little Faraday rotation) we clearly see a combination of the radial field in the inner region out to ln(R) of 0.5 – 0.6 and a more azimuthal one at larger radii. However, at this frequency the picture in the outer galaxy regions becomes rather noisy. A comparison of Figs. 7 a and b shows that Faraday rotation does not much change the global field picture in the inner region where the ionized gas density is highest and Faraday effects strongest. Thus, the 4.86 GHz data alone can be safely used in the galaxy outskirts.
Fig. 7b shows a very well ordered field in the polarized ring with the magnetic pitch angle $`\psi `$ keeping a constant sign over most of azimuthal angles (except a low signal-to-noise region at azimuths of $`0\mathrm{°}`$$`60\mathrm{°}`$ and ln(R) $`1`$). The value of $`\psi `$ is $`40\mathrm{°}`$ on average, with local variations. It resembles a somewhat distorted magnetic spiral with a substantial radial component. This, like in rapidly rotating spiral galaxies, may signify dynamo-type fields (Urbanik et al. urbanik (1997)), while the random field pushed away from the galaxy and squeezed by an expanding gaseous shell would yield the observed B-vectors parallel to the shell. Nevertheless, the pitch angles show some place-to-place changes, possibly due to processes like local outflows or compressions. The strongest distortion of the spiral - the region of nearly pure toroidal magnetic field at azimuthal angles $`270\mathrm{°}`$, ln(R) $`0.5`$ coincides with the densest HI clump and a region of star formation. The analysis of recent CO data (Kohle et al. in preparation) suggests strong gas compression possibly due to external interactions. We note also an opposite sign of Faraday rotation at both ends of the major axis, which is typical for axisymmetric magnetic fields.
The radial magnetic “fans” are structural elements not observed in spiral galaxies. They may be due to magnetic fields pulled out from the central star-forming region by gas outflows. Evidence for radial gas flows in NGC 4449 was indeed found by Martin (martin98 (1998), martin99 (1999)). However, in case of an initially random magnetic field (e.g. injected by supernovae) being stretched by gas flows, the “fans” would contain interspersed magnetic lines directed towards and outwards from the star-forming complex, yielding no significant Faraday rotation (see Sect. 4.3.1). Thus if the radial “fans” would result from the gaseous wind, a large-scale, coherent preexisting magnetic field would still be required, like one resulting from the dynamo process.
Alternatively, the observed magnetic field structure in NGC 4449 can be qualitatively explained by classical dynamo-generated fields. In addition to a toroidal field running around the disk, the classical dynamo process also generates a poloidal field with lines of force forming closed loop-like structures perpendicular to the disk plane and with diameters comparable to the galaxy radius (Donner & Brandenburg donner (1990)). They are due to a radial field component, B<sub>r</sub>, turning into a vertical one, B<sub>z</sub>, close to the centre and in the disk outskirts. The conservation of magnetic flux leads to B<sub>z</sub> being always much stronger in the central region than in the outer disk. In large spiral galaxies the vertical segments of the poloidal field loops with the strongest B<sub>z</sub> probably lie at heights $`2`$ – 3 kpc. This is too high to see the vertical magnetic field in synchrotron emission, as the latter has a vertical scale height of about 1 kpc (Hummel et al. hummel (1991)) due to a limited propagation range of cosmic ray electrons. As an exception NGC 4631 has a much larger scale height and dominating vertical fields in its inner regions (Hummel et al. hummel (1991)).
With its bright star-forming disk of about 4 kpc diameter NGC 4449 is several times smaller than normal spirals. If it had a classical poloidal dynamo-type magnetic field, its magnetic lines would make closed loop-like structures with a vertical size of about 1 – 1.5 kpc. The maximum B<sub>z</sub> would occur at some hundreds of parsecs above the galaxy’s plane, well within the propagation range of radio-emitting electrons, making vertical fields visible in emission. The intense star formation in NGC 4449 and its low gravitational potential may give rise to galactic winds which may additionally enhance the generation of vertical magnetic fields (Brandenburg et al. branden93 (1993)). With the inclination of NGC 4449 (Table 1) a strong poloidal field in the central part of NGC 4449, projected to the sky plane, may give rise to the observed radial magnetic “fans”. A detailed MHD model of magnetic field evolution in NGC 4449 is a subject of a separate study (Otmianowska-Mazur et al., in prep.). We note also that superimposed on the global magnetic field, smaller-scale ($`<`$ 0.5 kpc or $`30\mathrm{}`$ in our map) local phenomena (e.g. magnetized shells or giant magnetic loops caused by Parker instabilities, Klein et al. klein96 (1996)) may be present, as well. They may explain e.g. local RM reversals, like that seen in the eastern “fan” at RA$`{}_{2000}{}^{}12^h28^m12^s`$, Dec$`{}_{2000}{}^{}44\mathrm{°}04\mathrm{}30\mathrm{}`$.
Although the dynamo process constitutes some possibility to explain the magnetic field structure in NGC 4449, the question arises how the dynamo mechanism can work in this galaxy. Despite some evidence for the dynamo action strong regular magnetic fields are hard to explain by classical dynamo models which, given the weak signs of rotation of NGC 4449, yield growth rates of the regular magnetic field at least an order of magnitude smaller than in rapidly rotating spirals (see e.g. Brandenburg & Urpin branden98 (1998)). Estimates kindly provided by Dr Anvar Shukurov indicate that for the rotation speed and dimensions of NGC 4449 the classical, Coriolis force-driven $`\alpha `$-effect is too weak for the onset of either $`\alpha \omega `$ or $`\alpha ^2`$ dynamo (see Ruzmaikin et al. ruzmaikin (1988) for definitions). Faster field amplification is predicted by a recent concept of the dynamo driven by magnetic buoyancy and sheared Parker instabilities (e.g. Moss et al. moss (1999)). Crude estimates of its efficiency by A. Shukurov (priv. comm.) show that the $`\alpha ^2`$ dynamo process is easily excited throughout most of the galaxy’s body. However, what kind of structure is generated in such conditions remains still an open question and will be a subject of separate analytical and numerical studies.
Among other possibilities we can mention e.g. fast dynamos (Parker parker (1992)), interrelations between small-scale velocity and magnetic field perturbations caused by specific instabilities (Brandenburg & Urpin branden98 (1998)) or even magnetic field amplification without any $`\alpha `$-effect at all (Blackman blackman (1998)). As in these concepts ordered rotation is still needed it is not known how they would work in the complex velocity field of NGC 4449. In summary, our work provides arguments in support of non-standard magnetic field generation mechanisms, though some elements of its structure may be due to gas outflow processes. Still a lot of theoretical work is needed to understand how a classical mixture of poloidal and toroidal fields, similar to that in rapidly rotating spirals can arise in a slowly and chaotically rotating object. Nevertheless, it seems that the existence of strong, dynamically important magnetic fields in dwarf irregulars cannot be ignored.
## 5 Summary and conclusions
We performed a total power and polarization study of the dwarf irregular galaxy NGC 4449 at 8.46 and 4.86 GHz using VLA in its D-configuration. The object rotates slowly and chaotically, thus no large-scale regular magnetic fields were expected. To reach the maximum sensitivity to extended structures we combined our VLA data with the Effelsberg ones at 10.55 GHz and 4.85 GHz, respectively. Despite the slow and chaotic rotation of NGC 4449, unfavourable for dynamo-induced magnetic fields, we found it to possess strong regular, galaxy-scale fields.
The following results were obtained:
* NGC 4449 shows a large, partly polarized halo extending from its main plane up to 3.5 kpc, more than the isophotal major axis radius at 25$`{}_{}{}^{m}/(\mathit{}\mathrm{})`$.
* The radio-brightest peaks coincide with strongly star-forming regions. These regions show increased thermal fractions (up to 80%), however the nonthermal emission is enhanced there as well.
* The galaxy possesses regular magnetic fields reaching locally 6 – 8 $`\mu `$G, comparable to those in rapidly rotating spiral galaxies.
* NGC 4449 shows large domains of non-zero Faraday rotation measures indicating a genuine galaxy-scale regular magnetic field rather than random anisotropic ones with frequent reversals of their direction.
* The magnetic field structure consists of two basic elements: radial “fans” stretching away from the central star-forming complex and a magnetic ring at the radius of about 2.2 kpc. The magnetic field in the ring shows clear characteristics of a magnetic spiral with a substantial radial component signifying dynamo action.
* Both the radial “fans” and the polarized ring can still be explained in terms of a combination of sky-projected poloidal and toroidal dynamo-generated fields, taking into account the smaller size of NGC 4449 compared to normal massive spirals. Alternatively, magnetic “fans” could result from the gas outflow from the central star-forming complex. Even in this case a large-scale coherence of the magnetic field subject to stretching by outflows is required. Thus, some kind of dynamo action is needed, with a preference of non-standard (e.g. buoyancy-driven) dynamos. Whether and how this process can produce classical dynamo-like magnetic fields in a complex and chaotic velocity field of NGC 4449 remains yet unknown.
The detection of regular magnetic fields in spiral galaxies is important for understanding processes like turbulence, turbulent diffusion and the magnetic field generation in astrophysical plasmas; this is also of importance for plasma physics in general. It demonstrated that even in a highly turbulent medium large-scale regular fields can persist and grow quite efficiently. This has already boosted the development of dynamo theories applicable not only to a variety of astrophysical objects from planets to clusters of galaxies but also to laboratory plasmas. On the other hand, there was a widespread prejudice that all the mentioned concepts are restricted solely to rapidly rotating plasma bodies. Against these expectations we show that strong regular fields can also arise in slowly and chaotically rotating systems. Their dynamical role in dwarf irregulars, especially in processes of star-formation triggered by magnetic instabilities, filament formation and confinement or even accelerating galactic winds via cosmic ray pressure exerted on MHD waves (Breitschwerdt et al. breit (1991)), cannot be further neglected. NGC 4449 is the irregular galaxy with the best studied magnetic field so far. We believe that further progress needs more detailed models for such objects. Further detailed observations of the radio polarization of a larger number of irregulars with various morphological characteristics are also required.
###### Acknowledgements.
The Authors wish to express their thanks to Dr Dominik Bomans from Astronomisches Institut der Ruhr-Universität Bochum for providing us with his H$`\alpha `$ map in a numerical format. We are grateful to numerous colleagues from the Max-Planck-Institut für Radioastronomie (MPIfR) in Bonn for their valuable discussions during this work. We want to express our profound gratitude to Dr Elly M. Berkhuijsen from MPIfR for her critical reading of the manuscript and precious suggestions concerning its improvement. M.U. and K.Ch. are indebted to Professor R. Wielebinski (MPIfR) for the invitations to stay at this institute where substantial parts of this work were done. One of us (K.Ch.) is indebted to Professor Miller Goss from NRAO for his invitation to Soccorro and his assistance in some parts of data reduction. We are also grateful to colleagues from the Astronomical Observatory of the Jagiellonian University in Kraków for their comments. This work was supported by a grant from the Polish Research Committee (KBN), grant no. 962/P03/97/12. Large parts of computations were made using the HP715 workstation at the Astronomical Observatory in Kraków, partly sponsored by the ESO C&EE grant A-01-116 and on the Convex-SPP machine at the Academic Computer Centre ”Cyfronet” in Kraków (grant no. KBN/C3840/UJ/011/1996 and KBN/SPP/UJ/011/1996). |
warning/0001/math0001128.html | ar5iv | text | # Local tree-width, excluded minors, and approximation algorithms
## 1 Introduction
*Tree-width*, measuring the similarity of a graph with a tree, has turned out a to be an important notion both in structural graph theory and in the theory of graph algorithms. It is well known that planar graphs may have arbitrarily large tree-width. However, for every fixed $`d`$ the class of planar graphs of diameter at most $`d`$ has bounded tree-width. In other words, the tree-width of a planar graph can be bounded by a function of the diameter of the graph. This makes it possible to decompose planar graps into families of graphs of small tree-width in an orderly way. Such decompositions of planar graphs, better known under the name *outerplanar decompositions*, have been explored in various algorithmic settings . The main ideas go back to a fundamental article of Baker on approximation algorithms on planar graphs.
The *local tree-width* of a graph $`G=(V,E)`$ is the function $`\text{ltw}^G:`$ that associates with every $`r`$ the maximal tree-width of an $`r`$-neighborhood in $`G`$. More formally, we define the *$`r`$-neighborhood* $`N_r(v)`$ of a vertex $`vV`$ to be the set of all $`wV`$ of distance at most $`r`$ from $`v`$, and we let $`N_r(v)`$ denote the subgraph induced by $`G`$ on $`N_r(v)`$. Then, denoting the tree-width of a graph $`H`$ by $`\text{tw}(H)`$, we let
$$\text{ltw}^G(r):=\text{max}\left\{\text{tw}\left(N_r(v)\right)\right|vV\}.$$
We are mainly interested in classes of graphs of *bounded local tree-width*, that is, classes $`𝒞`$ for which there is a function $`f:`$ such that for all $`G𝒞`$ and $`r`$ we have $`\text{ltw}^G(r)f(r)`$. The class of planar graphs is an example. It has been observed by Eppstein that if a class $`𝒞`$ is closed under taking minors and has bounded local tree-width (Eppstein calls this the “diameter-treewidth property”), then the graphs in $`𝒞`$ admit a decomposition into graphs of small tree-width in the style of the outerplanar decomposition of planar graphs, and the planar-graph algorithms based on this decomposition generalize to graphs in $`𝒞`$. Eppstein gave a nice characterization of such classes; he proved that a minor closed class $`𝒞`$ of graphs has bounded local tree-width if, and only if, it does not contain all apex graphs.
The main graph-theoretic result of this paper, Theorem Theorem 4.2., can be phrased as follows: Let $`𝒞`$ be a minor closed class of graphs that does not contain all graphs. Then all graphs in $`𝒞`$ can be decomposed into a tree of graphs that, after removing a bounded number of vertices, have bounded local tree-width. (Of course the converse is also true, but trivial: If $`𝒞`$ is a minor closed class of graphs such that every graph in $`𝒞`$ admits such a decomposition, then $`𝒞`$ is not the class of all graphs.) The proof of this result is based on a deep structural characterization of graphs with excluded minors due to Robertson and Seymour .
We defer the precise technical statement of our decomposition theorem to Section 4 and turn to its applications now. In this paper, we focus on approximation algorithms. But let me mention that the theorem can also be used to re-prove a result of Alon, Seymour, and Thomas that graphs $`G`$ with an excluded minor have tree-width $`O(\sqrt{|G|})`$ (see Section 6).<sup>1</sup><sup>1</sup>1We have observed this in discussions with Reinhard Diestel and Daniela Kühn.
Actually, the main result of Alon, Seymour, and Thomas’s article is a separator theorem for graphs with an excluded minor, generalizing a well-known separator theorem due to Lipton and Tarjan for planar graphs. These separator theorems have numerous algorithmic applications, among them a polynomial time approximation scheme (PTAS) for the Maximum Independent Set problem on planar graphs and, more generally, classes of graphs with an excluded minor .
A different approach to approximation algorithms on planar graphs is Baker’s technique based on the outerplanar decomposition. It does not only give another PTAS for Maximum Independent Set, but also for other problems, such as Minimum Dominating Set, to which the technique based on the separator theorem does not apply.
We can use our decomposition theorem to extend Baker’s approach to arbitrary classes of graphs with an excluded minor. Our purpose here is to explain the technique and not to give an extensive list of problems to which it applies. We show in detail how to get a PTAS for Minimum Vertex Cover on classes of graphs with an excluded minor and then explain how this PTAS has to be modified to solve the problems Minimum Dominating Set and Maximum Independent Set. It should be no problem for the reader to apply the same technique to other optimization problems.
The paper is organized as follows: In Section 2 we fix our terminology and recall a few basic facts about tree-decompositions of graphs. Local tree-width is introduced in Section 3. In Section 4, we prove our decomposition theorem for classes of graphs with an excluded minor. Approximation algorithms are discussed in Section 5, and in Section 6 we briefly explain two other applications of the decomposition theorem.
## 2 Preliminaries
The vertex set of a graph $`G`$ is denoted by $`V^G`$, the edge set by $`E^G`$. Graphs are always assumed to be finite, simple, and undirected. We write $`vwE^G`$ to denote that there is an edge from $`v`$ to $`w`$. For a subset $`XV^G`$, we let $`X^G`$ denote the induced subgraph of $`G`$ with vertex set $`X`$. We let $`GX:=V^GX^G`$. For graphs $`G`$ and $`H`$, we let $`GH:=(V^GV^H,E^GE^H)`$. We often omit superscripts <sup>G</sup> if $`G`$ is clear from the context.
$`K_n`$ denotes the complete graph with $`n`$ vertices, and for an arbitrary set $`X`$, $`K_X`$ denotes the complete graph with vertex set $`X`$. A vertex set $`XV^G`$ in a graph $`G`$ is a *clique* if $`K_XG`$. The *clique number* $`\omega (G)`$ of a graph $`G`$ is the maximal size of a clique in $`G`$. For a class $`𝒞`$ of graphs, we let $`\omega (𝒞)`$ be the maximum of the clique numbers of all graphs in $`𝒞`$, or $`\mathrm{}`$, if this maximum does not exist.
Note that if $`𝒞`$ is closed under taking subgraphs and is not the class of all graphs, then $`\omega (𝒞)`$ is finite.
### Graph minors.
A *minor* of a graph $`G`$ is a graph $`H`$ that can be obtained from a subgraph of $`G`$ by contracting edges; we write $`HG`$ to denote that $`H`$ is a minor of $`G`$.
Note that $`HG`$ if, and only if, there is a mapping $`h:V^H\text{Pow}(V^G)`$ such that $`h(x)^G`$ is a connected subgraph of $`G`$ for all $`xV^H`$, $`h(x)h(y)=\mathrm{}`$ for $`xyV^H`$, and for every edge $`xyE^H`$ there exists an edge $`uvE^G`$ such that $`uh(x),vh(y)`$. We say that the mapping $`h`$ *witnesses* $`HG`$ and write $`h:HG`$.
A class $`𝒞`$ is *minor closed* if, and only if, for all $`G𝒞`$ and $`HG`$ we have $`H𝒞`$. We call $`𝒞`$ non-trivial if it is not the class of all graphs.
A class $`𝒞`$ is *$`H`$-free* if $`HG`$ for all $`G𝒞`$. We then call $`H`$ an *excluded minor* for $`𝒞`$. Note that a class $`𝒞`$ of graphs has an excluded minor if, and only if, there is an $`n1`$ such that $`𝒞`$ is $`K_n`$-free. Furthermore, this is equivalent to saying that $`𝒞`$ is contained in some non-trivial minor closed class of graphs.
Robertson and Seymour’s *Graph Minor Theorem* states that for every minor closed class $`𝒞`$ of graphs there is a finite set $``$ of graphs such that
$$𝒞=\{GH:HG\}.$$
For a nice introduction to graph minor theory we refer the reader to the last chapter of , a recent survey is .
### Tree-decompositions.
In this paper, we assume trees to be directed from the root to the leaves. If $`tuE^T`$ we call $`u`$ a *child* of $`t`$ and $`t`$ the *parent* of $`u`$. The root of a tree $`T`$ is always denoted by $`r^T`$.
A *tree-decomposition* of a graph $`G`$ is a pair $`(T,(B_t)_{tV^T})`$, where $`T`$ is a tree and $`(B_t)_{tV^T}`$ a family of subsets of $`V^G`$ such that $`_{tV^T}B_t^G=G`$ and for every $`vV^G`$ the set $`\{tvB_t\}`$ is connected. The sets $`B_t`$ are called the *blocks* of the decomposition. The *width* of $`(T,(B_t)_{tV^T})`$ is the number $`\text{max}\{B_ttV^T\}1`$. The *tree-width* of $`G`$, denoted by $`\text{tw}(G)`$, is the minimal width of a tree-decomposition of $`G`$.
The following lemma collects a few simple and well-known facts about tree-decompositions:
* 1. Let $`(T,(B_t)_{tV^T})`$ be a tree-decomposition of a graph $`G`$ and $`XV^G`$ a clique. Then there is a $`tV^T`$ such that $`XB_t`$.
2. Let $`G,H`$ be graphs such that $`V^GV^H`$ is a clique in both $`G`$ and $`H`$. Then $`\text{tw}(GH)=\text{max}\{\text{tw}(G),\text{tw}(H)\}`$.
3. Let $`G`$ be a graph and $`XV^G`$. Then $`\text{tw}(G)\text{tw}(GX)+|X|`$.
4. Let $`G,H`$ be graphs such that $`HG`$. Then $`\text{tw}(H)\text{tw}(G)`$.
Throughout this paper, for a tree-decomposition $`(T,(B_t)_{tV^T})`$ and $`tT\{r^T\}`$ with parent $`s`$ we let $`A_t:=B_tB_s`$. We let $`A_{r^T}:=\mathrm{}`$.
The *adhesion* of $`(T,(B_t)_{tV^T})`$ is the number
$$\text{ad}(T,(B_t)_{tV^T}):=\text{max}\{A_ttV^T\}.$$
The *torso* of $`(T,(B_t)_{tV^T})`$ at $`tV^T`$ is the subgraph
$$[B_t]:=B_t^GK_{A_t}\underset{u\text{ child of }t}{}K_{A_u},$$
or equivalently, the subgraph with vertex set $`B_t`$ in which two vertices are adjacent if, and only if, either they are adjacent in $`G`$ or they both belong to a block $`B_u`$, where $`ut`$. $`(T,(B_t)_{tV^T})`$ is a tree-decomposition of $`G`$ *over* a class $``$ of graphs if all its torsos belong to $``$.
Note that the adhesion of a tree-decomposition over $``$ is bounded by $`\omega ()`$. Actually, it can be easily seen that if a graph has a tree-decomposition over a minor-closed class $``$ then it has a tree-decomposition over $``$ of adhesion at most $`\omega ()1`$.
### Path decompositions.
A *path-decomposition* of a graph $`G`$ is a tree decomposition where the underlying tree is a path. Of course we can always assume that the path $`P`$ of a path decomposition $`(P,(B_p)_{pP})`$ has vertex set $`V^P=\{1,\mathrm{},m\}`$, for some $`m`$, and that the vertices occur on $`P`$ in their natural order (that is, we have $`i(i+1)E^P`$ for $`1i<m`$).
* Let $`G,H`$ be graphs and $`(\{1,\mathrm{},m\},(B_i)_{1im})`$ a path-decomposition of $`H`$ of width $`k`$. Let $`x_1\mathrm{}x_m`$ be a path in $`G`$ such that $`x_iB_i`$ for $`1im`$ and $`V^GV^H=\{x_1,\mathrm{},x_m\}`$. Then $`\text{tw}(GH)(\text{tw}(G)+1)(k+1)1.`$
Proof: Let $`(T,(C_t)_{tV^T})`$ be a tree-decomposition of $`G`$. Then $`(T,(C_t^{})_{tV^T})`$ with
$$C_t^{}=C_t\underset{\begin{array}{c}1im,\\ x_iC_t\end{array}}{}B_i$$
is a tree-decomposition of $`GH`$. $`\mathrm{}`$
## 3 Local tree-width
The distance $`d^G(x,y)`$ between two vertices $`x,y`$ of a graph $`G`$ is the length of the shortest path in $`G`$ from $`x`$ to $`y`$. For $`r1`$ and $`xG`$ we define the *$`r`$-neighborhood* around $`x`$ to be $`N_r^G(x):=\{yV^Gd^G(x,y)r\}`$.
* 1. The *local tree-width* of a graph $`G`$ is the function $`\text{ltw}^G:`$ defined by
$$\text{ltw}^G(r):=\text{max}\{\text{tw}(N_r^G(x))xV^G\}.$$
2. A class $`𝒞`$ of graphs has *bounded local tree-width* if there is a function $`f:`$ such that $`\text{ltw}^G(r)f(r)`$ for all $`G𝒞`$, $`r`$.
$`𝒞`$ has *linear local tree-width* if there is a $`\lambda `$ such that $`\text{ltw}^G(r)\lambda r`$ for all $`G𝒞`$, $`r`$.
* Let $`G`$ be a graph of tree-width at most $`k`$. Then $`\text{ltw}^G(r)k`$ for all $`r`$.
* Let $`G`$ be a graph of valence at most $`l`$, for an $`l1`$. Then $`\text{ltw}^G(r)l(l1)^{r1}`$ for all $`r`$.
The planar graph algorithms due to Baker and others that we mentioned in the introduction are based on the following result:
* The class of planar graphs has linear local tree-width. More precisely, for every planar graph $`G`$ and $`r1`$ we have $`\text{ltw}^G(r)3r`$.
In this paper, a *surface* is a compact connected 2-manifold with (possibly empty) boundary. The (orientable or non-orientable) *genus* of a surface $`S`$ is denoted by $`g(S)`$. An *embedding* of a graph $`G`$ in a surface $`S`$ is a mapping $`\mathrm{\Pi }`$ that associates distinct points of $`S`$ with the vertices of $`G`$ and internally disjoint simple curves in $`S`$ with the edges of $`G`$ in such a way that a vertex $`v`$ is incident with an edge $`e`$ if, and only if, $`\mathrm{\Pi }(v)`$ is an endpoint of $`\mathrm{\Pi }(e)`$.
* Let $`S`$ be a surface. Then the class of all graphs embeddable in $`S`$ has linear local tree-width. More precisely, there is a constant $`c`$ such that for all graphs $`G`$ embeddable in $`S`$ and for all $`r0`$ we have $`\text{ltw}^G(r)cg(S)r`$.
In the next subsection, we prove an extension of Proposition Proposition 3.5 (Eppstein ). that forms the bases of our decomposition theorem for graphs with excluded minors.
But before we do so, let me state another result due to Eppstein that characterizes the minor closed classes of graphs of bounded local tree-width. An *apex graph* is a graph $`G`$ that has a vertex $`vV^G`$ such that $`G\{v\}`$ is planar.
* Let $`𝒞`$ be a minor-closed class of graphs. Then $`𝒞`$ has bounded local tree-width if, and only if, $`𝒞`$ does not contain all apex graphs.
It is an interesting open problem whether there is a minor closed class of graphs of bounded local tree-width that does not have linear (or polynomially bounded) local tree-width.
### Almost embeddable graphs.
Let $`S`$ be a surface with non-empty boundary. The boundary of $`S`$ consists of finitely many connected components $`C_1,\mathrm{},C_\kappa `$, each of which is homeomorphic to the cycle $`S^1`$.
We now define a graph $`G`$ to be *almost embeddable* in $`S`$. Roughly, this means that we can obtain $`G`$ from a graph $`G_0`$ embedded in $`S`$ by attaching at most $`\kappa `$ graphs of path-width at most $`\kappa `$ to $`G_0`$ along the boundary cycles $`C_1,\mathrm{},C_\kappa `$ in an orderly way.
This notion plays an important role in the structure theory of graphs with excluded minors, to be outlined in the next subsection.
* Let $`S`$ be a surface with boundary cycles $`C_1,\mathrm{},C_\kappa `$. A graph $`G`$ is *almost embeddable* in $`S`$ if there are (possibly empty) subgraphs $`G_0,\mathrm{},G_\kappa `$ of $`G`$ such that
+ $`G=G_0\mathrm{}G_\kappa `$,
+ $`G_0`$ has an embedding $`\mathrm{\Pi }`$ in $`S`$,
+ $`G_1,\mathrm{},G_\kappa `$ are pairwise disjoint,
+ for $`1i\kappa `$, $`G_i`$ has a path decomposition $`(\{1,\mathrm{},m_i\},(B_j^i)_{1jm_i})`$ of width at most $`\kappa `$,
+ for $`1i\kappa `$ there are vertices $`x_1^i,\mathrm{},x_{m_i}^iV^{G_0}`$ such that $`x_j^iB_j^i`$ for $`1jm_i`$ and $`V^{G_0}V^{G_i}=\{x_1^i,\mathrm{},x_{m_i}^i\}`$,
+ for $`1i\kappa `$, we have $`\mathrm{\Pi }(V^{G_0})C_i=\{\mathrm{\Pi }(x_1^i),\mathrm{},\mathrm{\Pi }(x_{m_i}^i)\}`$, and the points $`\mathrm{\Pi }(x_1^i),\mathrm{},\mathrm{\Pi }(x_{m_i}^i)`$ appear on $`C_i`$ in this order (either if we walk clockwise or anti-clockwise).
* Let $`S`$ be a surface. Then the class of all graphs almost embeddable in $`S`$ has linear local tree-width.
Proof: Let $`G`$ be a graph that is almost embeddable in $`S`$. We use the notation of Definition Definition 3.7.. Let $`H_0`$ be the graph obtained from $`G_0`$ by adding new vertices $`z_1,\mathrm{},z_\kappa `$, and edges $`(z_i,x_j^i)`$, $`(x_j^i,x_{j+1}^i)`$, and $`(x_\kappa ^i,x_1^i)`$, for $`1i\kappa `$, $`1jm_i`$ (see Figure 1). Clearly, $`H_0`$ is still embeddable in $`S`$. For $`1i\kappa `$ we let $`H_i:=H_0G_1\mathrm{}G_i`$.
Let $`\lambda `$ such that for every graph $`G`$ embedabble in $`S`$ and every $`r`$ we have $`\text{ltw}^G(r)\lambda r`$ (such a $`\lambda `$ exists by Theorem Proposition 3.5 (Eppstein ).). For $`r`$ we let $`f_0(r):=\lambda r`$ and, for $`i`$, we let $`f_i(r):=(f_{i1}(r+1)+1)(\kappa +1)1`$. Then $`f_i`$ is a linear function for every $`i`$.
By induction on $`i0`$ we shall prove that for every $`r`$ and $`xV^{H_i}`$ we have
$$\text{tw}\left(N_r^{H_i}(x)\right)f_i(r).$$
(1)
For $`i=0`$, this is immediate. So we assume that $`i1`$ and that we have proved (1) for $`i1`$.
For all $`xH_i`$, we either have $`N_r^{H_i}(x)H_{i1}`$, or $`N_r^{H_i}(x)G_i`$, or $`N_r^{H_i}\{x_1^i,\mathrm{},x_{m_i}^i\}\mathrm{}`$.
If $`N_r^{H_i}(x)V^{H_{i1}}`$ then $`\text{tw}\left(N_r^{H_i}(x)^{H_i}\right)f_{i1}(r)f_i(r)`$.
If $`xV^{H_{i1}}`$ and $`N_r^{H_i}(x)V^{H_{i1}}`$, then $`N_{r1}^{H_i}(x)\{x_1^i,\mathrm{},x_{m_i}^i\}\mathrm{}`$. By the construction of $`H_0`$, this implies $`z_iN_r^{H_{i1}}(x)`$ and thus $`\{x_1^i,\mathrm{},x_{m_i}^i\}N_{r+1}^{H_{i1}}(x)`$.
By Lemma Lemma 2.2. and the induction hypothesis we get
$`\text{tw}\left(N_r^{H_i}(x)^{H_i}\right)`$ $`\text{tw}\left(N_{r+1}^{H_{i1}}(x)V^{G_i}^{H_i}\right)`$
$`(f_{i1}(r+1)+1)(\kappa +1)1=f_i(r).`$
If $`xV^{G_i}`$, then $`N_r^{H_i}(x)V^{H_{i1}}N_{r+1}^{H_{i1}}(z_i)`$. Thus by Lemma Lemma 2.2. and the induction hypothesis we have
$`\text{tw}\left(N_r^{H_i}(x)^{H_i}\right)`$ $`\text{tw}\left(N_{r+1}^{H_{i1}}(z_i)V^{G_i}^{H_i}\right)`$
$`(f_{i1}(r+1)+1)(\kappa +1)1=f_i(r).`$
$`\mathrm{}`$
Note that the local tree-width of a graph is not minor-monotone (that is, $`HG`$ does not imply $`\text{ltw}^H(r)\text{ltw}^G(r)`$ for all $`r`$). However, we do have
$$HG\text{ltw}^H\text{ltw}^G.$$
(2)
* Let $`S`$ be a surface. Then the class of all minors of graphs almost embeddable in $`S`$ has linear local tree-width.
Proof: Recall the proof of Proposition Proposition 3.8.. We use the same notation here. Suppose $`G^{}`$ is a minor of $`G`$. We can assume that $`G^{}`$ is a subgraph of a graph $`G^{\prime \prime }`$ obtained from $`G`$ only by contracting edges. Because of (2) we can even assume that $`G^{}=G^{\prime \prime }`$.
Let $`X=\{x_j^i1i\kappa ,1jm_i\}`$. Contracting edges with at least one endpoint not in $`X`$ is unproblematic, because the resulting graph is still almost embeddable in $`S`$.
So we can further assume that $`G^{}`$ is obtained from $`G`$ by contracting edges $`e_1`$, $`\mathrm{}`$, $`e_n`$ with both endpoints in $`X`$. Let $`H:=H_\kappa `$ (the graph obtained from $`G`$ by adding the vertices $`z_i`$ and corresponding edges as in Figure 1). Let $`H^{}`$ be the graph obtained from $`H`$ by contracting $`e_1,\mathrm{},e_n`$, and let $`h:H^{}H`$ witness these edge contractions.
The key observation is that for all $`x,yV^H^{}`$ and $`uh(x),vh(y)`$ we have
$$d^H(u,v)d^H^{}(x,y)+3\kappa 1$$
(3)
(no matter how large $`n`$ is). To see this, let $`P^{}`$ be a shortest path from $`x`$ to $`y`$ in $`H^{}`$. Let $`P`$ be a path from $`u`$ to $`v`$ in $`H`$ such that $`P^{}`$ is obtained from $`P`$ by contracting the edges $`e_1\mathrm{},e_n`$. Let us call such an edge an $`(i,j)`$-edge if it connects a vertex in $`\{x_1^i,\mathrm{},x_{m_i}^i\}`$ with a vertex in $`\{x_1^j,\mathrm{},x_{m_j}^j\}`$. Suppose that $`P=w_1\mathrm{}w_r`$. For $`1i\kappa `$, let $`w_s`$ and $`w_t`$, where $`1str`$, be the first and last vertex from $`\{x_1^i,\mathrm{},x_{m_i}^i\}`$ on $`P`$. If $`s<t`$ we replace the interval $`w_s\mathrm{}w_t`$ in $`P`$ by $`w_sz_iw_t`$. Doing this for $`1i\kappa `$ we obtain a new path $`Q`$ from $`u`$ to $`v`$ in $`H`$. This path $`Q`$ contains no at most $`2\kappa `$ edges that are not on $`P`$ and no $`(i,i)`$-edges. Furthermore, for $`1i<jn`$ the number of $`(i,j)`$-edges on $`Q`$ is at most $`(\kappa 1)`$. Because assume that $`Q`$ contains at least $`\kappa `$ such edges. Then there would be a “cycle” $`i=i_1,i_2,\mathrm{},i_l=i`$ such that for $`1j<l`$, $`Q`$ contains an $`(i_j,i_{j+1})`$-edge. However, this cycle would have been removed while transforming $`P`$ to $`Q`$.
Hence $`\text{length}(Q)\text{length}(P^{})+3\kappa 1`$, which proves (3).
(3) implies that for all $`r0`$, $`xV^H^{}`$, and $`uh(x)`$ we have
$$N_r^H^{}(x)N_{r+3\kappa 1}^H(u).$$
(4)
To see this, let $`yN_r^H^{}(x)`$. Then for all $`vh(y)`$, by (3) we have $`vN_{r+3\kappa 1}^H(u)`$. Thus $`h(N_r^H^{}(x))\text{Pow}(N_{r+3\kappa 1}^H(u))`$. Therefore the restriction of $`h`$ to $`N_r^H^{}(x)`$ witnesses $`N_r^H^{}(x)N_{r+3\kappa 1}^H(u)`$. This proves (4).
By (1) and (4) we get $`\text{tw}(N_r^H^{}(x))f_\kappa (r+3\kappa 1)`$. The statement of the lemma follows. $`\mathrm{}`$
## 4 Graphs with excluded minors
The following deep structure theorem for $`K_n`$-free graphs plays a central role in the proof of the Graph Minor Theorem. For a surface $`S`$ and $`\mu `$ we let $`𝒜(S,\mu )`$ be the class of all graphs $`G`$ such that there is an $`XV^G`$ with $`X\mu `$ such that $`GX`$ is almost embeddable in $`S`$.
* For every $`n`$ there exist $`\mu `$ and surfaces $`S,S^{}`$ such that all $`K_n`$-free graphs have a tree-decomposition over $`𝒜(S,\mu )𝒜(S^{},\mu )`$.
Further details concerning this theorem can be found in .
For $`\lambda ,\mu 0`$ we let
$`(\lambda )`$ $`:=\{GHGr0:\text{ltw}^H(r)\lambda r\},`$
$`(\lambda ,\mu )`$ $`:=\{GXV^G:\left(X\mu GX(\lambda )\right)\}.`$
Note that $`(\lambda ,\mu )`$ is minor closed and that $`\omega ((\lambda ,\mu ))=\lambda +\mu +1`$. Thus a tree-decomposition over $`(\lambda ,\mu )`$ has adhesion at most $`\lambda +\mu +1`$.
* Let $`𝒞`$ be a class of graphs with an excluded minor. Then there exist $`\lambda ,\mu `$ such that all $`G𝒞`$ have a tree-decomposition over $`(\lambda ,\mu )`$.
Proof: This follows immediately from Theorem Theorem 4.1 (Robertson and Seymour ). and Proposition Proposition 3.9.. $`\mathrm{}`$
For algorithmic applications we have in mind, Theorem Theorem 4.2. alone is not enough; we also have to compute a tree-decomposition of a given graph over $`(\lambda ,\mu )`$. Fortunately, Robertson and Seymour have proved another deep result that helps us with this task:
* Every minor closed class of graphs has a polynomial time membership test.
* Let $`𝒞`$ be a minor closed class of graphs.
Then there is a polynomial time algorithm that computes, given a graph $`G`$, a tree-decomposition of $`G`$ over $`𝒞`$, or rejects $`G`$ if no such tree-decomposition exists.
Proof: Note that the class $`𝒯`$ of all graphs that have a tree-decomposition over $`𝒞`$ is minor closed. Thus by Theorem Theorem 4.3 (Robertson and Seymour ). we have polynomial time membership tests for both $`𝒞`$ and $`𝒯`$.
Without loss of generality, we can assume that $`𝒞`$ is not the class of all graphs. Thus the clique number $`\omega :=\omega (𝒞)`$ is finite. Recall that every tree-decomposition over $`𝒞`$ has adhesion at most $`\omega `$. Our algorithm uses the following observation to recursively construct a tree-decomposition of the input graph $`G`$:
> $`G𝒯`$ if, and only if, $`G𝒞`$ or there is a set $`XV^G`$ such that $`|X|\omega `$, $`GX`$ has at least two connected components, and for all components $`C`$ of $`GX`$ we have $`XC^GK_X𝒯`$.
We omit the details. $`\mathrm{}`$
In particular, we are going to apply this result to the minor closed classes $`(\lambda ,\mu )`$.
## 5 Approximation algorithms
### Optimization problems.
An *NP-optimization problem* is a tuple $`(I,S,C,\text{opt})`$, consisting of a polynomial time decidable set $`I`$ of *instances*, a mapping $`S`$ that associates a non-empty set $`S(x)`$ of *solutions* with each $`xI`$ such that the binary relation $`\{(x,y)yS(x)\}`$ is polynomial time computable and there is a $`k`$ such that for all $`xI`$, $`yS(x)`$ we have $`yx^k`$, a polynomial time computable *cost* (or *value*) function $`C:\{(x,y)xI,yS(x)\}`$, and a *goal* $`\text{opt}\{\text{min},\text{max}\}`$.
Given an $`xI`$, we want to find a $`yS(x)`$ such that
$$C(x,y)=\text{opt}(x):=\text{opt}\{C(x,z)zS(x)\}.$$
Let $`xI`$ and $`ϵ>0`$. A solution $`yS(x)`$ for $`x`$ is *$`ϵ`$-close* if
$$(1ϵ)\text{opt}(x)C(x,y)(1+ϵ)\text{opt}(x).$$
A *polynomial time approximation scheme (PTAS)* for $`(I,S,C,\text{opt})`$ is a uniform family $`(A_ϵ)_{ϵ>0}`$ of approximation algorithms, where $`A_ϵ`$ is a polynomial time algorithm that, given an $`xI`$, computes an $`ϵ`$-close solution for $`x`$ in polynomial time. Uniformity means that there is an algorithm that, given $`ϵ`$, computes $`A_ϵ`$.
### The levels of graphs of bounded local tree-width.
For graph $`G`$, a vertex $`vV^G`$, and integers $`ji0`$ we let
$$L_v^G[i,j]:=\{wV^Gid^G(v,w)j\}.$$
To keep the notation uniform, we are actually going to write $`L_v^G[i,j]`$ for arbitrary $`i,j`$, with the understanding that $`L_v^G[i,j]:=\mathrm{}`$ for $`i>j`$ and $`L_v^G[i,j]:=L_v^G[0,j]`$ for $`i0`$.
* Let $`\lambda `$. Then for all $`G(\lambda )`$, $`vV^G`$, and $`i,j`$ with $`ij`$ we have $`\text{tw}\left(L_v^G[i,j]\right)\lambda (ji+1)`$.
Proof: First note that $`L_v^G[1,j]L_v^G[0,j]=N_j^G(v)`$, thus the claim holds for $`i1`$. For $`i2`$, consider the minor $`H`$ of $`G`$ obtained by contracting the connected subgraph $`L_v^G[0,i1]`$ to a single vertex $`v^{}`$. Then we have $`L_v^G[i,j]N_{ji+1}^H(v^{})`$, and the claim follows. $`\mathrm{}`$
### Minimum vertex cover.
Instances of Minimum Vertex Cover are graphs $`G`$, solutions are sets $`XV^G`$ such that for every edge $`vwE^G`$ either $`vX`$ or $`wX`$ (such sets $`X`$ are called *vertex covers*), the cost function is defined by $`C(G,X):=|X|`$, and the goal is min.
* For every $`k1`$, the restriction of Minimum Vertex Cover to instances of tree-width at most $`k`$ is solvable in linear time.
* Let $`𝒞`$ be a class of graphs with an excluded minor. Then the restriction of Minimum Vertex Cover to instances in $`𝒞`$ has a PTAS.
Proof: Applying Theorem Theorem 4.2., we choose $`\lambda ,\mu `$ such that every $`G𝒞`$ has a tree-decomposition over $`(\lambda ,\mu )`$. Let $`ϵ>0`$; we shall describe a polynomial time algorithm that, given a graph $`G𝒞`$, computes an $`ϵ`$-close solution for Minimum Vertex Cover on $`G`$. Uniformity will be clear from our description. Let $`k=\frac{1}{ϵ}`$ and note that $`\frac{k+1}{k}(1+ϵ)`$.
In a first step, let us prove that the restriction of Minimum Vertex Cover to instances in $`(\lambda )`$ has a PTAS.
Let $`G(\lambda )`$ and $`vV^G`$ arbitrary. For $`1ik`$ and $`j0`$ we let $`L_{ij}:=L_v^G[(j1)k+i,jk+i]`$. By Lemma Lemma 5.1., $`\text{tw}(L_{ij})\lambda (k+1)`$.
For $`1ik`$, $`j0`$ let $`X_{ij}`$ be a minimal vertex cover of $`L_{ij}`$. We let $`X_i:=_{j0}X_{ij}`$. Then $`X_i`$ is a vertex cover of $`G`$. Let $`X_{\text{min}}`$ be a minimal vertex cover for $`G`$. We have $`|X_{ij}||X_{\text{min}}L_{ij}|`$, because $`X_{\text{min}}L_{ij}`$ is also a vertex cover of $`L_{ij}`$. Hence
$$\underset{i=1}{\overset{k}{}}|X_i|\underset{i=1}{\overset{k}{}}\underset{j0}{}|X_{ij}|\underset{i=1}{\overset{k}{}}\underset{j0}{}|L_{ij}X_{\text{min}}|(k+1)|X_{\text{min}}|.$$
The last inequality follows from the fact that every $`vV^G`$ is contained in at most $`(k+1)`$ (successive) sets $`L_{ij}`$.
Choose $`m,1mk`$ such that $`|X_m|=\text{min}\{|X_1|,\mathrm{},|X_k|\}`$. Then
$$|X_m|\frac{k+1}{k}|X_{\text{min}}|(1+ϵ)|X_{\text{min}}|.$$
Since the $`X_{ij}`$ can be computed in polynomial time by Lemma Lemma 5.2 ()., $`X_m`$ can also be computed in polynomial time.
In a second step, we show how to extend this approximation algorithm to classes $`(\lambda ,\mu )`$ for $`\lambda ,\mu 0`$. Let $`G(\lambda ,\mu )`$ and $`UV^G`$ such that $`|U|\mu `$ and $`H:=GU(\lambda ,0)`$. The following extension of Lemma Lemma 5.2 (). can be proved by standard dynamic programming techniques (cf. ):
* For every $`k0`$, the following problem can be solved in linear time: Given a graph $`G`$, a subset $`UV^G`$ such that $`\text{tw}(GU)k`$, and a subset $`YU`$, compute a set $`XV^GU`$ of minimal order such that $`XY`$ is a vertex cover of $`G`$, if such a set exists, or reject otherwise.
For every $`YU`$ we shall compute an $`X(Y)\text{Pow}(V^GU)\{\}`$ such that either $`X(Y)Y`$ is a vertex cover of $`G`$ and
$$|X(Y)|(1+ϵ)\text{min}\{|X|XV^GU,XY\text{ vertex cover of }G\},$$
or $`X(Y):=`$ if no such $`X(Y)`$ exists. Using Lemma Lemma 5.4. instead of Lemma Lemma 5.2 ()., we can do this analogously to the first step.
Then we choose a $`Y_0U`$ such that $`|X(Y_0)Y_0|`$ is minimal. Here we define $`Z:=`$ for all $`Z`$ and $`||:=\mathrm{}`$. Then clearly $`X(Y_0)Y_0`$ is an $`ϵ`$-close solution for Minimum Vertex Cover on $`G`$. Moreover, since $`|U|\mu `$, there are at most $`2^\mu `$ sets $`YU`$, so $`X(Y_0)Y_0`$ can be computed in polynomial time (remember that $`\mu `$ is a constant only depending on the class $`𝒞`$).
In the third step, we extend our PTAS to graphs that have a tree-decomposition over $`(\lambda ,\mu )`$, i.e. to all graphs in $`𝒞`$.
So let $`G`$ be such a graph. We first compute a tree-decomposition $`(T,(B_t)_{tV^T})`$ of $`G`$ over $`(\lambda ,\mu )`$. Remember that by Lemma Lemma 4.4., this is possible in polynomial time. Recall that $`r^T`$ denotes the root of $`T`$ and that, for every $`tV^T`$ with parent $`u`$, we let $`A_t=B_tB_u`$. For every $`tV^T`$, we let $`S_t`$ be the subtree of $`T`$ with root $`t`$, that is, the subtree with vertex set $`\{st\text{ occurs on the path from }s\text{ to }r^T\}`$. We let $`C_t:=_{sS_t}B_t`$.
Inductively from the leaves to the root, for every node $`tV^T`$ and for every $`YA_t`$ we compute an $`X(t,Y)\text{Pow}(C_tA_t)\{\}`$ such that either $`X(t,Y)Y`$ is a vertex cover of $`C_t`$ and
$$|X(t,Y)|(1+ϵ)\text{min}\{|X|XY\text{ vertex cover of }C_t\},$$
or $`X(t,Y):=`$ if no such vertex set exists. Since a tree-decomposition over $`(\lambda ,\mu )`$ has adhesion at most $`\lambda +\mu +1`$ we have $`|A_t|\lambda +\mu +1`$, thus for every $`tV^T`$ we have to compute at most $`2^{\lambda +\mu +1}`$ sets $`X(t,Y)`$. For the root $`r^T`$ we have $`A_{r^T}=\mathrm{}`$, so $`X(r^T,\mathrm{})`$ is an $`ϵ`$-close solution for Minimum Vertex Cover on $`G`$.
Suppose that $`tV^T`$ and that for every child $`t^{}`$ of $`T`$ we have already computed the family $`X(t^{},)`$. Let $`UB_t`$ such that $`|U|\mu `$ and $`[B_t]U(\lambda )`$. Let $`W:=UA_t`$ and let $`ZW`$. Let $`X_{\text{min}}(Z)\text{Pow}(C_tW)\{\}`$ be a vertex set of minimal order such that $`X_{\text{min}}(Z)Z`$ is a vertex cover of $`C_t`$, or $`X(Z):=`$ if no such vertex set exists.
We show how to compute an $`X(Z)\text{Pow}(C_tW)\{\}`$ such that $`X(Z)Z`$ is a vertex cover of $`C_t`$ and $`|X(Z)|(1+ϵ)|X_{\text{min}}(Z)|`$, if $`X_{\text{min}}(Z)`$, or $`X(Z)=`$ otherwise. Then for every $`YA_t`$ we choose a $`ZW`$ such that $`YZ`$ with minimal $`|X(Z)(ZY)|`$ (among all $`ZY`$) and let $`X(t,Y):=X(Z)`$. Note that, since $`|U|\mu `$, for every $`Y`$ we have to compute at most $`2^\mu `$ sets $`X(Z)`$ to determine $`X(t,Y)`$.
So let us fix a $`ZW`$; we show how to compute $`X(Z)`$ in polynomial time.
If $`W=B_t`$ we let $`X(Z):=_{t^{}\text{ child of }t}X(t^{},A_t^{}Z)`$.
Otherwise, we choose an arbitrary $`vB_tW`$. For $`1ik`$ and $`j0`$ we let $`L_{ij}:=L_v^{[B_t]W}[(j1)k+i,jk+i]`$. Then $`\text{tw}(L_{ij})\lambda (k+1)`$. For $`1ik`$ and every child $`t^{}`$ of $`t`$ there is at least one $`j0`$ such that $`A_t^{}WL_{ij}`$, because $`A_t^{}`$ induces a clique in $`[B_t]`$. Let $`j^{}(i,t^{})`$ be the least such $`j`$ and $`L_{ij}^{}:=L_{ij}_{\begin{array}{c}t^{}\text{ child of }t\\ j^{}(i,t^{})=j\end{array}}C_t^{}A_t^{}`$.
For every $`XL_{ij}`$ we let
$$X^{}:=X\underset{\begin{array}{c}t^{}\text{ child of }t\\ j^{}(i,t^{})=j\end{array}}{}X(t^{},(XZ)A_t^{})$$
We compute an $`X_{ij}L_{ij}`$ with minimal $`|X_{ij}^{}|`$ such that $`X_{ij}Z`$ is a vertex cover of $`L_{ij}W`$ if such a vertex cover exists, and $`X_{ij}=`$ otherwise. The usual dynamic programming techniques on graphs of bounded tree-width show that each $`X_{ij}`$ can be computed in linear time if the numbers $`|X(t^{},Y)|`$ for the children $`t^{}`$ of $`t`$ are given (cf. Lemmas Lemma 5.2 (). and Lemma 5.4. and ). It is important here that every $`A_t^{}W`$ is a clique in $`L_{ij}`$ and thus by Lemma Lemma 2.1.(1) completely contained in a block of every tree-decomposition of $`L_{ij}`$.
We let $`X_i:=_{j0}X_{ij}`$ and $`X_i^{}:=_{j0}X_{ij}^{}`$. Then $`X_i^{}Z`$ is a vertex cover of $`C_t`$, if such a vertex cover exists, and $`X_i=`$ otherwise. We choose an $`i,1ik`$, such that $`|X_i^{}|=\text{min}\{|X_1^{}|,\mathrm{},|X_k^{}|\}`$ and let $`X(Z):=X_i^{}`$. Then $`X(Z)`$ can be computed in polynomial time.
Recall that $`X_{\text{min}}:=X_{\text{min}}(Z)C_tW`$ is a vertex set of minimal order such that $`X_{\text{min}}Z`$ is a vertex cover of $`C_t`$, if such a vertex cover exists, and $`X_{\text{min}}=`$ otherwise. It remains to prove that $`|X(Z)|(1+ϵ)|X_{\text{min}}|`$.
Recall that for every child $`t^{}`$ of $`t`$ we have
$$|X(t^{},(X_{\text{min}}Z)A_t^{})|(1+ϵ)|X_{\text{min}}C_t^{}A_t^{}|.$$
Our construction of the $`X_{ij}`$ and $`X_{ij}^{}`$ guarantees that for $`1ik,j0`$ we have
$$|X_{ij}^{}||X_{\text{min}}L_{ij}|+\underset{\begin{array}{c}t^{}\text{ child of }t\\ j^{}(i,t^{})=j\end{array}}{}|X(t^{},(X_{\text{min}}Z)A_t^{})|.$$
Then
$`k|X(Z)|`$ $`{\displaystyle \underset{i=1}{\overset{k}{}}}|X_i^{}|`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle \underset{j0}{}}|X_{ij}^{}|`$
$``$ $`{\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle \underset{j0}{}}\left(|X_{\text{min}}L_{ij}|+{\displaystyle \underset{\begin{array}{c}t^{}\text{ child of }t\\ j^{}(i,t^{})=j\end{array}}{}}|X(t^{},(X_{\text{min}}Z)A_t^{})|\right)`$
$``$ $`{\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle \underset{j0}{}}\left(|X_{\text{min}}L_{ij}|+{\displaystyle \underset{\begin{array}{c}t^{}\text{ child of }t\\ j^{}(i,t^{})=j\end{array}}{}}(1+ϵ)|X_{\text{min}}C_t^{}A_t^{}|\right)`$
$``$ $`(k+1)|X_{\text{min}}B_t|+k(1+ϵ)|X_{\text{min}}C_tB_t|.`$
This implies $`|X(Z)|(1+ϵ)X_{\text{min}}`$. $`\mathrm{}`$
### Minimum dominating set.
Instances of Minimum Dominating Set are graphs $`G`$, solutions are sets $`XV^G`$ such that for every $`vV^GX`$ there is a $`wX`$ such that $`vwE^G`$ (such sets $`X`$ are called *dominating sets*), the cost function is defined by $`C(G,X):=|X|`$, and the goal is min.
* Let $`𝒞`$ be a class of graphs with an excluded minor. Then the restriction of Minimum Dominating Set to instances in $`𝒞`$ has a PTAS.
Proof: We proceed very similarly to the proof of Theorem Theorem 5.3., the analogous result for Minimum Vertex Cover. Let $`\lambda ,\mu `$ such that every graph in $`𝒞`$ has a tree-decomposition over $`(\lambda ,\mu )`$. Let $`ϵ>0`$ and $`k:=\frac{2}{ϵ}`$.
Again, in the first step we consider the restriction of the problem to input graphs from $`(\lambda )`$. Given such a graph $`G`$, we choose an arbitrary $`vV^G`$. For $`1ik`$ and $`j0`$ we let $`L_{ij}:=L_v^G[(j1)k+i1,jk+i]`$. Then $`\text{tw}(L_{ij})\lambda (k+2)`$. Note that $`L_{ij}`$ and $`L_{i(j+1)}`$ overlap in two consecutive rows, which is different from the proof of Theorem Theorem 5.3.. The *interior* of $`L_{ij}`$ is the set $`L_{ij}^{}:=L_v^G[(j1)k+i,jk+i1]`$.
For $`1ik,j0`$ we let $`X_{ij}L_{ij}`$ be a vertex set of minimal order with the following property:
* For every $`wL_{ij}^{}X_{ij}`$ there is a $`xX_{ij}`$ such that $`(w,x)E^G`$.
Then for $`1ik`$ the set $`X_i:=_{j0}X_{ij}`$ is a dominating set of $`G`$. Let $`m`$ be such that $`|X_m|=\text{min}\{|X_1|,\mathrm{},|X_k|\}`$. Computing $`X_m`$ amounts to solving a variant of Minimum Dominating Set on instances of tree-width at most $`\lambda (k+2)`$; using the usual dynamic programming techniques, this can be done in linear time.
Since for every dominating set $`X`$ of $`G`$ the set $`XL_{ij}`$ has property $`()`$ we have $`X_{ij}XL_{ij}`$. Using this, we can argue as in the proof of Theorem Theorem 5.3. to show that $`X_m`$ is an $`ϵ`$-close solution.
Adapting the second and third step of the proof of Theorem Theorem 5.3., it is straightforward to extend this algorithm to arbitrary input graphs in $`𝒞`$. $`\mathrm{}`$
### Maximum independent set.
Instances of Maximum Independent Set are graphs $`G`$, solutions are sets $`XV^G`$ such that for all $`v,wX`$ we have $`vwE^G`$ (such sets $`X`$ are called *independent sets*), the cost function is defined by $`C(G,X):=|X|`$, and the goal is max.
* Let $`𝒞`$ be a class of graphs with an excluded minor. Then the restriction of Maximum Independent Set to instances in $`𝒞`$ has a PTAS.
Proof: Again we proceed similarly to the proof of Theorem Theorem 5.3.. Let $`\lambda ,\mu `$ such that every graph in $`𝒞`$ has a tree-decomposition over $`(\lambda ,\mu )`$. Let $`ϵ>0`$ and $`k=\frac{1}{ϵ}`$.
We describe how to treat input graphs in $`(\lambda )`$. Following the lines of the proof of Theorem Theorem 5.3., the extension to arbitrary $`G𝒞`$ is straightforward. Let $`G(\lambda )`$ and $`vV^G`$. For $`1ik`$ and $`j0`$ we let $`L_{ij}:=L_v^G[(j1)k+i,jk+i2]`$. Then $`\text{tw}(L_{ij})\lambda (k1)`$. Note that there are no edges between $`L_{ij}`$ and $`L_{i(j+1)}`$.
For $`1ik,j0`$ we let $`X_{ij}`$ be a maximal independent set of $`L_{ij}`$. Then $`X_i:=_{j0}X_{ij}`$ is an independent set of $`G`$. Let $`1mk`$ such that $`|X_m|=\text{max}\{|X_1|,\mathrm{},|X_k|\}`$. Since the restriction of Maximum Independent Set to graphs of bounded tree-width is solvable in linear time, such an $`X_m`$ can be computed in linear time.
Let $`X_{\text{max}}`$ be a maximum independent set of $`G`$. Then for $`1ik`$, $`j0`$ we have $`|X_{ij}||X_{\text{max}}L_{ij}|`$. Thus
$$k|X_m|\underset{i=1}{\overset{k}{}}|X_i|=\underset{i=1}{\overset{k}{}}\underset{j0}{}|X_{ij}|\underset{i=1}{\overset{k}{}}\underset{j0}{}|X_{\text{max}}L_{ij}|(k1)|X_{\text{max}}|,$$
which implies that $`X_m\frac{k1}{k}|X_{\text{max}}|(1ϵ)|X_{\text{max}}|`$. $`\mathrm{}`$
### Other problems.
Our approach can be used to find polynomial time approximation schemes for the restrictions of a number of other problems to classes of graphs with excluded minors, in particular for the other problems considered by Baker . I leave it to the reader to work out the details.
## 6 Other applications of Theorem Theorem 5.3.
### The tree-width of $`K_n`$-free graphs.
We re-prove a theorem of Alon, Seymour, and Thomas that the tree-width of a $`K_n`$-free graph $`G`$ is $`O(\sqrt{|G|})`$. This is joint work with Reinhard Diestel and Daniela Kühn.
* Let $`\lambda `$ and $`G(\lambda )`$. Then $`\text{tw}(G)3\sqrt{\lambda |G|}`$.
Proof: Let $`vV^G`$ arbitrary and, for $`i0`$, $`L_i:=\{wV^Gd^G(v,w)=i\}`$. Let $`m`$ be maximal such that $`L_m`$ is non-empty. We subdivide $`\{1,\mathrm{},m\}`$ into intervals $`I_1,J_1,I_2,\mathrm{},J_{l1},I_l,J_l`$ such that for $`1il`$ we have
* $`|L_j|\sqrt{\lambda |G|}`$ for all $`jI_i`$,
* $`|L_j|>\sqrt{\lambda |G|}`$ for all $`jJ_i`$.
Then $`\text{tw}(_{jI_i}L_j)2\sqrt{\lambda |G|}`$ and $`\text{tw}(_{jJ_i}L_j)\sqrt{\lambda |G|}`$ (because the length of $`J_i`$ is at most $`\sqrt{\frac{|G|}{\lambda }}`$). We can glue the decompositions together by adding to every block of a tree-decomposition of $`J_i`$ the last level of the previous $`I_i`$ and the first level of the next $`I_{i+1}`$ and obtain $`\text{tw}(G)3\sqrt{\lambda |G|}`$. $`\mathrm{}`$
* Let $`\lambda ,\mu `$ and $`G(\lambda ,\mu )`$. Then $`\text{tw}(G)3\sqrt{\lambda |G|}+\mu `$.
* Let $`G`$ be $`K_n`$-free. Then $`\text{tw}(G)O(\sqrt{|G|})`$.
### Deciding first-order properties.
In we give another algorithmic application of Theorem Theorem 4.2.. We show that for every class $`𝒞`$ of graphs with an excluded minor there is a constant $`c>0`$ such that for every property of graphs that is definable in first order logic there is an $`O(|G|^c)`$-algorithm deciding whether a given graph $`G𝒞`$ has this property.
For example, this implies that for every class $`𝒞`$ with an excluded minor there is a constant $`c`$ such that for every graph $`H`$ there is an $`O(|G|^c)`$-algorithm testing whether a given graph $`G𝒞`$ has a subgraph isomorphic to $`H`$.
## 7 Further research
We have never specified the exponents and coefficients of the polynomials bounding the running times of our algorithms; they seem to be enormous. So our algorithms are only of theoretical interest. The first important step towards improving the algorithms would be a practically applicable algorithm for computing tree-decompositions of graphs of small tree-width. On the graph theoretic side, it would probably help to prove Theorem Theorem 4.2. directly without using Robertson’s and Seymour’s Theorem Theorem 4.1 (Robertson and Seymour )..
The traveling salesman problem is another optimization problem that has a PTAS on planar graphs . It would be interesting to see if this problem has a PTAS on class of graphs with an excluded minor.
### Acknowledgements
I thank Reinhard Diestel and Jörg Flum for helful comments on earlier versions of this paper. |
warning/0001/math0001044.html | ar5iv | text | # Mirror symmetry and actions of braid groups on derived categories
## 1. Introduction
In this talk I will attempt to do two things: review a conjectural picture of mirror symmetry, which will no doubt crop up in many other talks, and explain some of its consequences which are proved (independently of the conjectures) in joint work with Mikhail Khovanov and Paul Seidel \[ST, KS\]. This is not the mirror symmetry of Gromov-Witten invariants and variations of Hodge structure (yet) but the more fundamental programmes of Kontsevich \[K\] and Strominger-Yau-Zaslow \[SYZ\] which should eventually lead back to the traditional predictions of mirror symmetry.
Section 3 explains the two programmes and their supposed link in the language of Fourier-Mukai transforms, which are reviewed and explained in Section 2. Though many of the objects used in the models are yet to be defined, we can attempt to deduce some consequences and prove these. In particular, symplectomorphisms of a Calabi-Yau manifold should induce autoequivalences of the bounded derived category $`D^b(X)`$ of coherent sheaves on the mirror $`X`$. We find many such autoequivalences, and actions of the braid group on $`D^b(X)`$ in particular cases (such as on resolutions of singularities). These should be mirror dual to braid groups of symplectomorphisms that arise from configurations of Lagrangian spheres. These spheres are often the vanishing cycles of the smoothing of another, dual, singularity, and we discuss this mirror symmetry of singularities. In Section 5 we mention briefly the relationship of this work to mutations of bundles on Fano manifolds, and finally in Section 6 we outline the faithfulness of the braid group action.
Acknowledgements. My main debts are to Mikhail Khovanov and Paul Seidel who allowed me to talk about this work as if it were my own. Kontsevich and Bridgeland and Maciocia have also discussed the twist (4.1). There are many more acknowledgements in \[ST\]. I would also like to thank the organisers of the 1999 Harvard Winter School on Mirror Symmetry for inviting me to take part.
## 2. Fourier-Mukai transforms
### Function transforms
We begin by introducing Fourier-Mukai transforms by analogy with function transforms. Suppose we are given a family $`F_p`$ of (complex-valued, say) functions or distributions on some space $`V`$, parametrised by some dual space $`\widehat{V}p`$.
Let $`F:V\times \widehat{V}`$ be the universal function (with $`F|_{V\times \{p\}}=F_p`$). For instance we could take $`V=^n,\widehat{V}=(^n)^{}`$ with $`F(x,p)=e^{ip.x}=F_p(x)`$. Similarly on any manifold $`V`$ we can take $`\widehat{V}=V`$ and $`F=\delta _\mathrm{\Delta }`$, the Dirac-delta of the diagonal $`\mathrm{\Delta }`$, with $`F_p(x)=\delta _p(x)`$.
Suppose now that $`\{F_p:p\widehat{V}\}`$ span (some class of) functions $`V`$, and are orthonormal with respect to some inner product (e.g. $`F_p\overline{F}_q=\delta _{pq}`$ in a distributional sense). Then for any $`f`$ in this class, $`f`$ is built up from the $`F_p`$ s, with the coefficient of $`F_p`$ being $`\widehat{f}(p)=_Vf\overline{F}_p`$, and we sum $`\widehat{f}(p)F_p`$ over $`p`$ to regain $`f`$. That is,
$$f=_{\widehat{V}}\left[_Vf\overline{F}_p\right]F_p𝑑p,$$
so $`f=(\widehat{f})^{}`$, where is the dual transform $`g^{}(x)=_{\widehat{V}}g(p)F(p,x)𝑑p`$.
For instance for $`F=e^{ip.x}`$ we get the Fourier Inversion theorem (though theorem is a little strong for the above treatment) for the Fourier transform $`\widehat{}`$ and its inverse . Similarly taking $`F=\delta _\mathrm{\Delta }`$ on $`V\times V`$ gives two applications of the identity map: $`\widehat{f}(p)=_Vf\delta _p=f(p)`$.
The way to look at this most relevant to sheaves is via the diagram
Then the transforms are
$`\widehat{f}`$ $`=`$ $`(\pi _2)_{}[\pi _1^{}f.\overline{F}],`$
(2.1) $`g^{}`$ $`=`$ $`(\pi _1)_{}[\pi _2^{}g.F].`$
Thus we pull up a function from $`V`$ to the product, multiply by $`\overline{F}`$, and push down the fibres $`V_p`$ of $`\pi _2`$ ($`\pi _{}`$ means integrate down the fibres of $`\pi `$): at each $`p`$ we integrate $`f`$ against $`\overline{F}_p`$ to get $`\widehat{f}(p)`$. Notice also that we can recover the functions $`F_p`$ parametrised by $`p`$ as the inverse transform of the Dirac-delta $`\delta _p`$,
(2.2)
$$F_p=\delta _p^{}.$$
We can also do this in a family. For instance on $`V\times Z`$ we can carry an extra parameter $`zZ`$ and take the usual transform at each point $`z`$, i.e. the transform just with respect to the variables $`x`$ and $`p`$,
(2.3)
$$f(x,z)\hat{}\widehat{f}(p,z).$$
This is just a relative transform on $`(V\times Z)\times _Z(\widehat{V}\times Z)=V\times \widehat{V}\times Z`$. Alternatively we can write it as a strict transform on $`(V\times Z)\times (\widehat{V}\times Z)`$. The function parametrised by $`(p,z)\widehat{V}\times Z`$ is $`F_p.\delta _z`$, where $`F_p`$ is (the pull-back to $`V\times Z`$ of) one of the family of functions on $`V`$, and $`\delta _z`$ is the Dirac-delta concentrated on $`V\times \{z\}`$. Again these functions
$$\{F_{p,z}=F_p.\delta _z:(p,z)\widehat{V}\times Z\}$$
form an orthonormal set, giving a transform with universal function $`F\delta _\mathrm{\Delta }`$, where $`F`$ is the pull-back of the universal function on $`V\times \widehat{V}`$ and $`\mathrm{\Delta }=V\times \widehat{V}\times Z(V\times Z)\times (\widehat{V}\times Z)`$ is the diagonal. The transform $`\widehat{f}(p,z)`$ is of course the same as (2.3).
Similarly we can have a non-trivial family $`XZ`$, with dual family $`\widehat{X}Z`$, and a relative transform on $`X\times _Z\widehat{X}`$ or a strict transform on $`X\times \widehat{X}`$. Again they are the same, with the universal function on $`X\times _Z\widehat{X}`$ pushed forward via the diagonal map $`X\times _Z\widehat{X}X\times \widehat{X}`$ to a function supported on this diagonal. This will seem much more natural in the setting of sheaves.
### Sheaf transforms
Suppose now we have a family $`F_p`$ of vector bundles, or coherent sheaves, or complexes of sheaves, on some complex manifold $`X`$, parametrised by some dual variety $`\widehat{X}p`$. We will also need to assume the existence of a universal object $`F`$ on $`X\times \widehat{X}`$ whose restriction to each $`X\times \{p\}`$ is $`F_p`$.
We will define the inverse transform first (denoted here by the unfortunate notation ; this is not the dual of a sheaf but its transform). We would like to think of a sheaf/bundle/complex $``$ on $`\widehat{X}`$ as giving a distribution of the $`F_p`$-components of a sheaf $`^{}`$ on $`X`$, i.e. $`^{}`$ should be the sum over $`p\widehat{X}`$ of the sheaves $`|_pF_p`$. So with respect to the diagram
we would like to set, as in (2.1),
$$^{}=(\pi _1)_{}\left[(\pi _2^{})F\right].$$
Since a sheaf assigns groups of sections to open sets, the right notion of pushdown $`(\pi _1)_{}`$ (the sum of all of the sections over $`p\widehat{X}`$, the analogue of integration down the fibres) is to assign to an open set $`U`$ all sections on $`\pi _1^1(U)=U\times \widehat{X}`$. In fact numbers of sections down fibres can jump, or be zero, due to the presence of higher cohomology. So we actually take the right derived functor $`𝐑(\pi _1)_{}`$ described in my last talk – this is the complex of sheaves obtained by resolving by injective sheaves and applying $`(\pi _1)_{}`$ to this complex. For similar reasons the tensor product should be the derived tensor product $`\stackrel{𝐋}{}`$, and the restriction to fibres $`𝒪_{X\times \{p\}}`$ we have mentioned before should also be taken in this derived sense.
Thus the correct setting is the bounded derived category of coherent sheaves $`D^b`$, with a universal object $`FD^b(X\times \widehat{X})`$. We set, for $`D^b(\widehat{X})`$,
(2.4)
$$^{}=𝐑\pi _1\left[(\pi _2^{})\stackrel{𝐋}{}F\right];$$
in most of the simple cases we shall be interested in one can think of vector bundles rather than complexes of sheaves, and normal pushdown and tensor product.
For instance the inverse transform of the structure sheaf $`𝒪_p`$ of a point $`p\widehat{X}`$ is
(2.5)
$$𝒪_p^{}=𝐑\pi _1\left[𝒪_{X\times \{p\}}\stackrel{𝐋}{}F\right]=\pi _1F_p=F_p,$$
the (derived) restriction of $`F`$ to $`X\times \{p\}`$, i.e. the object of $`D^b(X)`$ parametrised by $`p\widehat{X}`$; compare (2.2). Similarly the inverse transform of $`_i𝒪_{p_i}`$ is $`_iF_{p_i}`$.
The transform from $`D^b(X)`$ to $`D^b(\widehat{X})`$ is given by the formula (c.f. (2.1))
$$\widehat{}=𝐑\pi _2[(\pi _1^{}()\stackrel{𝐋}{}F^{}][n],$$
where here we have denoted by $`F^{}`$ the derived dual of $`F`$, $`𝐑om(F,𝒪_{X\times \widehat{X}})`$, and $`n=`$ dim $`X`$. Thus for $`D^b(X)`$ we have
$$\widehat{}=𝐑\pi _2\left[𝐑om(F,\pi _1^{})\right][n].$$
For and $`\widehat{}`$ to be actual inverses (equivalences is the correct categorical notion) we need the objects $`F_p=F\stackrel{𝐋}{}𝒪_{X\times \{p\}}`$ to be orthogonal \[BO\],
(2.6)
$$𝐑\mathrm{Hom}(F_p,F_q)0,pq,$$
orthonormal, in the sense that they are *simple*,
(2.7)
$$\mathrm{Hom}(F_p,F_p)=.\mathrm{id},$$
and to satisfy the dimension and partial Calabi-Yau conditions \[Ma\]
$`\mathrm{dim}X=\mathrm{dim}\widehat{X},`$
(2.8) $`F_p\omega _XF_pp.`$
Recently Bridgeland \[Br\] has shown that for $`X`$ and $`\widehat{X}`$ smooth projective varieties these conditions (2.6), (2.7), (2.8) are also sufficient for the Fourier-Mukai transform given by $`FD^b(X\times \widehat{X})`$ to be an equivalence. Here $`\omega _X`$ is the canonical bundle of $`X`$, and we are requiring this to be trivial on the support of each $`F_p`$; recall that $`X`$ is Calabi-Yau if $`\omega _X`$ is globally trivial.
###### Example 2.1.
Take $`\widehat{X}=X`$ and $`F=𝒪_\mathrm{\Delta }`$, the structure sheaf of the diagonal parametrising the sheaves $`F_x=𝒪_x`$. Then the transform is the identity, and the partial Calabi-Yau condition (2.8) is vacuous since $`𝒪_x\omega _X`$ is trivially isomorphic to $`𝒪_x`$.
###### Example 2.2.
The original Fourier-Mukai transform \[Mu\]. Let $`T`$ be an elliptic curve, and let $`\widehat{T}t`$ be its Jacobian parametrising degree 0 line bundles $`L_t`$ on $`T`$. We take $`F=L`$ to be a Poincaré line bundle on $`T\times \widehat{T}`$ (for instance if we fix a basepoint $`t_0T`$ then we can identify $`T`$ with $`\widehat{T}`$, and $`L`$ is given by the divisor $`\mathrm{\Delta }(\{t_0\}\times T)T\times T`$).
Then the $`L_t`$ s form an orthonormal set:
$$H^0(L_s^{}L_t)=H^1(L_s^{}L_t)=\{\begin{array}{cc}\hfill & s=t,\hfill \\ 0\hfill & st.\hfill \end{array}$$
So as $`T`$ is Calabi-Yau we get an invertible transform
$$\widehat{}=𝐑\pi _1\left[L\pi _2^{}\right]$$
with fibre over the generic point $`t\widehat{T}`$ (at which base-change holds) “$`H^0(L_t)H^1(L_t)`$” (really a complex with these two cohomologies). Thus, when $`H^1`$ vanishes, we replace $``$ by its spectrum of sections $`\{H^0(L_t):t\widehat{T}\}`$ of different twists. The inverse transform uses the dual line bundle, and the set-up is symmetric between $`T`$ and $`\widehat{T}`$.
Then using (2.5) we see that the transform of the structure sheaf $`𝒪_t`$ of a point $`t\widehat{T}`$ is the corresponding line bundle $`L_t`$, while $`r`$ points transform to the appropriate rank $`r`$ bundle (which is the sum of $`r`$ line bundles). We depict this as follows, drawing a basis of the fibre of the vector bundle over the right hand torus.
Thus, inversely, we have an algebraic gadget that takes a degree 0 semistable bundle on $`T`$ and splits it into its constituent pieces according to Atiyah’s classification; for every $`L_t`$ factor we get a point $`t\widehat{T}`$. The non-trivial extensions correspond to structure sheaves of double points, etc.
Similarly we can do this in a family. Consider an elliptic fibration $`XZ`$, possibly with singular fibres. Suppose we also have a section $`s`$ so that we can identify $`X`$ with its relative Jacobian. Then there is a Poincaré sheaf $``$ on $`X\times _ZX`$ corresponding to the Weil divisor $`\mathrm{\Delta }s(Z)\times _ZX`$. Rather than doing a relative transform on this singular space, we do a strict transform on $`X\times X`$ by pushing forward $``$ to $`\iota _{}`$ via $`\iota :X\times _ZXX\times X`$. Then setting
$$\widehat{}=𝐑\pi _1\left[\iota _{}\pi _2^{}\right]$$
we get a transform $`D^b(X)D^b(X)`$, a family version of the previous example, giving pictures like:
Multisection $`C`$ (+ line bundle on it) $``$ Vector bundle, deg 0 on fibres
where the covering degree $`r`$ of the multisection (its intersection with the fibre) is the rank $`r`$ of the vector bundle.
Thus vector bundles correspond to “spectral covers” in the usual way – the cover gives $`r`$ points on each fibre yielding the sum of $`r`$ line bundles on the dual fibre. Globally we patch these together in the base direction using the line bundle on $`C`$, and get a non trivial vector bundle which is not globally a sum of line bundles (because $`C`$ is not globally $`r`$ disjoint sections).
The rigorous statement is that there is an autoequivalence of $`D^b(X)`$ given by this transform. It is invertible since the universal sheaf $`\iota _{}`$ is concentrated on torus fibres, about which $`X`$ is locally Calabi-Yau so that $`\iota _{}\omega _X\iota _{}`$. For us, the point is simply that there is an algebraic gadget that converts spectral covers like $`C`$, with line bundles on them, into vector bundles (degree 0 and semistable on the fibres), and vice-versa. It does this for free, without any analysis of the singular fibres, since $`X\times X`$ is smooth (even though $`X\times _ZX`$ is not). This is at the expense of introducing the derived category, of course, but as we see here in many cases the transform takes sheaves to sheaves instead of a complex of sheaves.
## 3. Kontsevich vs. Strominger-Yau-Zaslow
We now turn to mirror symmetry and the two competing conjectural theories. Strominger-Yau-Zaslow \[SYZ\] suggest that (in an appropriate complex structure) a Calabi-Yau $`n`$-fold $`X`$ should admit a fibration by special Lagrangian tori ($`T^nX`$ is Lagrangian if the symplectic (Kähler) form restricts identically to zero on $`T^n`$, and special if the restriction of the imaginary part of the Calabi-Yau $`(n,0)`$-form is identically zero) with a special Lagrangian section. In this case the mirror $`\widehat{X}`$ should be the dual torus fibration.
Kontsevich \[K\], on the other hand, conjectures that there should be a natural exact equivalence of triangulated categories (an exact equivalence is one which preserves the distinguished triangles) between $`D^b(X)`$ and $`D^b(\mathrm{Fuk}(\widehat{X}))`$. The second category here is the derived category of the Fukaya category of the symplectic manifold $`\widehat{X}`$, and as such is not yet properly constructed (though see Fukaya’s talk for more progress on this). Other talks will explain more about the Fukaya category; all we need to know is that it is constructed from only the symplectic geometry of $`\widehat{X}`$, using (graded) Lagrangian submanifolds, local systems on them, and their Fukaya-Floer homology. In particular every Lagrangian cycle $`L`$ in $`\widehat{X}`$ (plus a grading), with a flat line bundle on it, should give an object in $`D^b(\mathrm{Fuk}(\widehat{X}))`$. The corresponding object $`_L`$ in $`D^b(X)`$ should have the same Homs, so that $`HF^{}(L,L)`$ should be quasi-isomorphic to $`𝐑\mathrm{Hom}(_L,_L)`$.
The (again conjectural) correspondence between the two pictures is now folklore and has been discussed by many people, see for instance \[AP, Ty\]. The basic idea is that a Lagrangian multisection $`L`$ in the fibration $`\widehat{X}`$ (with a flat line bundle on it, and intersection with the fibre $`r`$) should correspond to a rank $`r`$ holomorphic vector bundle on $`X`$ by an analytical version of the Fourier-Mukai transform, giving a diagram like the one above. That is, the intersection of $`L`$ with a fibre gives $`r`$ points corresponding to $`r`$ line bundles on the dual torus, as before. *Special* Lagrangian sections should perhaps correspond to bundles with Hermitian-Yang-Mills connections (i.e. *stable* bundles) as suggested in \[Va\] (both special and stable are stability conditions on the objects on either side under a natural group action, so this makes sense).
In 2 complex dimensions we can rigorously carry out this procedure, since we have the tools of algebraic geometry and Fourier-Mukai to deal with the singularities. We also have Yau’s solution of the Calabi conjecture, which gives a hyperkähler metric on a compact Calabi-Yau surface (i.e. a $`K3`$ or $`T^4`$). Thus for every complex structure $`I`$ there is a conjugate quaternionic partner $`J`$, and if we rotate the complex structure from $`I`$ to $`J`$ then the special Lagrangian cycles become the complex curves on $`X`$. Thus, after this hyperkähler rotation, the SYZ conjecture is concerned with an elliptically fibred surface with a section, and the mirror should be the dual fibration. Thus in this dimension the mirror is topologically the same as $`X`$ (in 3 dimensions there is a topology change at the singular fibres) and the correspondence we want is precisely the Fourier-Mukai transform described before. This gives an equivalence of categories $`D^b(X)D^b(X)`$, which is the appropriate even dimensional mirror symmetry, taking 0, 2, 4 branes (points, 2-cycles, and vector bundles) to even dimensional 0, 2, 4 branes (3 dimensional mirror symmetry as we have described should take 0, 2, 4, 6 branes to 3 branes).
We can ask what else we get from the Fourier-Mukai transform. By (2.5) a fibre, or more generally the pushforward of a line bundle on a fibre, gets taken to the corresponding point in the dual fibre:
This was of course the original idea of SYZ: that the moduli of special Lagrangian tori, plus flat line bundles, on $`\widehat{X}`$ should be isomorphic to the moduli of points on $`X`$, i.e. to $`X`$ itself.
On non-hyperkähler manifolds we cannot deduce anything much about either programme without hard analysis, and that is best left to other speakers here. But we can deduce some consequences, and try to prove these. The surface case which could be dealt with rigorously showed one thing: while the mirror $`\widehat{X}`$ was isomorphic to the original $`X`$, the mirror map was certainly *not* the identity: in fact it took points to fibres plus line bundles, rather than points, and induced a non trivial map on $`H^{}(X)`$ on taking Mukai vectors.
In particular we see that *mirror symmetry is not functorial on points* (a phrase I learnt from Paul Seidel); in fact, as Kontsevich envisaged, *(graded) symplectic automorphisms of the mirror $`\widehat{X}`$ should not induce holomorphic automorphisms of $`X`$, but autoequivalences of its derived category $`D^b(X)`$*.
This is what we shall concentrate on.
## 4. Autoequivalences of derived categories and braid groups
### Autoequivalences of derived categories
What are the autoequivalences of $`D^b(X)`$ ? There are the obvious ones given by translation of complexes $`[n]`$, tensoring with line bundles $`L`$, and those pulled back from automorphisms of $`X`$. Bondal and Orlov \[BO\] have shown that this is all of them for $`X`$ smooth and projective with ample canonical bundle or anticanonical bundle. In fact for there to be any more $`X`$ must be partially Calabi-Yau: Orlov \[Or\] has shown that any autoequivalence is set up by an object $`FD^b(X\times \widehat{X})`$ on the product as a Fourier-Mukai transform (2.4), which by (2.8) must then satisfy $`F_p\omega _XF_p`$ for all $`pX`$. And for Calabi-Yau manifolds Kontsevich’s proposal predicts there should be many such Fourier-Mukai transforms; in particular there should be one for every (graded) symplectomorphism of the mirror $`\widehat{X}`$.
For instance Seidel \[S\] has shown that given any Lagrangian sphere $`LS^n`$ in a symplectic manifold we may construct a symplectomorphism – the generalised Dehn twist about $`L`$. This is given as monodromy around a degeneration of the manifold in which the sphere is collapsed to (becomes the vanishing cycle of) an ordinary double point. Alternatively we can glue in the following local model twist on $`T^{}S^n`$. Give $`T^{}S^n`$ its standard symplectic structure and metric. Then $`\mu (\xi )=|\xi |`$, the length of a cotangent vector $`\xi `$, gives a Hamiltonian function which induces a circle action – the flow from $`\xi `$ (considered as a tangent vector using the metric) is the *unit speed* geodesic flow in the direction $`\xi `$ along $`S^n`$, lifted horizontally to $`TS^nT^{}S^n`$ as $`(x(t),\xi (t)=\dot{x}(t)|\xi (0)/\dot{x}(0)|)`$. This flow is clearly discontinuous across the zero section $`S^nT^{}S^n`$ (it is unit speed in opposite directions as we pass through the zero section) but since the geodesics have length $`2\pi `$, the flow through time $`\pi `$ *is* continuous, and gives the antipodal map. So we may define the generalised Dehn twist as the flow of any point $`\xi `$ through an angle varying smoothly from $`0=2\pi `$ as $`|\xi |\mathrm{}`$ to $`\pi `$ at the zero section $`|\xi |=0`$; see \[S\].
The Dehn twist has a canonical lift to the graded symplectomorphism group of $`\widehat{X}`$, and so should be dual to a Fourier-Mukai transform constructed from an element $`_L`$ of $`D^b(X)`$ (not $`D^b(X\times X)`$, notice) dual to the Lagrangian $`L`$. Since Homs should be the same on both sides, we know that $`𝐑\mathrm{Hom}(_L,_L)`$ should be isomorphic to $`HF^{}(L,L)H^{}(S^n)`$. Thus we might expect to be able to find an invertible Fourier-Mukai transform for every *spherical* $`D^b(X)`$, where
###### Definition 4.1.
$`D^b(X)`$ is $`N`$-spherical if $`𝐑\mathrm{Hom}(,)H^{}(S^N;)`$, where $`N=\mathrm{dim}X`$. That is
$$\mathrm{Ext}^i(,)=\{\begin{array}{cc}& i=0,N,\hfill \\ 0& i0,N.\hfill \end{array}$$
Thus for instance a simple (End($`)=.\mathrm{id}`$), rigid ($`\mathrm{Ext}^1(,)=0`$) sheaf $``$ on a smooth Calabi-Yau 3-fold is 3-spherical by Serre duality: $`\mathrm{Ext}^i(,𝒢)\mathrm{Ext}^{3i}(𝒢,)^{}`$.
###### Definition 4.2.
For an $`N`$-spherical $`D^b(X)`$ define the twist $`T_{}𝒢`$ of $`𝒢D^b(X)`$ to be the cone (total complex) of the evaluation map
(4.1)
$$\stackrel{𝐋}{}𝐑\mathrm{Hom}(,𝒢)𝒢.$$
Here we should pick suitable resolutions so the above becomes a genuine map of complexes, then take the total complex of this map. This defines $`T_{}𝒢`$ only up to quasi-isomorphism (cones are not functorial in $`D^b(X)`$); $`T_{}`$ can in fact be made into a functor $`D^b(X)D^b(X)`$ \[ST\].
So in simple cases $`T_{}𝒢`$ will be the kernel or cokernel of a map $`\mathrm{Hom}(,𝒢)𝒢`$. For $`X`$ a smooth projective variety, it is easy to see that $`T_{}`$ is the Fourier-Mukai transform given by the object
(4.2)
$$\{𝐑om(\pi _2^{},\pi _1^{})𝒪_\mathrm{\Delta }\}$$
of $`D^b(X\times X)`$. Here $`\pi _i`$ is projection to the $`i`$th factor $`X`$, the map is restriction to the diagonal $`\mathrm{\Delta }`$ followed by the trace map, and the braces mean we take the cone (total complex) of the map. However, it is more convenient for us to work with general twists in arbitrary derived categories; thus for instance our results apply to non-compact schemes without difficulty.
The categorical equivalent of the partial Calabi-Yau condition $`\omega _X`$ is the existence of a functorial duality (Serre duality for us, or a local form of it in the non-compact case) between $`𝐑\mathrm{Hom}(,𝒢)`$ and $`𝐑\mathrm{Hom}(𝒢,)[N]`$; equivalently the pairing
(4.3)
$$\mathrm{Ext}^i(,𝒢)\mathrm{Ext}^{Ni}(𝒢,)\stackrel{}{}\mathrm{Ext}^N(,)$$
should be perfect, where the last map uses the fact that $``$ is $`N`$-spherical.
###### Definition 4.3.
For $``$ $`N`$-spherical (4.1) and with a duality (4.3), there is a functor $`T_{}^{}:D^b(X)D^b(X)`$ \[ST\] such that the quasi-isomorphism class of $`T_{}^{}𝒢`$ is the total complex we denote
$$\{𝒢\stackrel{𝐋}{}𝐑\mathrm{Hom}(,𝒢)[N]\}$$
given by dualising a map of chain complexes representing the evaluation
$$𝒢\stackrel{𝐋}{}𝐑\mathrm{Hom}(𝒢,).$$
Again this is really a Fourier-Mukai transform with object the derived dual of (4.2) shifted by $`[N]`$.
The last three definitions, and the theorems below, can of course be formulated in the derived category of an arbitrary abelian category, linear over a field $`k`$, having enough injectives, and containing a spherical $``$ with a duality (4.3). We shall confine ourselves to derived categories of coherent sheaves (this does not have enough injectives and one must work with quasi-coherent complexes with coherent cohomology in the usual way; see \[ST\] for full details).
###### Theorem 4.4.
For $``$ spherical with a duality (4.3), $`T_{}`$ and $`T_{}^{}`$ are inverses.
Again we are being sketchy here. The precise statement \[ST\] takes place in the derived category of a $`k`$-linear abelian category, and “inverses” means that $`T_{}T_{}^{}`$ and $`T_{}^{}T_{}`$ are both naturally isomorphic to the identity functor.
Using Serre duality on a compact scheme, or a local form of it one can prove for $``$ compactly supported on a non-compact scheme, we get an invertible Fourier-Mukai transform in each of the following examples.
* Any simple, rigid sheaf $``$ on a Calabi-Yau 3-fold $`X`$. In particular the structure sheaf $`𝒪_X`$ gives a canonical transform, called the reflection functor by Mukai, which should be mirror dual to the Dehn twist about the (conjectural) special Lagrangian $`S^3`$ zero section of the SYZ fibration of $`\widehat{X}`$.
* Holomorphic $`2`$-spheres $`C`$ in a complex surface give Fourier-Mukai transforms, taking $`=𝒪_C`$. In particular we get Fourier-Mukai transforms for general-type surfaces containing such spheres, which might be surprising given the results of \[BO\]. These surfaces are locally Calabi-Yau along the spheres. The induced action on cohomology is the Picard-Lefschetz reflection in the corresponding $`2`$-vector.
* Spheres with normal bundle $`𝒪_^1(1)𝒪_^1(1)`$ in 3-folds also give transforms in the same way.
* Surfaces $`S`$ in 3-folds satisfying $`h^{0,1}(S)=0=h^{0,2}(S)`$, with the local Calabi-Yau condition $`\nu _S\omega _S`$.
Seidel proved in \[S\] that the Dehn twists along $`A_n`$-chains of Lagrangian spheres satisfy the braid relations. That is, if we have a chain of such spheres with consecutive pairwise intersections one transverse point, then we get a homomorphism from the braid group on $`(n+1)`$ strands $`B_{n+1}`$ into the symplectomorphism group of the ambient manifold. Moreover he showed the smoothing of an $`A_n`$ singularity on a complex surface (such as the smoothing of the standard $`SL(2,)`$ quotient singularity $`^{\mathrm{\hspace{0.17em}2}}/_n`$) contains such a chain of Lagrangian spheres.
For two Lagrangians $`L_i`$ intersecting transversely in a single point we have $`HF^{}(L_1,L_2;)`$, so we define
###### Definition 4.5.
A sequence of spherical objects $`_iD^b(X),i=1,\mathrm{},n`$ form an $`A_n`$-chain if they satisfy
$$\underset{k}{}\mathrm{dim}\mathrm{Ext}^k(_i,_j)=\{\begin{array}{cc}0& |ij|>1,\hfill \\ 1& |ij|=1.\hfill \end{array}$$
Thus there are no Homs between distinct $`F_i`$ s unless they are consecutive, in which case there is a unique map (up to scale) in some (arbitrary) degree.
###### Theorem 4.6.
Given an $`A_n`$-chain of spherical objects $`_iD^b(X)`$ with duality (4.3) there are the following natural isomorphisms between the corresponding functors $`T_i=T__i`$
$$\begin{array}{cc}T_iT_i^{}\mathrm{id}T_i^{}T_i\hfill & \\ T_iT_jT_iT_jT_iT_j\hfill & |ij|=1,\hfill \\ T_iT_jT_jT_i\hfill & |ij|>1.\hfill \end{array}$$
Thus they define a weak braid group action on $`D^b(X)`$. (We have not checked if the natural isomorphisms above satisfy the coherence relations of \[De\] to define a genuine $`B_{n+1}`$ action on the category $`D^b(X)`$.)
Thus we get braid group actions on derived categories of coherent sheaves in the following cases.
* $`A_n`$-chains of $`2`$-spheres $`C_i`$ (i.e. a sequence of $`2`$-spheres with consecutive pairwise intersections a reduced point) in quasi-projective surfaces give actions of $`B_{n+1}`$ with $`_i=𝒪_{C_i}`$. For instance the ALE spaces that are the resolutions of the $`SL(2,)`$ quotient singularities $`^{\mathrm{\hspace{0.17em}2}}/_n`$.
* If, as in the previous example, we have an $`A_n`$-chain of $`2`$-spheres $`C_i`$, we can twist instead by the corresponding line bundles $`L_i=𝒪(C_i)`$. Simple exact sequences show that these also form an $`A_n`$-chain of spherical objects in $`D^b(X)`$ for $`X`$ a $`K3`$ surface.
* Chains of surfaces $`S_i`$ in 3-folds, each with $`\nu _S\omega _S`$ and $`h^{0,1}(S)=0=h^{0,2}(S)`$, give braid group actions (for $`_i=𝒪_{S_i}`$) if $`S_iS_j=\mathrm{}`$ for $`|ij|>1`$ and $`S_iS_{i+1}`$ is a $`^1`$-fibre of one surface and a $`2`$-sphere in the other. Again such configurations arise in crepant resolutions of 3-fold singularities.
* Taking $`_i`$ to be line bundles on an elliptic curve, with the degrees of $`_i`$ and $`_j`$ differing by $`ij`$, we recover the original Fourier-Mukai transforms of \[Mu\], and from any two consecutive such line bundles we get an action of $`B_3`$, a central extension of $`SL(2,)`$. The central generator $`(T_1T_2)^6`$ acts as the translation $`[\mathrm{\hspace{0.17em}2}]`$ and the action is easily seen to be the $`SL(2,)`$ action of Mukai. There are mirror, symplectic analogues of these relations for Dehn twists on tori in \[ST\].
### Singularities
Notice the chains of $`2`$-spheres in surfaces are different in the holomorphic and symplectic cases – the former appear in resolutions of singularities, the latter in smoothings of the same singularities. Are the corresponding braid group actions, on $`D^b(X)`$ and $`D^b(\mathrm{Fuk}(X))`$ respectively, mirror dual ?
In the compact case we know what the mirror to the Lagrangian spheres should be in the presence of an SYZ fibration – the Fourier-Mukai transforms of their structure sheaves. If the spheres lie in the elliptic fibres then their transforms are themselves, and the braid group action on the symplectic side should be dual to the first example listed above. If however the spheres are sections of the fibration then their transforms are given by the line bundles $`𝒪(C_i)`$ (twisted by $`𝒪(s)`$, where $`s`$ is the image of the zero section, but this need not concern us), and the correct dual is the second example in the above list.
Either way mirror symmetry seems rather local in these cases, and there are other cases where the mirror dual of the (symplectic) smoothing of one singularity is the (holomorphic) resolution of another. In fact this is the general proposal of \[Mo\]: Morrison suggests that moving towards the discriminant locus in the complex structure moduli space of $`X`$ (i.e. degenerating $`X`$ to a singular Calabi-Yau) should be mirror dual to moving to a “boundary wall” of the complexified Kähler cone of $`\widehat{X}`$ (the annihilator of a face of the Mori cone), thus inducing an extremal contraction of $`\widehat{X}`$. Resolving the singularities of $`X`$ should then be mirror dual to smoothing the contracted $`\widehat{X}`$. In particular, in some generic situations, the smoothing of ordinary double points (with their Lagrangian $`S^3`$ vanishing cycles and corresponding Dehn twists on $`\widehat{X}`$) should be mirror to small resolutions of other ordinary double points (with exceptional $`^1`$ loci giving corresponding twists $`T_{𝒪_^1}`$ on $`D^b(X)`$).
Our results indicate that it might be reasonable to expect certain additional properties of singularities that are dual in this way. Namely, singularities whose smoothings have a Dynkin diagrams of Lagrangian $`S^3`$ vanishing cycles should be dual to singularities whose resolution has a similar diagram of (spherical structure sheaves of) irreducible components of its exceptional set (these form the nodes of the diagram, edges are provided by $`𝐑\mathrm{Hom}`$s).
For instance consider the smoothing of the 3-fold singularity
$$x^2+y^2+z^2+t^{2n}=0,$$
which contains an $`A_{2n1}`$-chain of Lagrangian $`S^3`$ vanishing cycles of the singularity giving a braid group of (graded) symplectomorphisms. Peturbing this into $`n`$ ordinary double points $`x^2+y^2+z^2+t^2=0`$ (which we can do on the symplectic side, we are only varying complex structure), analysing the effect on homology and its mirror, and using Morrison’s proposal, we are led in \[ST\] to ask whether the mirror should be given by the following geometry, which we know leads to a braid group action.
###### Proposition 4.7.
Suppose we have a chain of Fano surfaces $`\{S_{2k}\}_{k=1}^{n1}`$ in a smooth 3-fold $`X`$, with the local Calabi-Yau condition that their normal bundles are isomorphic to their canonical bundles. Suppose also that each $`S_{2k}`$ contains two disjoint $`(1)`$-$`^1`$s, $`C_{2k1}`$ and $`C_{2k+1}`$, in which it intersects $`S_{2k2}`$ and $`S_{2k+2}`$ respectively (and there are no more intersections, so only consecutive surfaces intersect). Then the sheaves
$$_{2k}=𝒪_{S_{2k}}\mathrm{and}_{2k+1}=𝒪_{C_{2k+1}}$$
form an $`A_{2n1}`$-chain, and so define an action of $`B_{2n}`$ on $`D^b(X)`$.
In particular if we take the surfaces $`S_{2k}`$ to be $`^2`$s blown up in two points (giving the two exceptional curves $`C_{2k1}`$ and $`C_{2k+1}`$ which we think of as the mirrors of the $`(2k1)`$th and $`(2k+1)`$th vanishing cycles in the smoothing of $`x^2+y^2+z^2+t^{2n}=0`$ according to Morrison’s proposal) then there is an extra $`(1)`$-curve $`C_{2k}`$ in $`S_{2k}`$ – the proper transform of the line joining the two blow-up points – which we can think of as the mirror of the $`2k`$th vanishing cycle. Such a configuration is easily shown \[ST\] to arise in smooth toric Calabi-Yau manifolds as the crepant resolution of a nasty singularity that we would like to think of as the dual of $`x^2+y^2+z^2+t^{2n}=0`$.
Another relevant example is Arnold’s strange duality (see for instance \[Pi\]), which is encompassed by mirror symmetry for $`K3`$ surfaces according to the work of a number of people (Aspinwall and Morrison, Kobayashi, Dolgachev, Ebeling, etc.). To every isolated surface singularity on Arnold’s list, described by three numbers $`b_1,b_2,b_3`$, there is a natural $`K3`$ compactification $`S`$ of the singularity containing at infinity a chain of $`2`$-spheres with intersection configuration given by the following Dynkin diagram $`T(b_1,b_2,b_3)`$:
Here the circles represent $`2`$-spheres, and edges give intersections of the spheres of intersection number 1, and the central $`2`$-sphere is counted in each $`b_i`$. The corresponding intersection matrix is not negative definite for the numbers $`b_i`$ in the list, so the spheres cannot be completely contracted, though they can be contracted to a smooth $`^1`$ with three points on it that are singular points of the surface (this surface is Pinkham’s original compactification of the affine surface singularity).
Let $`\{c_i\}_{i=1}^3`$ be the numbers dual to the $`b_i`$ s in strange duality (i.e. these are the $`b_i`$ s associated to the dual singularity). Then the intersection matrix of the smoothing of the original singularity has vanishing cycles given by $`T(c_i)H`$ (where $`H`$ is the hyperbolic $`(\begin{array}{cc}0& 1\\ 1& 0\end{array})`$). Together the $`2`$-spheres plus the vanishing cycles give all of the homology of the smoothed $`K3`$ surface: $`H_2(K3)T(b_i)T(c_i)H`$. The duality swaps the $`b_i`$ s and $`c_i`$ s, taking the homology $`T(c_i)H`$ generated by the vanishing cycles to the holomorphic homology $`T(c_i)H`$ in the mirror $`K3`$ (given by the resolution cycles and the hyperbolic $`H_0H_4`$).
Thus we can think of mirror symmetry as replacing the smoothing of one singularity by the resolution of another (though it is not an isolated singularity, it is the $`^1`$ with 3 surface singularities on it). Here we are not thinking of the duality as merely swapping Arnold’s singularities, but as a rather more global phenomenon on $`K3`$ surfaces, which therefore has to include the chains of $`2`$-spheres at infinity.
There is a generalised braid group action on the derived category of the $`K3`$ surface given by the $`T(b_1,b_2,b_3)`$ $`2`$-sphere configuration. Dually, there is probably the same group of Dehn twists around the Lagrangian vanishing cycles of the mirror $`K3`$, though it seems not to be known whether they can be found in the *geometric* intersection configuration $`T(b_1,b_2,b_3)`$ (they may have many more geometric intersections than their topological intersections suggest).
## 5. Mutations
The formula (4.1) for our twist is familiar to algebraists, in tilting theory, and those who work on exceptional bundles on Fano manifolds – see e.g. \[Ru\]. There the twists are called mutations, and act on certain modules over algebras (similar to those described in the next section) built from the bundles (they cannot give equivalences of derived categories by the result of \[BO\]). The bundles $``$ that one twists are also those with minimal $`\mathrm{Ext}^{}(,)`$; in the case of Fano manifolds this means $``$ is simple and has no higher Exts at all, and is called *exceptional*.
There are braid group actions of such twists on exceptional collections of bundles, but the relation with our work is far from clear, and it is possible there is none. Here we shall simply note a relationship between exceptional objects on Fano manifolds and spherical objects on Calabi-Yaus, motivated by two examples of \[Ku\].
###### Definition 5.1.
We say that a map $`f:XY`$ from a Calabi-Yau $`N`$-fold $`X`$ to a smooth projective Fano $`M`$-fold $`Y`$, of codimension $`c=NM`$ (which may be of any sign or zero), is *simple* if $`𝐑f_{}𝒪_X`$ is made up from $`𝒪_Y`$ and $`\omega _Y`$ in the sense that
* $`c>0`$ $`R^if_{}𝒪_X=\{\begin{array}{cc}𝒪_Y& i=0\hfill \\ 0& i0,c\hfill \\ \omega _Y& i=c\hfill \end{array}`$
* $`c=0`$ $`R^if_{}𝒪_X=\{\begin{array}{cc}𝒪_Y\omega _Y& i=0\hfill \\ 0& i0\hfill \end{array}`$
* $`c<0`$ There is an exact triangle $`\omega _Y[c]𝐑f_{}𝒪_X𝒪_Y`$.
Simple examples of maps $`f`$ are often of this type. For instance, for Calabi-Yaus fibred over a Fano base with generic fibre $`F`$ such that $`h^{0,i}(F)=0,0<i<c`$, relative Serre duality shows that the projection is simple in this sense.
Examples with $`c=0`$ are given by Calabi-Yaus double covering Fanos, branched over a double anticanonical divisor, while for $`c=1`$ we have a Calabi-Yau anticanonical divisor in a Fano manifold.
###### Theorem 5.2.
Suppose $`f:XY`$ is a simple map, as defined above, and $`D^b(Y)`$ is exceptional $`(\mathrm{Ext}^i(,)=`$ for $`i=0`$, and $`0`$ for $`i0)`$. Then $`𝐋f^{}D^b(X)`$ is spherical.
In the other direction, i.e. maps from a Fano to a Calabi-Yau, something can only be said in the case of a Fano divisor in a (locally) Calabi-Yau manifold. Namely,
###### Theorem 5.3.
Suppose that $`\iota :YX`$ is a smooth Fano divisor with normal bundle $`\nu _Y=\omega _Y`$ in a quasi-projective scheme $`X`$. If $`D^b(Y)`$ is exceptional then $`\iota _{}D^b(X)`$ is spherical.
This takes care of most of the examples of spherical sheaves given until now (by taking the exceptional sheaf to be $`𝒪_Y`$), and provides many more by pushing forward the exceptional collections of \[Ru\] to (local) Calabi-Yau manifolds.
## 6. Fidelity
Finally we briefly mention what is in many ways the main result of \[KS, ST\], namely that in dimension $`N2`$, the $`B_{n+1}`$ actions given by $`A_n`$-chains of spherical objects are *faithful*.
To do this it is clearly enough to show the induced $`B_{n+1}`$ action on the differential graded modules $`_i𝐑\mathrm{Hom}(_i,𝒢)`$, in the derived category of the differential graded algebra $`_{ij}𝐑\mathrm{Hom}(_i,_j)`$, is faithful. (Here the algebra and module structures are the obvious ones; replacing each $`_i`$ by a finite resolution of injectives, which one can prove is possible, it is just composition of morphisms.)
In fact a difficult result of \[KS\] is that the induced action on homology, i.e. the action on the graded modules $`_{ik}\mathrm{Ext}^k(_i,𝒢)`$ in the derived category of the graded algebra $`_{ijk}\mathrm{Ext}^k(_i,_j)`$, is faithful. (Since the $`_i`$ s form an $`A_n`$-chain these modules and algebras take a standard form, and the braid group action is the one considered in \[KS\].) However this is not enough to prove faithfulness of the action at the level of differential graded modules. The following result, though, is sufficient to provide a proof.
###### Theorem 6.1.
For $`\{_i\}_{i=1}^n`$ an $`A_n`$-chain of $`N`$-spherical objects, $`N2`$, the graded algebra $`A=_{ik}\mathrm{Ext}^k(_i,_i)`$ is *intrinsically formal*. That is, any differential graded algebra with $`A`$ as its cohomology is quasi-isomorphic to $`A`$.
This is proved by a lot of non-commutative obstruction theory that I will not go into here. That such machinery is really necessary is seen in the following example of non-faithfulness in dimension $`N=1`$.
Let $`T`$ be an elliptic curve, $`L`$ be a non-trivial degree zero line bundle, and denote by $`𝒪_p`$ the structure sheaf of a point $`pT`$. Notice that $`L`$ induces an automorphism $`\varphi _L`$ of $`T`$: identifying $`T`$ with its degree one line bundles, $`\varphi _L`$ is given by tensoring with $`L`$.
The sheaves $`𝒪=𝒪_T,𝒪_p`$ and $`L`$ form an $`A_3`$-chain of $`1`$-spherical sheaves.
###### Theorem 6.2.
The action of $`T_L^{}T_𝒪`$ is the action induced by the automorphism $`\varphi _L`$ by pullback. In particular if $`L`$ has order two $`(L^2=𝒪)`$ then $`(T_L^1T_𝒪)^2\mathrm{id}.`$ |
warning/0001/astro-ph0001493.html | ar5iv | text | # Analytic model for galaxy and dark matter clustering
## I Introduction
Correlations in dark matter contain a wealth of information about cosmological parameters. Their power spectrum is sensitive to parameters such as matter density, Hubble constant, primordial power spectrum slope and amplitude, massive neutrinos, baryon density etc. Determining the linear power spectrum of dark matter is one of the main goals of modern cosmology. There are several complications that prevent us at present from reaching this goal. First, on small scales the linear power spectrum is modified by nonlinear evolution which enhances its amplitude over the linear spectrum. It is important to understand this process, so that one can predict the relation between the two. This is necessary both to reconstruct the linear spectrum from a measured nonlinear one and to verify whether there are other mechanisms besides gravity that modify the clustering of dark matter on small scales. Examples of such are baryonic feedback effects on dark matter or nongravitational interactions between dark matter particles . Second, it is difficult to observe correlations in dark matter directly. Direct tracers such as peculiar velocity flows or weak lensing still suffer from low statistics and poorly understood systematics. Instead it is much easier to observe correlations between galaxies or correlations between galaxies and dark matter . While these are related to the dark matter correlations, the relation may not be simple. The goal of this paper is to address both issues with a model that is simple enough to allow analytic calculations without the use of N-body simulations, yet sufficiently accurate to be useful for predicting galaxy and dark matter power spectrum.
Our approach to dark matter clustering is based on the Press & Schechter model . In this picture at any given time all the matter in the universe is divided into virialized halos. These halos are correlated and have some internal density profile, which can be a function of halo mass. By specifying the halo mass function, their clustering strength and their halo profile we can determine the dark matter correlation function. The formalism for correlations inside halos has been developed by and applied to power law halos . We generalize this approach by including the correlations between halos and by using more realistic non-power law halo profiles whose shape may depends on the halo mass . We show in this paper that such a generalized model can provide very good agreement with results of numerical simulations over a wide range of scales .
The central question in extracting dark matter power spectrum from that of the galaxies is how well galaxies trace dark matter, the issue of bias. This has been addressed theoretically both with hydrodynamic and semi-analytic methods . The fact that the galaxy correlation function is a power law over several decades in scale, while power spectra in CDM models do not show such behaviour, already indicates that the bias is scale dependent. Moreover, galaxies come in different types and observational data show that they can be biased relative to one another . In our modelling of galaxy correlations we introduce two new functions, the mean number and the mean pair weighted number of galaxies inside the halo as a function of the halo mass. The importance of these has recently been emphasized in the context of pairwise velocity measurements and galaxy clustering . These play a key role in understanding the relation between galaxy and dark matter clustering. We explore the predictions for different choices of these relations and compare them to the results of semi-analytic models.
Galaxy-dark matter correlations can provide additional information on the clustering of galaxies and dark matter and the relation between them. Such correlations have been observed through gravitational lensing effects, for example using galaxy-galaxy lensing or correlations between foreground and background populations . Such measurements are often interpreted either in terms of an averaged density profile of a halo or in terms of a constant bias model . We discuss the applicability of these models and how they can be generalized to take into account effects such as broad range of halo masses and multiple galaxies inside halos.
## II Dark matter power spectrum
The halo model for power spectrum assumes all the matter is in a form of isolated halos with a well defined mass $`M`$ and halo profile $`\rho (r,M)`$. The halo profile is defined to be an average over all halos of a given mass and does not necessarily assume all halos have the same profile. The mass is determined by the total mass within the virial radius $`r_v`$, defined to be the radius where the mean density within it is $`\delta _{\mathrm{vir}}`$ times the mean density of the universe. Throughout the paper we will use $`\mathrm{\Lambda }CDM`$ model with $`\mathrm{\Omega }_m=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, normalized to $`\sigma _8=0.9`$ today. For this model $`\delta _{\mathrm{vir}}340`$, although we will also use $`\delta _{\mathrm{vir}}200`$ (the value for Einstein-de Sitter universe) for consistency with the results of some of the N-body simulations. The halo profile is spherically averaged and assumed to depend only on the mass of the halo. We will model the halo density profile in the form
$$\rho (r)=\frac{\rho _s}{(r/r_s)^\alpha (1+r/r_s)^{3+\alpha }}.$$
(1)
This model assumes that the profile shape is universal in units of scale radius $`r_s`$, while its characteristic density $`\rho _s`$ at $`r_s`$ or concentration $`c=r_v/r_s`$ may depend on the halo mass. The halo profile is assumed to go as $`r^3`$ in the outer parts and as $`r^\alpha `$ in the inner parts, with the transition between the two at $`r_s`$. The outer slope is fixed by the results of N-body simulations which generally agree in this regime. An example of such a profile is $`\alpha =1`$ . Other models have however been proposed with $`\alpha =1.5`$ or even $`\alpha >1`$ . In principle $`\alpha `$ could also be a function of mass scale and may steepen towards smaller mass halos with $`\alpha 1`$ for cluster halos and $`\alpha 1.5`$ for galactic halos . Similarly, concentration $`c`$ may depend on the mass and different authors find somewhat different dependence . We will explore how variations in the profile and concentration affect the power spectrum. Instead of $`r_s`$ we use the concentration parameter $`c=r_v/r_s`$ as a free parameter. Note that $`r_v`$ is related to $`M`$ via $`M=4\pi /3r_v^3\delta _{\mathrm{vir}}\overline{\rho }`$. Similarly we can eliminate $`\rho _s`$ and describe the halo only in terms of its virial mass $`M`$ and concentration $`c`$, because the integral over the halo density profile (equation 1) must equal the halo mass.
For a complete description we need the halo mass function $`dn/dM`$, describing the number density of halos as a function of mass. It can be written as
$$\frac{dn}{dM}dM=\frac{\overline{\rho }}{M}f(\nu )d\nu ,$$
(2)
where $`\overline{\rho }`$ is the mean matter density of the universe. We introduced function $`f(\nu )`$, which can be expressed in units in which it has a universal form independent of the power spectrum or redshift if written as a function of peak height
$$\nu =[\delta _c(z)/\sigma (M)]^2.$$
(3)
Here $`\delta _c`$ is the value of a spherical overdensity at which it collapses at $`z`$ ($`\delta _c=1.68`$ for Einstein-de Sitter model) and $`\sigma (M)`$ is the rms fluctuation in spheres that contain on average mass $`M`$ at initial time, extrapolated using linear theory to $`z`$. The form proposed by Press & Schechter (PS) is $`\nu f(\nu )=(\nu /2\pi )^{1/2}e^{\nu /2}`$. This has been shown to overpredict the halo abundance by a factor of 2 at intermediate masses below nonlinear mass scale $`M_{}`$ . A modified version of this form that fits better the N-body simulations is given by Sheth & Tormen (ST)
$$\nu f(\nu )=A(1+\nu ^p)\nu ^{1/2}e^{\nu ^{}/2},$$
(4)
where $`\nu ^{}=a\nu `$ with $`a=0,707`$ and $`p=0.3`$ as the best fitted values, which gives $`\nu f(\nu )\nu ^{0.2}`$ for small $`\nu `$. PS expression corresponds to $`a=1`$, $`p=0`$ giving $`\nu f(\nu )\nu ^{0.5}`$ for small $`\nu `$. The constant $`A`$ is determined by mass conservation, requiring that the integral over the mass function times the mass gives the mean density
$$\frac{1}{\overline{\rho }}\frac{dn}{dM}M𝑑M=f(\nu )𝑑\nu =1.$$
(5)
Note that we can still apply this equation even if some dark matter is not bound to any halo. In this case the mass function has a nonvanishing contribution in the limit $`M0`$.
The correlation function consists of two terms. On large scales the halos are correlated with each other. We assume the halo-halo correlation function follows the linear correlation function. Its amplitude depends on the bias for each halo. Halos more massive than the nonlinear mass scale $`M_{}`$ are more strongly clustered than the matter, while those with masses below $`M_{}`$ are less strongly clustered than the matter. A simple halo biasing scheme has been given by and generalized to the ST mass function by
$$b(\nu )=1+\frac{\nu 1}{\delta _c}+\frac{2p}{\delta _c(1+\nu ^p)}.$$
(6)
Since halos are not pointlike we need to convolve the halo-halo correlation function with the halo profiles of both halos to obtain the dark matter correlation function. The expressions simplify significantly in Fourier space, where convolution becomes a multiplication with the Fourier transform of the halo profile
$$\stackrel{~}{\rho }(k,M)=4\pi r^2𝑑r\rho (r,M)\frac{\mathrm{sin}(kr)}{kr}.$$
(7)
Note that this is normalized so that $`\stackrel{~}{\rho }(0,M)=M`$. It is convenient to renormalize it to unity by introducing a new variable $`y(k,M)=\stackrel{~}{\rho }(k,M)/M`$, so that $`y(0,M)=1`$ and $`y(k>0,M)<1`$. The mass of the halo rapidly increases as $`r^{3\alpha }`$ up to $`r=r_s`$, but increases only logarithmically between $`r_s`$ and $`r_v`$ if the outer profile is $`\rho (r)r^3`$. The dominant contribution to the mass therefore comes from radii around $`r_s`$. For $`kr_s1`$ we have $`y1`$. At $`kr_s1`$ there is a transition and $`y`$ begins to decrease with $`k`$, so that for $`kr_s1`$ we have $`y(k,M)(kr_s)^{(3+\alpha )}`$.
Because the expressions simplify significantly in Fourier space we will in the following only describe the power spectrum analysis. The halo-halo term is given by the integral over their mass function with the appropriate bias and the halo profile transform,
$$P_{\mathrm{dm}}^{hh}(k)=P_{\mathrm{lin}}(k)\left[f(\nu )𝑑\nu b(\nu )y[k,M(\nu )]\right]^2,$$
(8)
where $`P_{\mathrm{lin}}(k)`$ is the linear power spectrum and $`M`$ is related to $`\nu `$ via equation 3 using the relation between $`\sigma ^2(M)=4\pi P_{\mathrm{lin}}(k)W_R(k)k^2𝑑k`$ and $`M=4\pi R^3\overline{\rho }/3`$, where $`W_R(k)`$ is the Fourier transform of the top hat window with radius $`R`$. This gives $`M\nu ^{3/(n+3)}`$, where $`n`$ is the slope of the linear power spectrum at scale $`kR^1`$. We can also define the nonlinear mass scale $`M_{}`$ where $`\nu =1`$. Note that on galaxy and smaller scales $`n<2`$ and the relation between $`M`$ and $`\nu `$ is very steep, $`M\nu ^\gamma `$ with $`\gamma >3`$. The requirement that on large scales ($`k0`$, $`y1`$) the power spectrum reduces to the linear power spectrum imposes a nontrivial constraint on the bias distribution,
$$f(\nu )𝑑\nu b(\nu )=1.$$
(9)
This implies that if halos are biased ($`b>1`$) for $`M>M_{}`$ at least some of the halos with $`M<M_{}`$ must be antibiased ($`b<1`$) to satisfy this constraint. Most of the bias descriptions in the literature satisfy this constraint to within a few percent. The halo-halo term follows the linear power spectrum on large scales and drops below it on scales where finite extent of the halos become important (ie where $`y(k,M)<1`$). This term is shown in figure 1 and as expected is dominant on large scales.
In addition to the halo-halo correlation term there are also correlations between dark matter particles within the same halos. These are expected to dominate on small scales. We denote this the Poisson term, which is given by
$$P_{\mathrm{dm}}^P(k)=\frac{1}{(2\pi )^3}f(\nu )𝑑\nu \frac{M(\nu )}{\overline{\rho }}|y[k,M(\nu )]|^2,$$
(10)
The main difference between this term and the halo-halo term in equation 8 is that we have an additional mass weighting $`M/\overline{\rho }`$. This makes the dominant contribution to this term to come from the higher mass halos relative to the halo-halo term. On large scales ($`k0`$, $`y1`$) the Poisson term is independent of $`k`$ and behaves as white noise. It increases with $`k`$ more rapidly than the halo-halo term, which scales as the linear power spectrum (figure 1). The Poisson term declines below the white noise on small scales where the effects of the halo profile become important.
The total power spectrum is the sum of the two contributions,
$$P_{\mathrm{dm}}(k)=P_{\mathrm{dm}}^{hh}(k)+P_{\mathrm{dm}}^P(k).$$
(11)
To complete the calculation we need to model the dependence of $`c`$ on $`M`$. We will parametrize it as
$$c=c_0(M/M_{})^\beta .$$
(12)
Typical values for $`c_0`$ are around 10 at the nonlinear mass scale for $`\alpha =1`$ profile and about a third lower for $`\alpha =1.5`$ profile . Numerical studies also show that the concentration decreases slowly with the halo mass, making $`\beta `$ negative.
Figure 1 shows the individual contributions and the sum in comparison to the linear power spectrum and the nonlinear prediction from (PD). In top of the figure we used $`\alpha =1.5`$ and $`c(M)=6(M/M_{})^{0.15}`$. The latter fits the concentration mass dependence given in . Note that for consistency with we use $`\delta _{\mathrm{vir}}=200`$ in this case as opposed to $`\delta _{\mathrm{vir}}=340`$. In bottom of the figure we used the ST mass function and $`\alpha =1`$ with $`c(M)=10(M/M_{})^{0.2}`$, which is somewhat steeper than numerical studies predict as discussed below. The agreement in both cases is quite remarkable given the simple nature of the model. It correctly predicts the transition between the linear and nonlinear power spectrum, as well as reproduces well the slope at higher values of $`k`$. This shows that given a suitable choice of $`c(M)`$ both models can reproduce the nonlinear power spectrum. Conversely, the slope of the power spectrum at high $`k`$ is not directly determined by the inner slope of dark matter profiles, at least if the inner profiles are shallower than $`\alpha =1.5`$.
In the case of $`\alpha =1`$ profile the best fitted value for $`c_0=10`$ agrees well with , while $`\beta 0.2`$ is somewhat lower than $`\beta =0.07`$ and $`\beta =0.13`$ . If one adopts such shallow dependence of $`c(M)`$ with $`\beta 0.1`$ then for $`k>10h`$Mpc<sup>-1</sup> the predictions of the model are systematically below the PD model. Before concluding that this is caused by the galactic halos not being sufficiently compact we must investigate the possibility that the mass function is underestimated at small masses. Replacing ST with PS does not significantly affect the results. However, both PS and ST assume that each mass element belongs to only one halo, counting only the isolated halos. This is certainly a valid description on large scales, where the total halo mass determines the white noise amplitude of the power spectrum. On small scales the clumpiness caused by subhalos within the halos may become important. Recent numerical simulations have in fact shown that most of the small halos that merge into larger ones are not immediately destroyed, but stay around for some time until they are finally merged on the dynamical friction time scale . In such a case a given mass particle can be part of more than one halo at any given time. Because on very small scales the correlation function is dominated by the small halos it may make a difference whether the mass is smoothly distributed within the halos or some fraction of it is in the subhalos. However, the contribution to the total mass of the halo coming from the subhalos is below 10% . Recently, the mass function for subhalos from high resolution simulations was determined and it was shown that it is an order of magnitude below the one for isolated halos . One may conclude therefore that subhalos do not affect the mass function significantly and cannot resurrect $`\beta >0.13`$, $`\alpha =1`$ model.
Steepening the halo profile or changing concentration can both increase the power spectrum to agree better with N-body simulations (figure 1). This is because to increase the power on small scales one has to increase the amount of mass contained within a given radius. This can be achieved either by making the inner profile steeper than $`\alpha =1`$ or making the concentration parameter larger towards the smaller mass halos. The change in the slope would support the results in , where the universal profile has the inner slope close to $`\alpha 1.5`$, or in , where the profile is not universal and steepens as the halo mass is decreased, so that the inner slope changes from $`\alpha 1`$ on the cluster scales to $`\alpha 1.5`$ on the galactic scales. If the inner slope of the halo profile is $`\alpha 1`$ the concentration has a stronger mass dependence than in , although the discrepancy is not large. As shown in figure 1 both models can fit the nonlinear power spectrum on small scales remarkably well.
Further insight into the relation between the halos and the dark matter power spectrum can be obtained by investigating the contribution to the power spectrum from different mass intervals. This is shown in figure 2 for the Poisson term, using the two models from figure 1. On large scales the Poisson term is dominated by very massive clusters with $`M>10^{14}h^1M_{}`$. These halos dominate the nonlinear clustering on scales around and below $`k<1`$hMpc<sup>-1</sup>. On smaller scales the contribution from large clusters is suppressed because $`y(k,M)`$ begins to decrease from unity at $`kr_c^1c(M)M^{1/3}M^{0.5}`$. This occurs at lower $`k`$ for the higher mass halos. As a result around $`k10h`$Mpc<sup>-1</sup> the halos with $`10^{14}h^1M_{}>M>10^{13}h^1M_{}`$ dominate, while around $`k100h`$Mpc<sup>-1</sup> the halos with $`10^{13}h^1M_{}>M>10^{12}h^1M_{}`$ dominate. Note again that the inner slope plays a subdominant role in determining the amplitude of the power spectrum. Even if for a steeper slope the power spectrum from a given mass interval is decreasing less rapidly (for example for $`\alpha =1.5`$ it is asymptotically flat as opposed to decreasing as $`k^1`$ for $`\alpha =1`$), when this becomes important the smaller mass halos have already taken over as a dominant contribution to the power spectrum. The nonlinear power spectrum therefore does not reflect the inner slope of the halo profile, but rather the halo mass function and the radius at which the mass enclosed within this radius begins to deviate significantly from the total halo mass. In both models the halos with $`M>10^{11}h^1M_{}`$ dominate the power spectrum for $`k<100h`$Mpc<sup>-1</sup>. Any modifications in the linear power spectrum on mass scales below $`M10^{11}h^1M_{}`$ would therefore show up in the dark matter correlation function only on kiloparsec scales and below.
It is interesting to explore in more detail the quasi-linear regime, where $`P^P(k)\mathrm{const}`$. This approximation is valid up to $`4\pi k^3P(k)10`$ or $`k1h`$Mpc<sup>-1</sup>. On scales larger than these the power spectrum can be approximated as a sum of a linear power spectrum and a constant term, whose amplitude is given as an integral over the mass function (equation 10 with $`y=1`$). From figure 2 one can see that the amplitude of this integral is dominated by the massive halos, $`M>10^{14}h^1M_{}`$. It is important to emphasize that this amplitude depends only on the integral over the power spectrum and not on the details of the power spectrum itself. Even if there are sharp features in the linear power spectrum, such as for example baryonic wiggles , these would not show up as features in the quasi-linear power spectrum. Instead, they would be integrated over into a single number, corresponding to the mass weighted integral over the mass function (equation 10). This argument is in agreement with the results of N-body simulations which indeed show that any baryonic features are erased in the nonlinear regime. This suggests that while the PD model breaks down for such spectra, our model could also be applied in such a case. This also applies to the spectra with truncated power on small scales . We plan to investigate this further in the future.
## III Galaxy power spectrum
We now apply the above developed model to the galaxies. We assume all the galaxies form in halos, which is a reasonable assumption given that only very dense enviroments which have undergone nonlinear collapse allow the gas to cool and to form stars. The key new parameters we introduce are the mean number of galaxies per halo as a function of halo mass, $`N(M)`$, and the mean pair weighted number of galaxies per halo, $`N(N1)^{1/2}(M)`$. Just as in the case of dark matter these functions are well defined even if the assumption that the statistical properties of galaxy population depend only on the halo mass and not on its enviroment is not satisfied , as long as the averaging is performed over all possible enviroments. The resulting power spectrum on small scales where the Poisson term dominates is independent of this assumption. On large scales where correlations between the halos are important violation of this assumption may lead to a change in the strength of the halo-halo term.
We furthermore assume that there each halo has a galaxy at its center, while the rest of the galaxies in the halos are distributed in the same way as the dark matter, so $`y(k,M)`$ remains unchanged. This is only the simplest model and one could easily generalize it to profiles that differ from the dark matter. Any such complications are important on small scales, while on large scales ($`k<1h`$Mpc<sup>-1</sup>) all that is relevant is the total number of galaxies inside the halo. The normalization equation 9 becomes
$$\frac{N}{M}f(\nu )𝑑\nu =\frac{\overline{n}}{\overline{\rho }},$$
(13)
where $`\overline{n}`$ is the mean density of galaxies in the sample.
The halo-halo correlation term is given by
$$P_{\mathrm{gg}}^{hh}(k)=P_{\mathrm{lin}}(k)\left[\frac{\overline{\rho }}{\overline{n}}f(\nu )𝑑\nu \frac{N}{M}b(\nu )y(k,M)\right]^2.$$
(14)
This should be modified somewhat because the central galaxy does not contribute a $`y(k,M)`$ term, but this is only important on small scales where the halo-halo term is negligible. On large scales where $`y1`$ this term gives the constant bias model
$$P_{\mathrm{gg}}^{hh}(k)=b^2P_{\mathrm{lin}}(k),$$
(15)
where the mean galaxy bias $`b`$ is given by
$$b=\frac{\overline{\rho }}{\overline{n}}f(\nu )𝑑\nu \frac{N}{M}b(\nu ).$$
(16)
The Poisson term is given by
$$P_{\mathrm{dm}}^P(k)=\frac{\overline{\rho }^2}{(2\pi )^3\overline{n}^2}\frac{M}{\overline{\rho }}f(\nu )𝑑\nu \frac{N(N1)}{M^2}|y(k,M)|^p.$$
(17)
We use the approximation with $`p=2`$ if $`N(N1)>1`$, because in the limit the number of pairs is large it is dominated by the halo galaxies, and $`p=1`$ if $`N(N1)<1`$, because in the opposite limit the number of pairs in this case is dominated by the central galaxy paired with a halo galaxy. Following the usual convention we use $`N(N1)`$ instead of $`N^2`$, since we subtract out the shot noise term arising from the discrete nature of galaxies (such a term does not depend on the halo profile $`y(k,M)`$). Comparing equations 8 and 10 with equations 14 and 17 we see that there is no difference between the two only if $`N/M`$ and $`N(N1)^{1/2}/M`$ are independent of $`M`$, there are many galaxies per halo and the galaxies are distributed as the dark matter within the halo. For such conditions the power spectrum of galaxies is identical to the power spectrum of dark matter.
To test the model above we use semi-analytic models of galaxy formation developed in . These models use N-body simulations to identify the halos and their progenitors. Gas is assumed to follow dark matter initially so that it heats up during the collapse to the virial temperature of the halo. Because of the high density it can efficiently cool and subsequently concentrate at the center of the halo. Stars are formed from this cold gas on the dynamical time scale. The parametrized star formation efficiency and the stellar population synthesis models are used to assign magnitudes in different color bands to the formed galaxies. The small halos with galaxies in them subsequently merge into larger halos and exist as individual galaxies until they merge with the central galaxy on the dynamical friction time scale. The output of these models is a catalog of halos and their masses. For each halo the output consists of a list of galaxies, their positions and luminosities in different bands. From such a catalog one can reconstruct the 3-d distribution of galaxies and dark matter, as well as $`N`$ and $`N(N1)^{1/2}`$ averaged over a given range of halo masses for any desired galaxy selection criterion. The goal of our comparison is to compare the galaxy power spectrum predicted from our model using $`N(M)`$ and $`N(N1)^{1/2}(M)`$ from semi-analytic models to the galaxy power spectrum obtained directly from these models. This is a meaningful comparison even if semi-analytic models do not correctly describe the nature. If we determine that the model contains all the necessary ingredients to predict the galaxy correlations we can then try to obtain these ingredients by other means, either through direct observations or better modelling. This can also be applied in the other direction: from observations of galaxy power spectrum (and galaxy-dark matter power spectrum discussed in the next section) we can determine the ingredients of our model, which must be satisfied by any theoretical galaxy formation model.
A generic outcome of theoretical models such as these is that the amount of cold gas that can be transformed to stars increases as a function of the mass slower than the halo mass itself, because larger halos are hotter and the gas takes longer to cool . In such models one would expect $`N/M`$ to decline with $`M`$. This is shown in figure 3 where $`N`$ and $`N(N1)^{1/2}`$ is plotted versus $`M`$ for galaxies selected only on the basis on absolute magnitude ($`M_B<19.5`$). Both functions have similar dependence for $`M>10^{14}h^1M_{}`$. When the number of galaxies per halo begins to drop below unity the two functions begin to deviate from one another and $`N(N1)^{1/2}`$ drops below $`N`$. This is because only the halos with two or more galaxies contribute to $`N(N1)^{1/2}`$, while single galaxy halos also contribute to $`N`$. However, both functions increase less rapidly than the mass for $`M>10^{13}h^1M_{}`$.
Using $`N(N1)^{1/2}/M`$ and $`N/M`$ for $`M_B<19.5`$ from figure 3 in equations 13-17 we obtain the galaxy power spectrum shown in figure 4. We only show results for $`\alpha =1`$ model, but the $`\alpha =1.5`$ model gives essentially identical results. Also shown is the dark matter power spectrum and its two contributions, as well as the measured APM and scaled IRAS galaxy power spectrum compiled in . First thing to note is the good agreement between our analytical model and the simulations. The agreement is significantly better for this model than for the model where there is no central galaxy, which would give a stronger decline in power on small scales. The galaxy power spectrum is almost a perfect power law over several decades in scale, in agreement with observations and in contrast to the dark matter power spectrum, whose slope gradually decreases with $`k`$. The slope of the galaxy power spectrum is in agreement with the observed slope $`k^3P(k)k^{1.8}`$ and this slope persists in the analytic model down to kpc scales.
It is useful to introduce bias $`b(k)`$, defined as the square root of the ratio between galaxy and dark matter power spectrum,
$$b(k)=[P_{\mathrm{gg}}(k)/P_{\mathrm{dm},\mathrm{dm}}(k)]^{1/2}.$$
(18)
The bias $`b(k)`$ is approximately constant and close to unity on large scales, decreases and becomes less than unity between $`0.3\mathrm{hMpc}^1<k<6\mathrm{h}\mathrm{M}\mathrm{p}\mathrm{c}^1`$ and then increases for large $`k`$. The bias is therefore scale dependent and nonmonotonic, both of which as shown below are generic predictions of this model. On very large scales the power spectrum is dominated by the correlations between the halos and the internal structure of halos can be neglected. This gives the constant bias on large scales, which for the galaxy type considered here is close to unity. On smaller scales the halo Poisson term becomes important both for galaxies and dark matter. However, if $`N(N1)^{1/2}/MM^\psi `$ with $`\psi <0`$ the Poisson term for galaxies is lower than the Poisson term for dark matter in the limit $`y(k,M)=1`$. This is because the halo Poisson term is larger if halos are rarer. If $`\psi <0`$ the dominant contribution in galaxy power spectrum is shifted to lower mass halos, which are more abundant and this reduces the Poisson term relative to dark matter. Another important factor that reduces the galaxy Poisson term is that $`N`$ exceeds $`N(N1)^{1/2}`$ below $`M10^{14}h^1M_{}`$. $`N(M)`$ determines the mean density of galaxies $`\overline{n}`$ in equation 13. This suppresses the Poisson term in equation 17 even if $`\psi =0`$. Suppression of the galaxy Poisson term relative to the dark matter delays the onset of nonlinear power in the galaxy power spectrum relative to the dark matter, which is clearly seen in figure 4. It gives a natural explanation for the position of the inflection point in the observed galaxy power spectrum without the need to introduce phenomenological double power law spectra . While our model is already in a good agreement with the data an even better fit would be achieved with a somewhat smaller Poisson term, which would require $`\psi `$ to be even lower or $`N`$ to exceed $`N(N1)^{1/2}`$ even more. This would further delay the onset of the nonlinear clustering.
On even smaller scales the halo profile $`y(k,M)`$ becomes important, since it begins to decrease from unity at a scale that corresponds to a typical size of the halo, which is smaller for the lower mass halos. Since the galaxy power spectrum is weighted towards lower mass halos relative to dark matter the term $`y(k,M)`$ begins to be important in suppressing the Poisson term at a smaller scale. In addition, if each halo hosts a central galaxy then switching to $`p=1`$ for small halos also makes the suppression by the halo profile less important. Smaller suppression of the galaxy power spectrum relative to the dark matter results in an increase of bias with $`k`$. This argues that the decrease of $`b(k)`$ on intermediate scales and the increase on small scales are generic predictions. The overall result of this is an approximate power law of the galaxy power spectrum over several decades. Such a power law arises quite generically in a CDM family of models where $`\psi <0`$ (or $`N>N(N1)^{1/2}`$) and where each halo hosts a central galaxy. We note that the latter is required to preserve the power law behaviour to very small scales. A model where $`p=2`$ for all halo masses turns below the power law in the power spectrum at $`k>50h`$Mpc<sup>-1</sup>, similar to the dark matter.
The conclusion above that the bias first declines with $`k`$ and then rises again applies to a normal galaxy population. If one selects red galaxies on the basis of color or ellipticals on the basis of morphology then one may expect a different bias dependence, since red or elliptical galaxies are preferentially found in more massive halos, such as groups and clusters. Figure 3 shows that the galaxies selected by $`M_BM_V>0.8`$ and $`M_B<19.5`$ are dominant in halos with $`M>10^{14}h^1M_{}`$, while their relative fraction declines rapidly below that. The dependence of $`N(N1)^{1/2}`$ with $`M`$ is much steeper in this case so that $`\psi >0`$ for $`M<10^{14}h^1M_{}`$. In addition $`NN(N1)^{1/2}`$ across the entire range of halo masses, a consequence of the fact that most of the red galaxies are not central galaxies, which in these models show recent star formation and are therefore not red.
Figure 5 shows the comparison between analytic predictions and results from semi-analytic models . We use mass dependence from figure 3 in equations 13-17 for both $`M_B<19.5`$ and $`M_BM_V>0.8`$ galaxy selection. Also shown are the dark matter power spectrum from the model and from the GIF simulations which were used for semi-analytic models. Qualitatively the agreement is excellent, specially for dark matter and $`M_B<19.5`$ galaxies, while for red galaxies semi-analytic models predict a somewhat higher amplitude. Part of the disagreement is caused simply by dark matter spectrum not being in agreement with PD (our models are chosen so that they agree with PD) . We do not show small scales ($`k>20h`$Mpc<sup>-1</sup>) where limited resolution of N-body simulations prevents a meaningful comparison. The remaining discrepancy for red galaxies between $`0.5h`$Mpc$`{}_{}{}^{1}<k<20h`$Mpc<sup>-1</sup> can only be explained by them not tracing exactly dark matter distribution in halos with $`M>10^{13}h^1M_{}`$. The red galaxies must be more centrally concentrated than dark matter in semi-analytic models in order that their power spectrum has a higher amplitude than predicted from our model. This is in agreement with direct analysis of galaxy distribution inside halos using the same simulations , where it was found that red galaxies in $`\mathrm{\Lambda }CDM`$ model tend to be more centrally concentrated that dark matter. Galaxies that form first end up more towards the center of the cluster because the violent relaxation during the merging is incomplete.
In the case of the red galaxies the bias starts with a value larger than unity on large scales. This is because most of the red galaxies are in clusters which are biased relative to the dark matter following equation 6. Bias first rises with $`k`$ and then declines. This is just the opposite from the scale dependence of the normal galaxies and is a consequence of $`\psi >0`$ for $`M<10^{14}h^1M_{}`$ and $`NN(N1)^{1/2}`$. This gives rise to the Poisson term larger for the galaxies than for the dark matter on large scales. This conclusion is again independent of the distribution of the galaxies inside the halos. This is confirmed in figure 5 where on large scales our model agrees very well with the semi-analytic predictions. Because the galaxies are preferentially in larger halos relative to the dark matter $`y(k,M)`$ suppression is more important and the bias declines on smaller scales. This is seen in the power spectrum from the simulations. In our model it begins to rise again on even smaller scales because $`p`$ switches to unity for $`M<10^{14}h^1M_{}`$, resulting in a smaller suppression by the halo profile. This effect is not seen in the simulations, presumably because of their limited resolution.
It is important to note that bias may never be really constant even on scales above $`100h^1`$Mpc. For the red sample it changes by 30% between $`k=0.01h`$Mpc<sup>-1</sup> and $`k=0.1h`$Mpc<sup>-1</sup>. This is because the Poisson term does not become much smaller than the halo-halo term even on very large scales, a consequence of the fact that the slope of $`P_{\mathrm{lin}}(k)`$ and thus the halo-halo term itself becomes flat and even positive on very large scales (approaching $`n1`$ on very large scales). Since even at the turnover of the power spectrum (where $`n0`$) the Poisson term for the red galaxies is of the order of 20% of the halo-halo term the bias does not become constant and begins to increase again on scales larger than the scale of the turnover. In fact on very large scales ($`k<10^3h`$Mpc<sup>-1</sup>) the red galaxy power spectrum becomes white noise, although these scales are already approaching the size of the visible universe. It should be noted that this description is valid on large scales only for galaxies which do not obey mass and momentum conservation. For the dark matter mass and momentum conservation require that the Poisson term vanishes on large scales and any spectrum generated by a local process should decrease faster than $`P(k)k^4`$ as $`k0`$ . Galaxies do not obey mass and momentum conservation and can have the Poisson contribution, so the qualitative scale dependence of bias remains as predicted above.
We have concentrated on the power spectrum above because it is the quantity that can be most directly compared to the theoretical predictions. The same analysis could however be applied to the correlation function as well. The power law dependence of the power spectrum would also result in a power law correlation function, so the conclusions would remain unchanged. The main difference in the real space is that the Poisson term is localized to scales smaller than the typical halo scale and vanishes on scales above that. In this case bias would be scale dependent up to this typical scale (of order few Mpc), but would become scale independent on scales above that. There is no need to model the Poisson term on large scales at all. In this sense the real space correlation function offers some advantages over the power spectrum, where one must attempt to remove the Poisson term in the power spectrum by modelling it as a constant term on large scales.
Our predictions agree with the results of semi-analytic models, indicating that the here proposed model is sufficient to extract the key ingredients to model the galaxy clustering. This means one does not need to rely on N-body simulations as long as the ingredients of the model are specified. If one can extract $`N`$, $`N(N1)^{1/2}(M)`$ and $`y(k,M)`$ directly from the data one can sidestep the theoretical modelling of this relation and predict the galaxy power spectrum directly . It is in principle possible to obtain such information at least for the massive halos by combining dynamical information on galaxy groups and clusters, such as X-ray temperature, velocity dispersion or weak lensing mass, with the number of galaxies in these clusters. Existing data such as CNOC survey indeed find that $`N/M`$ for galaxies with $`M_K<18.5`$ is systematically lower in massive clusters with $`\sigma >1000`$ km/s than in poorer clusters. The current data are sparse, but new large surveys such as SDSS and 2dF will enable one to extract such information with a much better statistics. This could allow one to determine within our model the dark matter power spectrum from the galaxy power spectrum directly.
Another direction to obtain $`N/M`$ is to require consistency with other measurements that combine dynamical and galaxy information. Galaxy-dark matter correlations discussed in the next section are one possibility. Another are pairwise velocity dispersion measurements. If $`N/M`$ declines with $`M`$ then the pairwise velocity dispersion for the galaxies will be lower than for the dark matter . This is because there will be more pairs of galaxies in smaller halos relative to the dark matter. Smaller halos have smaller velocity dispersions and smaller relative velocities between the particles. This can explain the lower amplitude of pairwise velocity dispersion in the LCRS data compared to the N-body simulations . The required value of $`\psi 0.1`$ has indeed the same sign as required to reproduce the delayed onset of nonlinear clustering and the power law in galaxy power spectrum. It would be interesting to see whether a single set of functions $`N(M)`$, $`N(N1)^{1/2}(M)`$ can provide a unified description to both galaxy clustering and pairwise velocities within the CDM models.
## IV Dark matter-galaxy cross-correlation
Dark matter-galaxy cross-correlations are measured whenever a galaxy is cross-correlated with a tracer of the dark matter. Examples of this are galaxy-galaxy lensing , where one is measuring correlation between galaxies and cosmic shear, and correlations between foreground and background galaxies or quasars , where correlations (or anti-correlations) are induced by magnification bias of background objects. In both cases one is measuring the correlations between the galaxies and dark matter along the line of sight, which can be expressed as a convolution over the galaxy-dark matter cross-correlation power spectrum.
Galaxy-dark matter cross-correlations have been modelled in the past using either a bias model relating them to the dark matter or galaxy power spectrum or using galaxies sitting at the centers of the galactic size halos . In the first description assuming galaxy-dark matter cross-correlations measure bias $`b(k)P_{\mathrm{dm}}(k)`$, which in combination with the galaxy power spectrum $`b^2(k)P_{\mathrm{dm}}(k)`$ can give both $`b(k)`$ and $`P_{\mathrm{dm}}(k)`$. Such a model is a reasonable description on large scales, but must break down on small scales where galaxies do not trace dark matter and there is no guarantee that the scale dependent bias that relates $`P_{\mathrm{gal},\mathrm{dm}}(k)`$ and $`P_{\mathrm{gal}}(k)`$ can be used to extract $`P_{\mathrm{dm}}(k)`$.
Second model describes cross-correlations in terms of galaxies sitting at the centers of their halos and interprets the results in terms of the averaged halo profile . There are two potential problems with this approach. First, there may be more than one galaxy inside the halo, which is specially important for large halos (figure 3). Since not all galaxies can lie at the halo center this can affect the interpretation of the cross-correlations in terms of the halo profile. Second, just as in the case of the dark matter the contribution to the power spectrum comes from a range of halo masses and one cannot model the galaxy-dark matter cross-correlation simply as a typical $`L_{}`$ galaxy halo profile. The strength of the correlations is determined both by the dark matter profile of the halos as well as by the halo mass function, so the slope of the correlation function that one is ultimately measuring with galaxy-galaxy lensing and foreground-background galaxy correlations need not be directly related to the dark matter profile . Model developed in previous sections may be applied to the dark matter-galaxy cross-correlation power spectrum to quantify these issues in more detail.
Galaxy-dark matter cross-correlation power spectrum has halo-halo and halo Poisson terms. First term describes the correlations between galaxies and dark matter in neighbouring halos and is dominant on large scales. Second term includes the correlations between the galaxies and dark matter in the same halo and dominates on small scales. The halo-halo term is given by
$$P_{\mathrm{g},\mathrm{dm}}^{hh}(k)=P_{\mathrm{lin}}(k)\left[\frac{\overline{\rho }}{\overline{n}}f(\nu )𝑑\nu \frac{N}{M}b(\nu )y[k,M(\nu )]\right]\left[f(\nu )𝑑\nu b(\nu )y[k,M(\nu )]\right].$$
(19)
On large scales where this term dominates it reduces to constant bias model, $`P_{\mathrm{g},\mathrm{dm}}^{hh}(k)=bP_{\mathrm{lin}}(k)`$. The Poisson term in the model where galaxies trace dark matter inside the halos except for the central galaxy sitting at its center is given by
$$P_{\mathrm{g},\mathrm{dm}}^P(k)=\frac{1}{(2\pi )^3\overline{n}}f(\nu )𝑑\nu N|y(k,M)|^p.$$
(20)
Here $`p=2`$ for $`N>1`$ and $`p=1`$ for $`N<1`$.
Figure 6 shows the results for the cross-correlation power spectrum for the same galaxy selection as in figures 3 and 4. For regular galaxies selected by an absolute magnitude (top panel) the cross-correlation spectrum is similar to the galaxy power spectrum. If we define the cross-correlation coefficient as
$$r(k)=\frac{P_{\mathrm{g},\mathrm{dm}}(k)}{[P_{\mathrm{dm}}(k)P_\mathrm{g}(k)]^{1/2}},$$
(21)
then we see from figure 6 that it is approximately unity up to $`k1h`$Mpc<sup>-1</sup> and increases for higher $`k`$. Note that the cross-correlation coefficient is not restricted to $`|r(k)|<1`$ because we have subtracted out the shot noise term from the galaxy power spectrum following the usual approach . Because on small scales the galaxy and cross-correlation spectra are comparable and exceed the dark matter spectrum the cross-correlation coefficient grows to large values in this model. Comparison with the semi-analytic results again shows very good agreement up to the resolution limit of the simulations. Bottom of figure 6 shows the results for the red galaxies. In this case the cross-correlation spectrum falls in between the dark matter and the galaxy spectrum, so that $`r1`$ down to very small scales. This is again in agreement with semi-analytic results which show $`r1`$ throughout the entire range of $`k`$.
The main reason for $`r1`$ on small scales is that $`N(N1)^{1/2}N`$ (figure 3). The difference between the two functions is more significant for the normal than for the red galaxies, which is why the cross-correlation coefficient begins to deviate from unity at larger scales for $`M_B<19.5`$ than for $`M_BM_V>0.8`$. Because in this regime $`N(N1)^{1/2}<N`$ this leads to $`r(k)>1`$, as seen in figure 6. It is interesting to note from figure 3 that for the red galaxies the two functions agree very well even below unity and this leads to $`r(k)1`$ down to very small scales. When this happens one can reconstruct the dark matter power spectrum from the galaxy and cross-correlation spectrum even if most of the dark matter halos are not directly observed. Unfortunately one cannot extract these two functions without first identifying the dark matter halos, so this prediction cannot be directly verified from observational data using the galaxy information only.
Second source of stochasticity is the presence of central galaxy. For those halos where $`N(N1)^{1/2}<1`$ or $`N<1`$ only one power of $`y(k,M)`$ is used as opposed to two in the case of the dark matter. This induces some stochasticity even if $`N=N(N1)^{1/2}`$, because it enhances the galaxy-galaxy and galaxy-dark matter spectrum above the dark matter-dark matter spectrum. Another source of stochasticity would be $`\psi 0`$, which would make correlations at a given scale being dominated by different mass range in the case of the dark matter and the galaxies. Calculations where only this effect is present give $`r1`$ over a wide range of scales, showing that this cannot be a significant source of stochasticity, at least for reasonable values of $`\psi `$.
Our model predicts that even if the constant bias model is not valid, its generalization $`r=1`$ model may be a reasonable approximation at least down to 1 Mpc scales. An example are the red galaxies (bottom of figure 6), which have very strong scale dependent bias, yet $`r1`$ over a wide range of scales. In this sense determining the dark matter power spectrum from the measurements of galaxy-galaxy spectrum and galaxy-dark matter spectrum under the assumption of $`r=1`$ may have a larger range of validity than the constant bias model. This relies on the assumption $`N=N(N1)^{1/2}`$ predicted from these models. This prediction can be verified at least for the more massive halos directly from observations, for example by using galaxy counts in cluster catalogs to extract $`N`$ and $`N(N1)^{1/2}`$. Such an approach would provide an alternative way to determine $`r(k)`$ directly from the data.
The model developed here can also be used to clarify the interpretation of the galaxy-dark matter cross-correlation in terms of an averaged density profile of a typical galaxy. Figure 7 shows the contribution to the cross-correlation spectrum from the different halo mass intervals, similar to figure 2 for the dark matter. For $`k<20h`$Mpc<sup>-1</sup>, corresponding approximately to scales larger than 100h<sup>-1</sup>kpc in real space, one cannot interpret the correlations in terms of the shape of a single halo profile, but instead as the convolution of these over the halo mass function, multiplied with the number of galaxies per halo. Observed correlations on large scales do not necessarily mean that the halo of an $`L_{}`$ galaxy extends to large distances. Instead, it is more likely that one is observing correlations arising from the group and cluster size halos, which exceed the correlations contributed from the galactic size halos on larger scales. This cannot be corrected in any simple manner by taking into account the correlation function of the galaxies , which attempts to model the presence of other nearby halos. Even if the galaxy correlations vanished one would still need to take into account the halo mass function and the fact that different halos dominate on different scales. More detailed discussion of these points will be presented elsewhere .
On smaller scales the transition to $`N(M)<1`$ implies that $`y(k,M)`$ suppression is less important because $`p=1`$. This is further enhanced by the flattening of $`N`$ below $`M10^{13}h^1M_{}`$ as seen in figure 3. In addition, galaxies selected on the basis of their absolute magnitude cannot exist in very small halos, so the mass function has a strong cutoff below $`10^{12}h^1M_{}`$. Thus on scales with $`k>20h`$Mpc<sup>-1</sup> the galaxy-dark matter cross-correlation may be better interpreted in terms of the average profile of $`10^{12}h^1M_{}<M<10^{13}h^1M_{}`$ halos. However, this may not be a robust prediction since the semi-analytic predictions in figure 3 are highly uncertain over this mass range. A small change in $`N(M)`$ may lead to a larger influence of the mass function on the power spectrum, making the correspondence between the halo profile and the power spectrum less certain. In general one should be cautious in interpreting the shape of the galaxy-dark matter correlation function in terms of an averaged dark matter profile.
## V Conclusions
We developed an analytic model for computing the power spectrum of the dark matter, galaxies and their cross-correlation based on the Press-Schechter model. In this model all the matter in the universe is divided into virialized halos. These halos cluster and have some internal profile. The total power spectrum is the sum of the halo clustering term and the halo Poisson term, which accounts for the correlations within the halos. We assume that the halo profiles are self-similar regardless of the initial conditions, but with the mass dependent concentration parameter, as suggested by high resolution simulations . The model agrees well with the results of N-body simulations for the $`\mathrm{\Lambda }CDM`$ model. We are able to find a good agreement for inner slopes $`\alpha =1`$ and $`\alpha =1.5`$, indicating that the shape of the nonlinear power spectrum cannot by itself distinguish between the two.
The model can in principle be applied to any cosmological model, including those with a cutoff in the linear power spectrum on small scales or with some features in the power spectrum. While this will be explored in more detail in a future paper we wish to emphasize here that the mass function, which is sensitive to the linear power spectrum, has a direct effect on the nonlinear power spectrum through the halo abundance, so that not all of the information on the linear power is lost in the nonlinear regime. For example, if the linear power spectrum is cut-off on small scales and if inner profile $`\alpha >1.5`$ as suggested by the simulations then the correlation function or $`k^3P(k)`$ must have a turnover on small scales. This differs from the CDM models which predict the nonlinear correlation function to continue to grow on small scales. If we wish to eliminate the halos with $`M<10^{11}h^1M_{}`$ then this would suppress the power on scales below 10kpc (figure 2). This effect therefore becomes significant on scales smaller than those resolved in a recent study of such truncated power spectrum models .
Our main conclusion regarding the galaxy power spectrum is that a simple model for the dependence of the linear and pair weighted number of galaxies inside halo as a function of the halo mass can explain most of the properties of the galaxy clustering seen in more complicated models based on the N-body simulations. A power law in the galaxy correlation function with slope $`1.8`$ is a generic prediction of the model where the number of galaxies inside the halo increases less rapidly with mass than the halo mass itself, mean number of galaxies exceeds pair weighted average and there is a central galaxy in each halo. The decline of number of galaxies per unit mass as a function of mass is predicted by the galaxy formation models and has been observed in clusters . It is also required to explain the pairwise velocity dispersion results . For such galaxies bias first decreases below unity, because the Poisson term is smaller for them than for dark matter. This naturally explains the later onset of nonlinearity in galaxy power spectrum compared to the dark matter, which reconciles the discrepancy between the data and the CDM models . Conversely, there is no need to invoke poorly motivated models such as double power law model . On large scales bias converges to a constant for these galaxies.
Red or elliptical galaxies, which are more abundant in massive halos, show a different relation: their number inside the halos increases on average more rapidly than the halo mass. In this case bias increases with $`k`$ above the turnover in the power spectrum ($`k0.01h`$Mpc<sup>-1</sup>), because their Poisson term is larger than that of dark matter. In fact, the Poisson term may be so strong that it may not be negligible compared to the halo clustering term even on very large scales and one may not converge to the constant bias model.
Galaxy-dark matter correlations can also be predicted by this model. In this case one must specify the average number of galaxies per halo as a function of halo mass. Here again our model reproduces the main features present in the N-body simulations with semi-analytic galaxy formation . Galaxy-dark matter cross-correlations can be measured with galaxy-galaxy lensing or correlations between foreground and background galaxies and may provide a way to break some of the uncertainties present with the galaxy clustering. For example, we have shown that even if the constant bias may not be a good approximation, cross-correlation coefficient may nevertheless be close to unity down to Mpc scales, which would allow one to extract the dark matter power spectrum from the knowledge of the galaxy and cross-corelation spectrum on scales larger than this. The main source of stochasticity ($`r1`$) arises from the pair weighted number of galaxies inside the halo differing from the mean number of galaxies and from the (possible) existence of central galaxies in the halos.
We have emphasized that caution must be applied when interpreting the cross-correlations such as galaxy-galaxy lensing in terms of an averaged density profile of a halo. As we have shown different halo masses dominate on different scales and the correlation function reflects this combined effect of all the halos. For example, correlations at a few hundred kpc observed by galaxy-galaxy lensing are more likely to be caused by group and cluster sized halos at $`r_s`$ distances than by galaxy sized halos at $`r_v`$ distances. More detailed work is needed to extract the structure and extent of the dark matter halos from such observations.
Perhaps the most promising direction to explore in the future is to extract the functional dependences that parametrize our model directly from the observations. If one can determine the linear and pair weighted number of galaxies as a function of halo mass and their distribution inside the halos then one can determine the galaxy power spectrum directly within this model. Similarly if one can compare the mean number of galaxies with the pair weighted number as a function of halo mass then one can predict the galaxy-dark matter cross-correlation coefficient. This is certainly feasible for clusters, which dominate the Poisson term on large scales. Current data are sparse , but new surveys such as SDSS or 2dF should provide sufficient statistics to make this feasible. This approach would provide an independent estimate of the scale dependence of bias and correlation coefficient on large scales. It will also provide important constraints that would need to be satisfied by any viable galaxy formation model.
I ackowledge the support of NASA grant NAG5-8084. I thank G. Kauffmann and S. White for a detailed reading of the manuscript and for providing results of GIF N-body and semi-analytic simulations and J. Guzik for help with them. I also thank R. Sheth and R. Scoccimarro for useful conversations and help in initial stages of this project and J. Peebles and U. Pen for useful discussions. |
warning/0001/hep-ph0001076.html | ar5iv | text | # Flavor asymmetry in polarized proton-deuteron Drell-Yan process
## 1 Introduction
Antiquark distributions used to be considered flavor symmetric for the light antiquark distributions ($`\overline{u}=\overline{d}`$). However, it became apparent that they are significantly different $`^\mathrm{?}`$ because of the experimental findings of Gottfried-sum-rule violation and proton-deuteron asymmetry in Drell-Yan experiments. Although meson-cloud-type models seem to be the promising explanation, various models have been proposed to explain the experimental data. In order to test the theoretical ideas, we need further experimental information. In particular, polarized flavor asymmetries should provide crucial information for determining the physics mechanism behind the flavor asymmetric distributions.
Longitudinally-polarized parton distributions have been investigated mainly through the structure functions $`g_1`$ for the proton, neutron (or <sup>3</sup>He), and deuteron.$`^\mathrm{?}`$ Semi-inclusive data were also obtained; however, they are not accurate enough at this stage to provide any significant constraint on the polarized distributions.$`^\mathrm{?}`$ Therefore, the antiquark distributions are assumed to be flavor symmetric in almost all the parametrizations.$`^\mathrm{?}`$ We expect that the situation will become clearer by the RHIC-Spin and other experiments. As far as the transversity is concerned, there is no experimental information yet. Although experiments could be done at RHIC and HERA, nobody knew how to measure the light-antiquark flavor asymmetry because the transversity distributions cannot be measured in the inclusive lepton scattering and W-production experiments due to the chiral-odd property. Reference 4 proposed that the polarized proton-deuteron ($`pd`$) Drell-Yan process could be used in combination with the proton-proton ($`pp`$) Drell-Yan for extracting the flavor asymmetry information. We discuss such possibility in this talk.
In Sec. 2, relations between the polarized $`pd`$ Drell-Yan process and the polarized parton distributions are introduced. Then, we discuss how to extract the light antiquark distributions from the polarized $`pd`$ Drell-Yan cross section in Sec. 3. Our studies are summarized in Sec. 4.
## 2 Proton-deuteron Drell-Yan process
Unpolarized Drell-Yan processes have been studied as an alternative method to lepton scattering for finding parton distributions in the nucleon and nuclei. The unpolarized and polarized $`pp`$ Drell-Yan processes have been investigated theoretically for a long time. In addition, the unpolarized $`pd`$ Drell-Yan is used experimentally for extracting the flavor-asymmetry information ($`\overline{u}/\overline{d}`$). If the proton and deuteron are polarized in the $`pd`$ Drell-Yan process, we could investigate a polarized version of the flavor asymmetry. This topic is discussed in Sec. 3. However, it is not straightforward to express the $`pd`$ cross section in terms of structure functions. In particular, it was not clear how the tensor structure is involved in the polarized cross section. Reference 5 clarified this point.
In formulating the polarized $`pd`$ Drell-Yan, we tried two different methods.$`^\mathrm{?}`$ The first one uses the Jacobi-Wick helicity formalism with the spin-density matrices. The essential difference from the $`pp`$ reaction is that there exist rank-two tensors due to the spin-1 nature of the deuteron. The conditions of Hermiticity, parity conservation, and time-reversal invariance are imposed on possible structure functions. Then, we found that there are 108 structure functions in general. If they are integrated over the virtual-photon transverse momentum $`\stackrel{}{Q}_T`$, there exist only 22 structure functions. In the second method, the hadron tensor is expressed in terms of possible combinations of momentum and spin vectors by imposing the same three conditions. Only the limiting case $`Q_T0`$ is considered in this analysis, and we also obtained the same 22 structure functions. These finite functions should be physically significant ones which could be investigated by the polarized $`pd`$ Drell-Yan process.
Considering the present situation on the proton spin physics, we think that the 22 functions are still too many to be investigated seriously. Furthermore, the physics meaning of these functions, particularly the new ones which do not exist in the $`pp`$ reaction, is not clear. Therefore, the $`pd`$ Drell-Yan was also analyzed in a parton model.$`^\mathrm{?}`$ The hadron tensor is expressed by correlation functions for the process $`q+\overline{q}\mathrm{}^++\mathrm{}^{}`$. Then, the correlation functions are expanded in terms of the sixteen $`4\times 4`$ matrices: $`\mathrm{𝟏},\gamma _5,\gamma ^\mu ,\gamma ^\mu \gamma _5,\sigma ^{\mu \nu }\gamma _5`$ and kinematical Lorentz vectors and pseudovectors. Then, we found finite structure functions in the parton model. There is a new polarization asymmetry $`A_{UQ_0}`$ with the unpolarized proton and the tensor-polarized deuteron. It is given by the parton distributions in the proton and deuteron as
$$A_{UQ_0}=\frac{_ae_a^2\left[f_1(x_1)\overline{b}_1(x_2)+\overline{f}_1(x_1)b_1(x_2)\right]}{_ae_a^2\left[f_1(x_1)\overline{f}_1(x_2)+\overline{f}_1(x_1)f_1(x_2)\right]},$$
(1)
where $`f_1(x)`$ and $`\overline{f}_1(x)`$ are unpolarized quark and antiquark distributions, and $`b_1(x)`$ and $`\overline{b}_1(x)`$ are tensor-polarized distributions. The $`b_1`$ structure function is known in lepton scattering; $`^\mathrm{?}`$ however, the Drell-Yan process provides important information on the antiquark tensor polarization $`\overline{b}_1`$. However, this topic is no more discussed in the following because it is not the major purpose of this paper to investigate the tensor structure. We refer the reader to Ref. 5 for more details. In the following, we discuss the double longitudinal and transverse spin asymmetries in connection with the flavor asymmetry.
First, according to the general formalism,$`^\mathrm{?}`$ the difference between the longitudinally-polarized $`pd`$ cross sections is given by
$$\mathrm{\Delta }\sigma _{pd}=\sigma (_L,1_L)\sigma (_L,+1_L)\frac{1}{4}\left[2V_{0,0}^{LL}+(\frac{1}{3}cos^2\theta )V_{2,0}^{LL}\right],$$
(2)
where $`\sigma (pol_p,pol_d)`$ indicates the cross section with the proton polarization $`pol_p`$ and the deuteron one $`pol_d`$. The longitudinally polarized structure functions $`V_{0,0}^{LL}`$ and $`V_{2,0}^{LL}`$ are defined in Ref. 5. The $`\theta `$ is the polar angle of the lepton $`\mathrm{}^+`$. Then, the structure functions are related to the polarized parton distributions in the parton-model analysis.$`^\mathrm{?}`$ The $`\stackrel{}{Q}_T`$-integrated results indicate
$$\mathrm{\Delta }\sigma _{pd}\underset{a}{}e_a^2\left[\mathrm{\Delta }q_a(x_1)\mathrm{\Delta }\overline{q}_a^d(x_2)+\mathrm{\Delta }\overline{q}_a(x_1)\mathrm{\Delta }q_a^d(x_2)\right],$$
(3)
where $`\mathrm{\Delta }q_a^d`$ and $`\mathrm{\Delta }\overline{q}_a^d`$ are the longitudinally-polarized quark and antiquark distributions in the deuteron.
In the transverse-polarization asymmetry, the situation is more complicated in the sense that four structure functions ($`V_{0,0}^{TT}`$, $`V_{2,0}^{TT}`$, $`U_{2,2}^{TT}`$, and $`U_{2,1}^{UT}`$) contribute. However, it becomes a simple expression if the parton model is used by neglecting higher-twist contributions:
$`\mathrm{\Delta }_T\sigma _{pd}`$ $`=\sigma (\varphi _p=0,\varphi _d=0)\sigma (\varphi _p=0,\varphi _d=\pi )`$
$`{\displaystyle \underset{a}{}}e_a^2\left[\mathrm{\Delta }_Tq_a(x_1)\mathrm{\Delta }_T\overline{q}_a^d(x_2)+\mathrm{\Delta }_T\overline{q}_a(x_1)\mathrm{\Delta }_Tq_a^d(x_2)\right],`$ (4)
where $`\mathrm{\Delta }_Tq`$ and $`\mathrm{\Delta }_T\overline{q}`$ are quark and antiquark transversity distributions, and $`\varphi `$ is the azimuthal angle of a polarization vector. In this way, we found that the cross-section difference is written in terms of the longitudinally-polarized and transversity distributions.
## 3 Light-antiquark flavor asymmetry
Because the expressions of Eqs. (3) and (4) are the same as the unpolarized one if the polarized distributions are replaced by the unpolarized ones, the polarized flavor asymmetries could be extracted from the polarized $`pp`$ and $`pd`$ Drell-Yan cross sections as it has been investigated in the unpolarized case. In order to discuss the $`pp`$ and $`pd`$ cross sections in connection with the flavor asymmetry, we define the ratio
$$R_{pd}\frac{\mathrm{\Delta }_{\left(T\right)}\sigma _{pd}}{2\mathrm{\Delta }_{\left(T\right)}\sigma _{pp}}=\frac{_ae_a^2\left[\mathrm{\Delta }_{\left(T\right)}q_a\left(x_1\right)\mathrm{\Delta }_{\left(T\right)}\overline{q}_a^d\left(x_2\right)+\mathrm{\Delta }_{\left(T\right)}\overline{q}_a\left(x_1\right)\mathrm{\Delta }_{\left(T\right)}q_a^d\left(x_2\right)\right]}{2_ae_a^2\left[\mathrm{\Delta }_{\left(T\right)}q_a\left(x_1\right)\mathrm{\Delta }_{\left(T\right)}\overline{q}_a\left(x_2\right)+\mathrm{\Delta }_{\left(T\right)}\overline{q}_a\left(x_1\right)\mathrm{\Delta }_{\left(T\right)}q_a\left(x_2\right)\right]},$$
(5)
where $`\mathrm{\Delta }_{(T)}=\mathrm{\Delta }`$ or $`\mathrm{\Delta }_T`$ depending on the longitudinal or transverse case. There is another issue in calculating the cross sections because of nuclear corrections. However, they are ignored in the following discussions since they are not expected to be the essential part. If experimental data are taken in future, such corrections should be taken into account properly.
We show our numerical-analysis results for the ratio $`R_{pd}`$ in Fig. 1. The parton distributions are taken from Ref. 7 at $`Q^2`$=1 GeV<sup>2</sup>, where flavor asymmetry ratio is introduced as $`r_{\overline{q}}\mathrm{\Delta }_{(T)}\overline{u}/\mathrm{\Delta }_{(T)}\overline{d}=`$0.7, 1.0, or 1.3. Then, the distributions are evolved to $`Q^2=M_{\mu \mu }`$=25 GeV<sup>2</sup> by the leading-order DGLAP evolution equations. The center-of-mass energy is taken as $`\sqrt{s}`$=50 GeV with a fixed target experiment in mind. In Fig. 1, the solid and dashed curves are longitudinally- and transversely-polarized ratios, respectively. Because the transversity distributions are assumed to be the same as the corresponding longitudinally-polarized ones at $`Q^2`$=1 GeV<sup>2</sup>, the transverse ratios are almost the same as the longitudinal ones. There are large differences between the curves for $`r_{\overline{q}}`$=0.7, 1.0, and 1.3, so that it should be possible to extract the longitudinally-polarized and transversity flavor asymmetries from the ratios $`R_{pd}^{(L)}`$ and $`R_{pd}^{(T)}`$. The differences are especially large in the large-$`x_F`$ region with the following reason. If two extreme limits ($`x_F=x_1x_2\pm 1`$) are taken in Eq. (1) with the assumption $`\mathrm{\Delta }_{(T)}u_v(x1)\mathrm{\Delta }_{(T)}d_v(x1)`$, the ratio becomes
$`R_{pd}(x_F+1)`$ $`={\displaystyle \frac{1}{2}}\left[\mathrm{\hspace{0.17em}1}+{\displaystyle \frac{\mathrm{\Delta }_{(T)}\overline{d}(x_2)}{\mathrm{\Delta }_{(T)}\overline{u}(x_2)}}\right]_{x_20},`$ (6)
$`R_{pd}(x_F1)`$ $`={\displaystyle \frac{1}{2}}\left[\mathrm{\hspace{0.17em}1}+{\displaystyle \frac{\mathrm{\Delta }_{(T)}\overline{d}(x_1)}{4\mathrm{\Delta }_{(T)}\overline{u}(x_1)}}\right]_{x_10}.`$ (7)
These equations suggest that the flavor-asymmetric distribution $`\mathrm{\Delta }_{(T)}\overline{u}\mathrm{\Delta }_{(T)}\overline{d}`$ can be extracted by finding the deviation from 1 at $`x_F+1`$ or from 5/8 at $`x_F1`$. However, $`R_{pd}`$ should be more sensitive to the flavor asymmetry at large $`x_F`$ due to the factor of 1/4 in Eq. (7) in comparison with Eq. (6).
Next, we show the numerical results in Fig. 2 for larger energies $`\sqrt{s}`$=200 and 500 GeV by considering a collider option. The solid, dashed, dotted curves are for $`\sqrt{s}`$=50, 200, and 500 GeV, respectively. Although there are large variations in the medium-$`x_F`$ region, the ratios in the large- and small-$`x_F`$ regions stay the same. Furthermore, we studied the parametrization-model dependence,$`^\mathrm{?}`$ and the results indicate that the large- and small-$`x_F`$ ratios are again rather independent of the parametrization. Therefore, these regions are appropriate for investigating the flavor asymmetry. In addition, the variations in the medium-$`x`$ region indicate that the details of the polarized parton distributions could be investigated in this region.
In this way, we find that it is possible to extract both $`\mathrm{\Delta }\overline{u}/\mathrm{\Delta }\overline{d}`$ and $`\mathrm{\Delta }_T\overline{u}/\mathrm{\Delta }_T\overline{d}`$ from the cross-section ratios, particularly in the large-$`x_F`$ region. Our suggestion should be important for the transversity because it cannot be measured in the inclusive lepton scattering and W-production processes.
At this stage, no actual experimental measurement is planned. However, there are certain possibilities at Fermilab, HERA-N, and JHF for measuring the polarized $`pd`$ Drell-Yan cross sections by using a fixed deuteron target. In addition, the polarized deuteron could be accelerated in principle at RHIC. However, it is not easy to attain the longitudinal polarization due to the small magnetic moment unless someone has a smart idea for the longitudinal polarization.$`^\mathrm{?}`$ In any case, we should be able to investigate at least the transverse part at RHIC. As far as the tensor polarization is concerned in Eq. (1), we may combine the transverse cross sections with the unpolarized one for getting the tensor-polarized cross section. Because there are a variety of interesting topics on polarized deuteron reactions, we hope that experimental possibilities are seriously studied.
## 4 Summary
First, we briefly discussed the general framework of the polarized $`pd`$ Drell-Yan process. Then, we explained the relation between the flavor-asymmetry ratio $`\mathrm{\Delta }_{(T)}\overline{u}/\mathrm{\Delta }_{(T)}\overline{d}`$ and the Drell-Yan cross-section ratio $`\mathrm{\Delta }_{(T)}\sigma _{pd}/[2\mathrm{\Delta }_{(T)}\sigma _{pp}]`$. Our numerical analysis suggested that the polarized flavor asymmetry could be extracted from the $`pd`$ and $`pp`$ cross-section measurements particularly in the large-$`x_F`$ region. At this stage, this is the only proposal for extracting the transversity asymmetry $`\mathrm{\Delta }_T\overline{u}/\mathrm{\Delta }_T\overline{d}`$ due to the chiral-odd property.
## Acknowledgments
S.K. was partly supported by the Grant-in-Aid for Scientific Research from the Japanese Ministry of Education, Science, and Culture under the contract number 10640277.
## References |
warning/0001/hep-lat0001018.html | ar5iv | text | # Calibration of Smearing and Cooling Algorithms in SU(3)-Color Gauge Theory
## I Introduction
Cooling is a well established method for locally suppressing quantum fluctuations in gauge field configurations. The locality of the method allows topologically nontrivial field configurations to survive numerous iterations of the cooling algorithm. Application of such algorithms are central to removing short-distance fluctuations responsible for causing renormalization constants to significantly deviate from one as is the case for the topological charge operator .
An alternative yet somewhat related approach is the application of APE smearing to the gauge field configurations. Because the algorithm is applied simultaneously to all link variables, without any annealing in the process of a lattice update, the application of the algorithm may be viewed as an introduction of higher dimension irrelevant operators designed to improve the operator based on unsmeared links. This approach also leads to the desirable effect of renormalization constants approaching one. The net effect may be viewed as introducing a form factor to gluon vertices that suppresses large $`q^2`$ interactions and, as a result, suppresses lattice artifacts at the scale of the cutoff.
The focus of this paper is to calibrate the relative rates of cooling and smearing as measured by the action and the topological charge. In doing so we aim to establish relationships between various algorithms and their internal parameters. It is well known that errors in the standard Wilson action eventually destroy (anti)instanton configurations and this effect is represented in the topological charge as sharp transitions from one integer to another. We find the timing analysis of these topological charge defects to be central to understanding the behavior of the various algorithms.
To bridge the gap between cooling and smearing and discover the origin of variations between these algorithms, we present a new algorithm for smearing the gauge fields which we refer to as Annealed $`U`$ Smearing or AUS smearing. In addition to comparing cooling and smearing, we also calibrate the effects of different smearing fractions, describing the balance between the unsmeared link and the smeared link.
There have been numerous studies of the dependence of vacuum matrix elements on the number of smoothing sweeps. Of particular note is the observation of cooling as a diffusive procedure , such that the distance $`d`$ affected by cooling is related to the square-root of the number of sweeps over the lattice, $`n`$. This random walk ansatz was further investigated by introducing the lattice spacing $`a`$, such that $`d=a\sqrt{n}`$ takes on a physical length . With this simple formula, it is possible to effectively relate results from lattices of different $`\beta `$ ($`a`$) and different smoothing sweeps $`n`$. Extensive work on removing the dependence of particular quantities on the number of smoothing sweeps was carried out in Ref. . There, highly improved actions are fine tuned to stabilize the instanton action and size under cooling. Moreover, a modification of the updating algorithm allows one to define a physical length scale beyond which the physics is largely unaltered. Closer to the focus of this investigation, some exploratory comments favoring the possibility of calibrating under-relaxed cooling were made in Ref. . We will lend further weight to this possibility in the following analysis.
This paper is divided as follows: Section II, briefly describes the gauge action and summarizes the numerical simulations examined in this study. Section III presents the relaxation algorithms and topological charge definition on the lattice. Sections IV A and IV B present the numerical results for the action and topological charge trajectories respectively. A summary of our findings is provided in Section V.
## II Lattice Gauge Action for $`SU_c(3)`$
### A Gauge Action.
The Wilson action is defined as,
$$S_G=\beta \underset{x}{}\underset{\mu <\nu }{}\frac{1}{3}e\mathrm{Tr}(1U_{\mu \nu }(x)),$$
(1)
where the operator $`U_{\mu \nu }(x)`$ is the standard plaquette,
$$U_{\mu \nu }(x)=U_\mu (x)U_\nu (x+\widehat{\mu })U_\mu ^{}(x+\widehat{\nu })U_\nu ^{}(x).$$
(2)
Gauge configurations are generated using the Cabbibo-Marinari pseudo–heat–bath algorithm with three diagonal $`SU_c(2)`$ subgroups. Simulations are performed using a parallel algorithm on a Thinking Machines Corporations (TMC) CM-5 with appropriate link partitioning. For the Wilson action we partition the link variables in the standard checkerboard fashion.
Configurations are generated on a $`16^3\times 32`$ lattice at $`\beta =5.70`$ and a $`24^3\times 36`$ lattice at $`\beta =6.00`$. Configurations are selected after 5000 thermalization sweeps from a cold start, and every 1000 sweeps thereafter. Lattice parameters are summarized in Table I. When required, we use the plaquette definition of the mean link variable.
## III Cooling and Smearing Algorithms
Relaxation techniques have been the subject of a large amount of study in lattice QCD over the past few years . Numerous algorithms have been employed to suppress the high-frequency components of the gauge field configurations in investigating their semi-classical content. In the following three subsections we consider three different algorithms that perform this task.
Perhaps it should be noted that cooling algorithms have been designed to render instantons stable over many hundreds of sweeps. For example, De Forcrand et al. have established a five-loop improved gauge action that achieves this. However, our interest here is in using the defects in the topological charge evolution under standard cooling and smearing algorithms based on the Wilson gauge action to calibrate these algorithms.
### A Cooling
Standard cooling minimizes the action locally at each link update. The preferred algorithm is based on the Cabbibo-Marinari pseudo-heat-bath algorithm for constructing $`SU_c(3)`$-color gauge configurations. A brief summary of a link update follows.
At the $`SU_c(2)`$ level the algorithm is transparent. An element of $`SU_c(2)`$ may be parameterized as, $`U=a_0I+i\stackrel{}{a}\stackrel{}{\sigma }`$, where $`a`$ is real and $`a^2=1`$. Let $`\stackrel{~}{U}_\mu `$ be one of the six staples associated with creating the plaquette associated with a link $`U_\mu `$.
$$\stackrel{~}{U}_\mu =U_\nu (x+\widehat{\mu })U_\mu ^{}(x+\widehat{\nu })U_\nu ^{}(x).$$
(3)
We define
$`{\displaystyle \underset{\alpha =1}{\overset{6}{}}}\stackrel{~}{U}_\alpha =k\overline{U}`$ $`\mathrm{where}`$ $`\overline{U}SU_c(2)\mathrm{and}k^2det\left({\displaystyle \underset{\alpha =1}{\overset{6}{}}}\stackrel{~}{U}_\alpha \right).`$ (4)
The feature of the sum of $`SU_c(2)`$ elements being proportional to an $`SU_c(2)`$ element is central to the algorithm. The local $`SU_c(2)`$ action is proportional to
$$e\mathrm{Tr}(1U\overline{U}),$$
(5)
and is locally minimized when $`e\mathrm{Tr}(U\overline{U})`$ is maximized, i.e. when
$$e\mathrm{Tr}(U\overline{U})=e\mathrm{Tr}(I).$$
(6)
which requires the link to be updated as
$$UU^{}=\overline{U}^1=\overline{U}^{}=\frac{\left(_{\alpha =1}^6\stackrel{~}{U_\alpha }\right)^{}}{k}.$$
(7)
At the $`SU_c(3)`$ level, one successively applies this algorithm to various $`SU_c(2)`$ subgroups of $`SU_c(3)`$, with $`SU_c(2)`$ subgroups selected to cover the $`SU_c(3)`$ gauge group.
We explored cooling gauge field configurations using two diagonal $`SU_c(2)`$ subgroups. The cooling rate is slow and the action density is not smooth, even after 50 sweeps over the lattice. Addition of the third diagonal $`SU_c(2)`$ subgroup provides acceptably fast cooling and a smooth action density. These $`SU_c(2)`$ subgroups are acting as a covering group for the $`SU_c(3)`$ gauge group. One would therefore expect that repeated updates of $`SU_c(2)`$ subgroups will better approximate the ultimate $`SU_c(3)`$ link.
### B Smearing
Here we consider two algorithms for smearing the gauge links. APE smearing is now well established. Our alternate algorithm AUS smearing is designed to remove instabilities of the APE algorithm and provide an intermediate algorithm sharing features of both smearing and cooling.
#### 1 APE Smearing
APE smearing is a gauge equivariant prescription for averaging a link $`U_\mu (x)`$ with its nearest neighbors $`U_\mu (x+\widehat{\nu }),\nu \mu `$. The linear combination takes the form:
$$U_\mu (x)U_\mu ^{}(x)=(1\alpha )U_\mu (x)+\frac{\alpha }{6}\stackrel{~}{\mathrm{\Sigma }}^{}(x;\mu ),$$
(8)
where $`\stackrel{~}{\mathrm{\Sigma }}(x;\mu )=\left(_{\nu =1}^6\stackrel{~}{U_\nu }(x)\right)_\mu `$, is the sum of the six staples defined in Eq. (3). The parameter $`\alpha `$ represents the smearing fraction. The algorithm is constructed as follows: (i) calculate the staples, $`\stackrel{~}{\mathrm{\Sigma }}(x;\mu )`$; (ii) calculate the new link variable $`U_\mu ^{}(x)`$ given by Eq. (8) and then reunitarize $`U_\mu ^{}(x)`$; (iii) once we have performed these steps for every link, the smeared $`U_\mu ^{}(x)`$ are mapped into the original $`U_\mu (x)`$. This defines a single APE smearing sweep which can then be repeated.
The celebrated feature of APE smearing is that it can be realized as higher-dimension operators that might appear in a fermion action for example. Indeed, “fat link” actions based on the APE algorithm are excellent candidates for an efficient action with improved chiral properties . The parameter space of fat link actions is described by the number of smearing sweeps $`n_{\mathrm{ape}}`$ and the smearing coefficient $`\alpha `$. It is the purpose of this investigation to explore the possible reduction of the dimension of this parameter space from two to one.
#### 2 AUS Smearing
The Annealed $`U`$ Smearing (AUS) algorithm is similar to APE smearing in that we take a linear combination of the original link and the associated staples, $`\stackrel{~}{\mathrm{\Sigma }}^{}(x;\mu )`$. However in AUS smearing we take what was a single APE smearing sweep over all the links on the lattice and divide it up into four partial sweeps based on the direction of the links. One partial sweep corresponds to an update of all links oriented in one of the four Cartesian directions denoted $`\mu =1,..,4`$. Hence in a partial sweep we calculate the reuniterized $`U_\mu ^{}(x)`$ for all $`\mu `$–oriented links and update all of these links $`[U_\mu ^{}(x)U_\mu (x)]`$ at the end of each partial sweep.
Thus the difference between AUS smearing and APE smearing is that in APE smearing no links are updated until all four partial sweeps are completed, while in AUS smearing, updated smeared information is cycled into the calculation of the next link direction, a process commonly referred to as annealing. In a sense AUS smearing is between cooling and APE smearing in that cooling updates one link at a time, AUS smearing updates one Cartesian direction at a time, and APE smearing updates the whole lattice at the same time.
As for APE smearing, the reunitarization of the $`SU_c(3)`$ matrix is done using the standard row by row orthonormalization procedure: begin by normalizing the first row; then update the second row by $`row2=row2(row2row1)row1`$; normalize the second row; finally set $`row3`$ equal to the cross product of $`row1`$ and $`row2`$.
For $`SU_c(2)`$ gauge theory, cooling and smearing are identical for the case of $`\alpha =1`$ when the links are updated one at a time. The AUS algorithm changes the level of annealing and provides the opportunity to alter the degree of smoothing. At the same time the AUS algorithm preserves the gauge equivariance of the APE algorithm for $`SU_c(3)`$ gauge theory . As we shall see, it also removes the instability of the APE algorithm at large smearing fractions $`\alpha `$. This latter feature may provide a more efficient method for finding the fundamental modular region of Landau Gauge in trying to understand the effects of Gribov copies in gauge dependent quantities such as the gluon-propagator .
### C Topological Charge
We construct the lattice topological charge density operator analogous to the standard Wilson action via the clover definition of $`F_{\mu \nu }`$. The topological charge $`Q_\mathrm{L}`$ and the topological charge density $`q_\mathrm{L}(x)`$ are defined as follows
$$Q_\mathrm{L}\underset{x}{}q_\mathrm{L}(x)=\frac{g^2}{32\pi ^2}ϵ_{\mu \nu \rho \sigma }\underset{x}{}\mathrm{Tr}(F_{\mu \nu }(x)F_{\rho \sigma }(x)).$$
(9)
where the field strength tensor is
$`a^2gF_{\mu \nu }(x)={\displaystyle \frac{i}{8}}\left[\left(𝒪_{\mu \nu }(x)𝒪_{\mu \nu }^{}(x)\right){\displaystyle \frac{1}{3}}\mathrm{Tr}\left(𝒪_{\mu \nu }(x)𝒪_{\mu \nu }^{}(x)\right)\right]+𝒪(a^4),`$ (10)
and $`𝒪_{\mu \nu }(x)`$ is
$`𝒪_{\mu \nu }(x)`$ $`=`$ $`U_\mu (x)U_\nu (x+\widehat{\mu })U_\mu ^{}(x+\widehat{\nu })U_\nu ^{}(x)`$ (11)
$`+`$ $`U_\nu (x)U_\mu ^{}(x+\widehat{\nu }\widehat{\mu })U_\nu ^{}(x\widehat{\mu })U_\mu (x\widehat{\mu })`$ (12)
$`+`$ $`U_\mu (x\widehat{\mu })U_\nu ^{}(x\widehat{\mu }\widehat{\nu })U_\mu (x\widehat{\mu }\widehat{\nu })U_\nu (x\widehat{\nu })`$ (13)
$`+`$ $`U_\nu ^{}(x\widehat{\nu })U_\mu (x\widehat{\nu })U_\nu (x+\widehat{\mu }\widehat{\nu })U_\mu ^{}(x).`$ (14)
Lattice operators possess a multiplicative lattice renormalization factor, $`Q_L=𝒵_Q(\beta )Q`$ which relates the lattice quantity $`Q_L`$ to the continuum quantity $`Q`$. Perturbative calculations indicate $`𝒵_Q(\beta )15.451/\beta +𝒪(1/\beta ^2)`$ . This large renormalization causes a problem when one is working at $`\beta 6.0`$, as $`𝒵_Q(\beta )1`$. This implies that the topological charge is almost impossible to calculate directly. However when one applies cooling or smearing techniques to remove the problem of the short range quantum fluctuations giving rise to $`𝒵_Q(\beta )1`$, one can resolve near integer topological charge. One can apply the operator $`Q_L`$ to cooled configurations or smeared configurations. The latter case may be regarded as employing an improved operator in which the smeared links are understood to give rise to additional higher dimension irrelevant operators designed to provide a smooth approach to the continuum limit.
## IV Numerical Simulations
We analyzed two sets of gauge field configurations generated using the Cabbibo-Marinari pseudo-heat-bath algorithm with three diagonal $`SU_c(2)`$ subgroups as described in Table I. Analysis of a few configurations proves to be sufficient to resolve the nature of the algorithms in question. The two sets are composed of ten $`16^3\times 32`$ configurations and five $`24^3\times 36`$ configurations. For each configuration we separately performed 200 sweeps of cooling, 200 sweeps of APE smearing at four values of the smearing fraction and 200 sweeps of AUS smearing at six values of the smearing fraction.
For clarity, we define the number of times an algorithm is applied to the entire lattice as $`n_c`$, $`n_{\mathrm{ape}}(\alpha )`$ and $`n_{\mathrm{aus}}(\alpha )`$ for cooling, APE smearing and AUS smearing respectively. We monitor both the total action normalized to the single instanton action $`S_0=8\pi ^2/g^2`$ and the topological charge $`Q_\mathrm{L}(x)`$ of Eq. (9) and observe their evolution as a function of the appropriate sweep variable and smearing fraction $`\alpha `$.
For APE smearing we consider four different values for the smearing fraction $`\alpha `$ including 0.30, 0.45, 0.55, and 0.70. Larger smearing fractions reveal an instability in the APE algorithm where the links are rendered to noise.
The origin of this instability is easily understood in fat-link perturbation theory where the smeared vector potential after $`n`$ APE smearing steps is given by
$$A_\mu ^{(n+1)}(q)=\underset{\nu }{}\left[f^n(q)\left(\delta _{\mu \nu }\frac{\widehat{q}_\mu \widehat{q}_\nu }{\widehat{q}^2}\right)+\frac{\widehat{q}_\mu \widehat{q}_\nu }{\widehat{q}^2}\right]A_\nu ^{(n)}(q),$$
(15)
reflecting the transverse nature of APE smearing. Here
$$\widehat{q}_\mu =\frac{2}{a}\mathrm{sin}\left(\frac{aq_\mu }{2}\right)$$
(16)
and
$$f(q)=1\frac{\alpha }{6}\widehat{q}^2.$$
(17)
For $`f(q)`$ to act as a form factor at each vertex in perturbation theory over the entire Brillouin zone
$$\frac{\pi }{a}<q_\mu \frac{\pi }{a},$$
(18)
one requires $`1f(q)1`$ which constrains $`\alpha `$ to the range $`0\alpha 3/4`$.
The annealing process in AUS smearing removes this instability. Hence, the parameter set for AUS smearing consists of the APE set plus an extra two, $`\alpha =0.85`$ and 1.00.
### A The Action Analysis.
#### 1 APE and AUS Smearing Calibration
The action normalized to the single instanton action $`S/S_0`$ can provide some insight into the number of instantons left in the lattice as a function of the sweep variable for each algorithm. However, our main concern is the relative rate at which the algorithms perform. In Fig. 1, we show $`S/S_0`$ as a function of cooling sweep on the $`24^3\times 36`$ lattice for five different configurations. The close proximity of the five curves is typical of the configuration dependence of the normalized action $`S/S_0`$.
Figures 2 and 3 report results for APE and AUS smearing respectively. Here we focus on one of the five configurations, noting that similar results are found for the other configurations. Each curve corresponds to a different value of the smearing fraction $`\alpha `$.
A similar analysis of the $`16^3\times 32`$ lattice at $`\beta =5.70`$ is also performed yielding analogous results. In fact, taking the physical volumes of the two lattices into account reveals qualitatively similar action densities after 200 sweeps.
Note that in Fig. 3, the curve associated with the smearing parameter $`\alpha =1.00`$, crosses over the one generated at $`\alpha =0.85`$, when the sweep number, $`n_{\mathrm{aus}}(\alpha )`$, is approximately 40 sweeps. Thus $`\alpha =1`$ is not the most efficient smearing fraction for untouched gauge configurations. However, as the configurations become smooth, $`\alpha =1.00`$ becomes the most efficient.
To calibrate the rate at which the algorithms reduce the action, we record the number of sweeps $`n(\alpha )`$ required for the smeared action to cross various thresholds $`S_T`$. This is repeated for each of the smearing fractions $`\alpha `$ under consideration. In establishing the relative $`\alpha `$ dependence for the number of sweeps $`n(\alpha )`$, we first consider a simple linear relation between the number of sweeps required to cross $`S_T`$ at one $`\alpha `$ compared to another $`\alpha ^{}`$, i.e.
$$n(\alpha ^{})=c_0+c_1n(\alpha ).$$
(19)
Anticipating that $`c_0`$ will be small if not zero, we divide both sides of this equation by $`n(\alpha )`$ and plot as a function of $`n(\alpha )`$. Deviations from a horizontal line will indicate failings of our linear assumption.
Fig. 4 displays results for $`\alpha ^{}`$ fixed to 0.55 for the APE smearing algorithm and Fig. 5 displays analogous results for AUS smearing. We omit thresholds that result in $`n(0.55)<10`$ as these points will have integer discretization errors exceeding 10%. For $`\alpha =0.70`$ discretization errors the order of 10% are clearly visible in both figures.
For the smearing fraction $`\alpha 0.85`$ both plots show little dependence of $`n(\alpha ^{}=0.55)/n(\alpha )`$ on $`n(\alpha )`$. This supports our simple ansatz of Eq. (19) and indicates $`c_0`$ is indeed small as one would expect.
The non-linear behavior for $`\alpha =1.00`$ in AUS smearing, in Fig. 5, reflects the cross over in Fig. 3, for $`\alpha =1.00`$. For unsmeared configurations, $`\alpha =1.00`$ appears to be too large. Further study in $`SU_c(2)`$ is required to resolve whether the origin of the smearing inefficiency lies in the sum of staples lying too far outside the $`SU_c(3)`$ gauge group for useful reunitarization, or whether the annealing of the links in AUS smearing is insufficient relative to cooling.
To determine the $`\alpha `$ dependence of $`c_1`$, we plot $`c_1=n(\alpha ^{}=0.55)/n(\alpha )`$ as a function of $`\alpha `$. Here the angular brackets denote averaging over all the threshold values considered; i.e. averaging over data in the horizontal lines of Figs. 4 and 5.
Fig. 6 for APE smearing and Fig. 7 for AUS smearing indicate the relationship between $`c_1`$ and $`\alpha `$ is linear with zero intercept to an excellent approximation. Indeed, ignoring the point at $`\alpha =1`$ for AUS smearing, one finds the same coefficients for the $`\alpha `$ dependence of APE and AUS smearing. When the fits are constrained to pass through the origin, one finds a slope of 1.818 which is the inverse of $`\alpha ^{}=0.55`$. Hence we reach the conclusion that
$$\frac{n_{\mathrm{ape}}(\alpha ^{})}{n_{\mathrm{ape}}(\alpha )}\frac{\alpha }{\alpha ^{}}\mathrm{and}\frac{n_{\mathrm{aus}}(\alpha ^{})}{n_{\mathrm{aus}}(\alpha )}\frac{\alpha }{\alpha ^{}}.$$
(20)
This analysis based on the action suggests that a preferred value for $`\alpha `$ does not really exist. In fact it has been recently suggested that one should anticipate some latitude in the values for $`n(\alpha )`$ and $`\alpha `$ that give rise to effective fat-link actions . What we have done here is established a relationship between $`n(\alpha )`$ and $`\alpha `$, thus reducing what was potentially a two dimensional parameter space to a one dimensional space. This conclusion will be further supported by the topological charge analysis below.
To summarize these finding we plot the ratios
$$\frac{\alpha ^{}n_{\mathrm{ape}}(\alpha ^{})}{\alpha n_{\mathrm{ape}}(\alpha )}\mathrm{and}\frac{\alpha ^{}n_{\mathrm{aus}}(\alpha ^{})}{\alpha n_{\mathrm{aus}}(\alpha )},$$
(21)
designed to equal 1 in Figs. 8 and 9 for APE and AUS smearing respectively. Figs. 10 and 11 report the final results of a similar analysis for APE and AUS smearing respectively at $`\beta =5.7`$. Here the $`\alpha =1.0`$ results are omitted from the AUS smearing results for clarity.
#### 2 Cooling Calibration
Here we repeat the previous analysis, this time comparing cooling with APE and AUS smearing. We make the same linear ansatz for the relationship and plot $`n_c/n_{\mathrm{ape}}(\alpha )`$ versus $`n_{\mathrm{ape}}(\alpha )`$ for APE smearing in Fig. 12. Fig. 13 reports the ratio $`n_c/n_{\mathrm{aus}}(\alpha )`$ versus $`n_{\mathrm{aus}}(\alpha )`$ for AUS smearing. At small numbers of smearing sweeps, large integer discretization errors the order of 25% are present, as $`n_c`$ is as small as 4. With this in mind, we see an independence of the ratio on the amount of cooling/smearing over a wide range of smearing sweeps for both APE and AUS smearing. This supports a linear relation between the two algorithms.
Averaging the results provides
$$\frac{n_c}{n_{\mathrm{ape}}(0.55)}=0.330\mathrm{and}\frac{n_c}{n_{\mathrm{aus}}(0.55)}=0.340,$$
(22)
or more generally
$$n_c0.600\alpha n_{\mathrm{ape}}(\alpha )\mathrm{and}n_c0.618\alpha n_{\mathrm{aus}}(\alpha ).$$
(23)
Hence we see that cooling is much more efficient at smoothing than the smearing algorithms requiring roughly half the number of sweeps for a given product of $`n`$ and $`\alpha `$.
Equating the equations of Eq. (23) provides the following relation between APE and AUS smearing
$$\alpha n_{\mathrm{ape}}(\alpha )1.03\alpha ^{}n_{\mathrm{aus}}(\alpha ^{}),$$
(24)
summarizing the near equivalence of the two smearing algorithms. It is important to recall that the annealing of the AUS smearing algorithm removes the instability of APE smearing encountered at large smearing fractions, $`\alpha `$. Thus AUS smearing offers a more stable gauge equivariant smoothing of gauge fields. It offers faster smoothing due to its ability to handle larger smearing fractions. This algorithm may be of use in studying Gribov ambiguities in Landau gauge fixing .
A similar analysis of the $`16^3\times 32`$ lattice at $`\beta =5.7`$ provides
$$n_c0.572\alpha n_{\mathrm{ape}}(\alpha )\mathrm{and}n_c0.604\alpha n_{\mathrm{aus}}(\alpha ),$$
(25)
with
$$\alpha n_{\mathrm{ape}}(\alpha )1.06\alpha ^{}n_{\mathrm{aus}}(\alpha ^{}).$$
(26)
Here the change in the coefficient relating APE smearing to cooling appears to be proportional to $`\beta `$.
### B The Topological Charge Density Analysis.
Typical evolution curves for the lattice topological charge operator of Eq. (9) are shown in Figs. 14, 15 and 16 for APE smearing, AUS smearing and cooling respectively. These data are obtained from a typical gauge configuration on our $`24^3\times 36`$ lattice at $`\beta =6.00`$. The configuration used here is the same representative configuration illustrated in Figs. 2 and 3 of the action analysis.
The main feature of these figures is that the smearing/cooling algorithms produce a similar trajectory for the topological charge. For most cases only the rate at which the trajectory evolves changes. In these cases, one can use these trajectories as another way in which to calibrate the rates of the algorithms.
After a few sweeps, the topological charge converges to a near integer value with an error of about 10%, typical of clover definitions of $`Q`$ at $`\beta =6.0`$. The standard Wilson action is known to lose (anti)instantons during cooling due to $`𝒪(a^2)`$ errors in the action which act to tunnel through the single instanton action bound. Given the intimate relationship between cooling and smearing discussed in Section III it is not surprising to see similar behavior in the topological charge trajectories of the smearing algorithms considered here. The sharp transitions from one integer to another indicate the loss of an (anti)instanton.
These curves look quite different for other configurations. Fig. 17 displays trajectories for cooling on five different gauge configurations. However, the feature of similar trajectories for the various cooling/smearing algorithms remains at $`\beta =6.0`$. To calibrate the cooling/smearing algorithms, we select thresholds at topological charge values where the trajectories are making sharp transitions from one near integer to another. Table II summarizes the thresholds selected and the number of sweeps required to pass though the various thresholds. Note that the second configuration, $`C_2`$, data entry in Table II corresponds to the sampling of the curves illustrated in Figs. 14, 15 and 16.
At $`\beta =5.7`$ we did not find analogous trajectories, suggesting that the coarse lattice spacing and larger errors in the action prevent one from reproducing a similar smoothed gauge configuration using different algorithms. In this case the two-dimensional aspect of the smearing parameter space remains for studies of the topological sector. That this might be the case is hinted at in Fig. 11 in which the ratio of smearing results for $`\beta =5.7`$ lattices is not as closely constrained to one as for the $`\beta =6.0`$ results in Figs. 9.
As in the action analysis, we report the ratio of $`\alpha =0.55`$ results to other $`\alpha `$ value results within APE and AUS smearing. Figs. 18 and 19 summarize the results for APE and AUS smearing respectively. Again we see an enhanced spread of points at small numbers of smearing sweeps due to integer discretization errors.
Data for large numbers of sweeps at small and large $`\alpha `$ values are absent in Fig. 19. The absence of points for small $`\alpha `$ values is simply due to the thresholds not being crossed within 200 smearing sweeps. The absence of points for large $`\alpha `$ values reflects the divergence of the topological charge evolution from the most common trajectory among the algorithms. This divergence is also apparent in Fig. 5 where the $`\alpha =1.0`$ AUS smearing results fail to satisfy the linear ansatz.
While the data from the topological charge evolution is much more sparse, one can see reasonable horizontal bands forming supporting a dominant linear relationship between various $`\alpha `$ values. Averages of the bands are reported in Table III along with previous results from the action analysis. With the exception of the $`\alpha =1.0`$ AUS smearing results, the agreement is remarkable, leading to the same conclusions of the action analysis summarized in Eq. (20).
In calibrating cooling via the topological charge defects, we report the ratio $`n_c/n_{\mathrm{ape}}(0.55)`$ as a function of $`n_{\mathrm{ape}}(\alpha )`$ and $`n_c/n_{\mathrm{aus}}(0.55)`$ as a function of $`n_{\mathrm{aus}}(\alpha )`$ for APE and AUS smearing respectively. The results are shown in Figs. 20 and 21.
The data is too poor to determine anything beyond a linear relation between $`n_c`$ and $`n_{\mathrm{ape}}(\alpha )`$ or $`n_{\mathrm{aus}}(\alpha )`$. Averaging the results provides
$$\frac{n_c}{n_{\mathrm{ape}}(0.55)}=0.30(3)\mathrm{and}\frac{n_c}{n_{\mathrm{aus}}(0.55)}=0.35(3),$$
(27)
in agreement with the earlier action based results of Eq. (22).
As a final examination of the relations we have established among APE smearing, AUS smearing and cooling algorithms, we illustrate the topological charge density in Fig. 22. Large positive (negative) winding densities are shaded red (blue) and symmetric isosurfaces aid in rendering the shapes of the densities. In fixing the $`x`$ coordinate to a constant, a three-dimensional slice of a four-dimensional $`24^3\times 36`$ gauge field configuration is displayed. Fig. 22(a) illustrates the topological charge density after 21 APE smearing steps at $`\alpha =0.7`$. The other three quadrants display APE smearing, AUS smearing and cooling designed to reproduce Fig. 22(a) according to the relations of Eqs. (20), (23) and (24). The level of detail in the agreement is remarkable.
## V Conclusion
The APE smearing algorithm is now widely used in a variety of ways in lattice simulations. It is used to smear the spatial gauge-field links in studies of glueballs, hybrid mesons, the static quark potential, etc. It is used in constructing fat-link actions and in constructing improved operators with smooth transitions to the continuum limit. We have shown that to a good approximation the two-dimensionful parameter space of the number of smearing sweeps $`n_{\mathrm{ape}}(\alpha )`$ and the smearing fraction $`\alpha `$ may be reduced to a single dimension via the constraint
$$\frac{n_{\mathrm{ape}}(\alpha ^{})}{n_{\mathrm{ape}}(\alpha )}=\frac{\alpha }{\alpha ^{}}.$$
(28)
satisfied for $`\alpha `$ and $`\alpha ^{}`$ in the range 0.3 to 0.7. This result is in agreement with fat-link perturbation theory expectations, and survives for up to 200 sweeps over the lattice. We expect this relation to hold over the entire APE smearing range $`0<\alpha <3/4`$. We find the same relation for AUS smearing provided $`\alpha 0.85`$. For AUS smearing, annealing of the links is included as one cycles through the Lorentz directions of the link variables in the smearing process.
The smoothing of gauge field configurations is often a necessary step in extracting observables in which the renormalization constants differ significantly from one, as is the case for the topological charge operator. It is also often used to gain insights into the nonperturbative features of the field theory which give rise to the observed phenomena. We have determined that cooling, APE smearing and AUS smearing produce qualitatively similar smoothed gauge field configurations at $`\beta =6.0`$ provided one calibrates the algorithms as follows:
$$n_c0.600\alpha n_{\mathrm{ape}}(\alpha ),n_c0.618\alpha n_{\mathrm{aus}}(\alpha )\mathrm{and}\alpha n_{\mathrm{ape}}(\alpha )1.03\alpha ^{}n_{\mathrm{aus}}(\alpha ^{}).$$
(29)
The topological charge analysis serves to confirm the action analysis results at $`\beta =6.0`$ and further support these relations. However, it also reveals that at $`\beta =5.7`$ different trajectories are taken. At $`\beta =5.7`$ different algorithms produce smoothed gauge field configurations with similar action, but different topological properties.
As most modern simulations are performed at $`\beta 6.0`$, the relations of Eqs. (28) and (29) will be most effective in reducing the exploration of the parameter space. It is now possible to arrive at optimal smearing by fixing the number of smearing sweeps and varying the smearing fraction, or vice versa. Finally, it is now possible to directly compare the physics of smeared and cooled gauge field configurations in a quantitative sense.
We have used the action and topological charge trajectories to calibrate the various smoothing algorithms. While the topological charge may be related to the topological susceptibility, it may be of future interest to introduce other physical quantities such as the string tension as a measure of smoothing .
It is well known that interpreting the physics of cooled configurations based on the standard Wilson action is somewhat of an art. The difficulty is that the $`𝒪(a^2)`$ errors in the action spoil an instanton under cooling by reducing the action below the one-instanton bound. A consequence of this is a lack of universality. For example, the action varies smoothly to zero while the topological sector jumps as (anti)instantons are destroyed. One may then attempt to define the concept of an optimal amount of smoothing, and perhaps also extrapolate back to zero smoothing steps. Unfortunately, all of these difficulties are apparent in the APE smearing analysis results as well.
Because of the intimate connection between smearing and cooling, as emphasized in the discussion of Section III B 2, it is natural to consider improved smearing, where the Hermitian conjugate of the improved staple is used in the smearing process. It will be interesting to see if the benefits of improved cooling carry over into improved smearing and this work is currently in progress.
## Acknowledgment
Thanks are extended to Tom Degrand for helpful comments on fat-link perturbation theory. Thanks also to Francis Vaughan of the South Australian Centre for Parallel Computing and the Distributed High-Performance Computing Group for support in the development of parallel algorithms implemented in Connection Machine Fortran (CMF) and for generous allocations of time on the University of Adelaide’s CM-5. Support for this research from the Australian Research Council is gratefully acknowledged. AGW also acknowledges support form the Department of Energy Contract No. DE-FG05-86ER40273 and by the Florida State University Supercomputer Computations Research Institute which is partially funded by the Department of Energy through contract No. DE-FC05-85ER25000. |
warning/0001/cond-mat0001058.html | ar5iv | text | # Dimensional Crossover in Quasi One-Dimensional and High 𝑇_𝑐 Superconductors
## I Abelian Bosonization and the Spectral Functions
We begin by considering the properties of a single chain in the absence of any interchain coupling; we treat this problem using Abelian bosonization, which is based on the fact that the properties of an interacting 1DEG at low energies and long wavelength are asymptotically equal to those of a set of two independent bosonic fields, one representing the charge and the other the spin degrees of freedom in the system. The widely discussed separation of charge and spin in this problem is formally the statement that the Hamiltonian density $`_j`$ can be expressed as
$$=_c+_s,$$
(2)
where the chain index is implicit, and the charge and spin pieces of the Hamiltonian are each of the sine Gordon variety,
$$_\alpha =\frac{v_\alpha }{2}\left[K_\alpha (_x\theta _\alpha )^2+\frac{(_x\varphi _\alpha )^2}{K_\alpha }\right]+V_\alpha \mathrm{cos}(\sqrt{8\pi }\varphi _\alpha ),$$
(3)
where $`\alpha =c,s`$ for the charge and spin fields, respectively, $`\theta _\alpha `$ is the dual field to $`\varphi _\alpha `$, or equivalently $`_x\theta _\alpha `$ is the momentum conjugate to $`\varphi _\alpha `$. We consider a sufficiently incommensurate 1DEG and therefor set $`V_c=0`$ since it arises from Umklapp scattering. Of course, if the Umklapp scattering is crucial to explain doped insulator behavior, its role cannot be neglected. Where there is no spin gap, or at temperatures large compared to $`\mathrm{\Delta }_s`$, we can likewise set $`V_s=0`$.
When $`V_s`$ is relevant, (perturbatively, this means $`K_s<1`$) the spin gap is dynamically generated, i.e. it depends both on $`V_s`$ and the ultraviolet cutoff in the problem, $`\mathrm{\Lambda }`$, according to the scaling relation $`\mathrm{\Delta }_sv_s\mathrm{\Lambda }[V_s/v_s\mathrm{\Lambda }^2]^{1/(22K_s)}`$. At the gapless fixed point, spin-rotational invariance requires $`K_s=1`$, at which point $`V_s`$ is perturbatively marginal. It is marginally irrelevant for repulsive interactions ($`K_s>1`$) and marginally relevant for attractive interactions ($`K_s<1`$). Thus, the long distance spin physics is described by $`_s`$ with $`V_s=0`$ and $`K_s=1`$ for a gapless spin-rotationally invariant phase. Where there is a spin gap in a spin rotationally invariant system, it is exponentially small for weak interactions, $`\mathrm{\Delta }_s\sqrt{V_sv_s}\mathrm{exp}[v_s\mathrm{\Lambda }^2/2\pi V_s]`$.
In order to compute correlation functions, we use the Mandelstam representation of the fermion field operators
$$\psi _{\lambda ,\sigma }(x)=𝒩_\sigma \mathrm{exp}\left[i\lambda k_Fxi\mathrm{\Phi }_{\lambda ,\sigma }(x)\right],$$
(4)
where $`𝒩_\sigma `$ contains both a normalization factor (which depends on the ultraviolet cutoff) and a “Klein” factor (which can be implemented in many ways) so that $`𝒩_\sigma `$ anticommutes with $`𝒩_\sigma ^{}`$ for $`\sigma \sigma ^{}`$ and commutes with it for $`\sigma =\sigma ^{}`$. In addition
$$\mathrm{\Phi }_{\lambda ,\sigma }=\sqrt{\pi /2}\left[(\theta _c\lambda \varphi _c)+\sigma (\theta _s\lambda \varphi _s)\right],$$
(5)
where $`\lambda =1`$ for left moving electrons, $`\lambda =+1`$ for right moving electrons, and $`\sigma =\pm 1`$ refers to spin polarization. From Eq. (4), it is a straightforward (and standard) exercise to obtain the boson representations of all interesting electron bilinear and quartic operators. Physically, $`\varphi _c`$ and $`\varphi _s`$ are, respectively, the phases of the $`2k_F`$ CDW and SDW fluctuations, and $`\theta _c`$ is the superconducting phase. The long-wavelength components of the charge ($`\rho `$) and spin ($`S_z`$) densities are given by
$`\rho (x)=`$ $`{\displaystyle \underset{\lambda ,\sigma }{}}\psi _{\lambda ,\sigma }^{}\psi _{\lambda ,\sigma }{\displaystyle \frac{2k_F}{\pi }}=\sqrt{{\displaystyle \frac{2}{\pi }}}_x\varphi _c,`$ (6)
$`S_z(x)=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\lambda ,\sigma }{}}\sigma \psi _{\lambda ,\sigma }^{}\psi _{\lambda ,\sigma }=\sqrt{{\displaystyle \frac{1}{2\pi }}}_x\varphi _s.`$ (7)
When analyzing results for this model, it is always important to remember that the parameters which enter the field theory are renormalized, and are related to the microscopic interactions in a very complicated manner. For instance, although for a single-component 1DEG with repulsive interactions $`V_s`$ is always irrelevant, for multicomponent 1DEG’s, and for the “1DEG in an active environment”, it is common to find a dynamically generated spin gap, even when the microscopic interactions are uniformly repulsive.
The bosonized expressions for all electron operators are readily extended to an array of chains by adding a chain index to the Bose fields and to the Klein factors; the Klein factors on different chains must now anticommute with each other. Where single particle interchain hopping is relevant, the Klein factors appear explicitly in the bosonized Hamiltonian. Where only pair hopping and collective interactions between neighboring chains need be included in the low energy physics, the Klein factors cancel in $`H`$.
While it is generally simpler to derive results concerning the spectrum, it is important for comparison with experiment to compute actual correlation functions. Specifically, we will consider the transverse spin dynamic structure factor
$`\stackrel{~}{𝒮}(x,t;T)`$ $``$ $`\mathrm{S}_{2k_F}^{x}{}_{}{}^{}(x,t)\mathrm{S}_{2k_F}^x(0,0)`$ (8)
$`+`$ $`\mathrm{S}_{2k_F}^{y}{}_{}{}^{}(x,t)\mathrm{S}_{2k_F}^y(0,0),`$ (9)
where
$$𝐒_{2k_F}=\frac{1}{2}\underset{\sigma ,\sigma ^{}}{}\psi _{1,\sigma }^{}𝝉_{\sigma \sigma ^{}}\psi _{1,\sigma ^{}},$$
(10)
and the $`𝝉`$ are Pauli matrices. We will also consider the one-hole Green function,
$`\stackrel{~}{G}^<(x,t)`$ $`\psi _{1,}^{}(x,t)\psi _{1,}(0,0),`$ (11)
the singlet-pair correlator,
$$\stackrel{~}{\chi }(x,t)\psi _{1,}^{}(x,t)\psi _{1,}^{}(x,t)\psi _{1,}(0,0)\psi _{1,}(0,0),$$
(12)
and the various spectral functions, $`𝒮`$, $`G^<`$, and $`\chi `$, obtained by Fourier transforming these correlators. As a consequence of separation of charge and spin, $`\stackrel{~}{𝒮}`$, $`\stackrel{~}{G}^<`$, and $`\stackrel{~}{\chi }`$ are expressible as a product of spin and charge contributions, and consequently $`𝒮`$, $`G^<`$, and $`\chi `$ are convolutions. For instance,
$$G^<(k,\omega )=\frac{dq}{2\pi }\frac{d\nu }{2\pi }G_s(kq,\omega \nu )G_c(q,\nu ).$$
(13)
## II High temperature: Luttinger liquid behavior
At temperatures large compared to $`T_c`$ and the spin gap, $`\mathrm{\Delta }_s`$, (or at all temperatures in systems in which $`T_c=\mathrm{\Delta }_s=0`$), the 1DEG exhibits “Luttinger liquid” behavior. Because the Luttinger liquid is a quantum critical system, the response functions have a scaling form. Specifically, this implies that
$$G^<(k,\omega ;T)=T^{2\gamma _c+2\gamma _s1}G^<(k/T,\omega /T;1),$$
(14)
where we define for $`\alpha =c`$ or $`s`$
$$\gamma _\alpha =\frac{1}{8}(K_\alpha +K_\alpha ^12),$$
(15)
and that so long as $`K_s1`$,
$$𝒮(k,\omega ;T)=T^{(K_s^1+K_c2)}𝒮(k/T,\omega /T;1).$$
(16)
Note that here, and henceforth, we will measure $`k`$ relative to $`k_F`$ and $`2k_F`$, respectively, when computing the scaling functions $`G^<`$ and $`𝒮`$. If the system is spin-rotationally invariant, $`K_s=1`$ in the above expressions.
The form of these scaling functions can be computed analytically in many cases; this has recently been accomplished in Ref. C. They may or may not have a peak at energies small compared to the bandwidth, depending on certain exponent inequalities. Where there is a peak, it occurs at positive energies $`\omega =\pm v_\alpha k+(\mathrm{const}.)T`$, but the peak width, however defined, does not narrow in proportion to $`T`$ at low temperatures; such a peak does not correspond to a quasiparticle.
## III Intermediate temperature: The Luther-Emery liquid
When $`V_s`$ is relevant, the spin sine Gordon theory scales to a strong-coupling fixed point, and the excitations are massive solitons, in which $`\varphi _s`$ changes by $`\pm \sqrt{\pi /2}`$ (i.e. $`S_z=\pm 1/2`$). This problem is most simply treated in terms of spin fermion fields,
$$\mathrm{\Psi }_{s,\lambda }^{}=𝒩_s\mathrm{exp}[i\sqrt{\pi /2}(\theta _s2\lambda \varphi _s)].$$
(17)
The refermionized form of the Hamiltonian is then
$`_s=`$ $`i\stackrel{~}{v}_s[\mathrm{\Psi }_{s,1}^{}_x\mathrm{\Psi }_{s,1}\mathrm{\Psi }_{s,1}^{}_x\mathrm{\Psi }_{s,1}]`$ (20)
$`+\stackrel{~}{\mathrm{\Delta }}_s[\mathrm{\Psi }_{s,1}^{}\mathrm{\Psi }_{s,1}+\mathrm{H}.\mathrm{c}.]`$
$`+g_s\mathrm{\Psi }_{s,1}^{}\mathrm{\Psi }_{s,1}^{}\mathrm{\Psi }_{s,1}\mathrm{\Psi }_{s,1},`$
where
$`\stackrel{~}{v}_s`$ $`=`$ $`v_s\left({\displaystyle \frac{1}{4K_s}}+K_s\right),`$ (21)
$`\stackrel{~}{\mathrm{\Delta }}_s`$ $`=`$ $`{\displaystyle \frac{\pi V_s}{\mathrm{\Lambda }}},`$ (22)
$`g_s`$ $`=`$ $`2\pi v_s\left({\displaystyle \frac{1}{4K_s}}K_s\right).`$ (23)
For $`K_s=1/2`$, which is known as the Luther-Emery point, the refermionized model is non-interacting and massive, with a gap $`\mathrm{\Delta }_s=\stackrel{~}{\mathrm{\Delta }}_s`$. Assuming there is a single massive phase of the sine Gordon theory, the Luther-Emery model will exhibit the same asymptotic behavior as any other model in this phase. Formally, the Luther-Emery point can be thought of as a strong-coupling fixed point Hamiltonian, and $`g_s`$, which vanishes at the fixed point, is the amplitude of a leading irrelevant operator. We will henceforth compute correlation functions at the Luther-Emery point, and then comment on the effects of deviations from this point.
Now, in computing the various spectral properties of the system, we can distinguish two regimes of temperature: at temperatures large compared to $`\mathrm{\Delta }_s`$, the spin gap is negligible, and the results for the Luttinger liquid apply. If the temperature is small compared to the spin gap, then we can compute the spin contributions to the various correlation functions in the zero-temperature limit, and only make exponentially small errors of order $`\mathrm{exp}(\mathrm{\Delta }_s/T)`$. The spin piece of the transverse spin response function can be expressed in terms of the spin fermion fields
$$\stackrel{~}{𝒮}_s(x,t)=\mathrm{\Psi }_{s,1}^{}(x,t)\mathrm{\Psi }_{s,1}^{}(x,t)\mathrm{\Psi }_{s,1}(0,0)\mathrm{\Psi }_{s,1}(0,0).$$
(24)
Since the theory reduces, at the Luther-Emery point, to a theory of free massive fermions, the corresponding spectral function can be readily computed with the result, for $`T=0`$,
$$𝒮_s(k,\omega )=\frac{\omega ^24E_s^2(k/2)}{4v_s^2|q_1E_s(q_2)q_2E_s(q_1)|}\mathrm{\Theta }[\omega 2E_s(k/2)],$$
(25)
where the spin soliton spectrum is
$$E_s(k)=\sqrt{v_s^2k^2+\mathrm{\Delta }_s^2},$$
(26)
and $`q_{1,2}`$ are the solutions of the quadratic equation $`\omega +E_s(q)+E_s(k+q)=0`$. Explicitly
$$q_{1,2}=\frac{k}{2}\pm \frac{\omega }{2v_s}\sqrt{1+\frac{4\mathrm{\Delta }_s^2}{v_s^2k^2\omega ^2}}.$$
(27)
The spin piece of the one hole Green function is more complicated, since it involves nonlocal operators in the refermionized form:
$$\stackrel{~}{G}_s(x,t)=U_s^{}(x,t)\mathrm{\Psi }_{s,1}^{}(x,t)\mathrm{\Psi }_{s,1}(0,0)U_s(0,0),$$
(28)
where the vertex operator $`U_s(x)=e^{i\sqrt{\pi /2}\varphi _s(x)}`$ with $`\varphi _s(x)=\sqrt{\pi /2}_\lambda ^x𝑑y\mathrm{\Psi }_{s,\lambda }^{}\mathrm{\Psi }_{s,\lambda }`$. From kinematics, it follows that this Green function consists of a coherent one spin soliton piece and an incoherent multisoliton piece:
$`G_s(k,\omega )=Z_s(k)\delta [\omega E_s(k)]+G_s^{(multi)}(k,\omega ),`$ (29)
where the multisoliton piece is proportional to $`\mathrm{\Theta }[\omega 3E_s(k/3)]`$. (Deviations from the Luther-Emery point in the case $`g_s>0`$ will result in the formation of a spin soliton-antisoliton bound state, a “breather”, which can shift the threshold energy for multisoliton excitations somewhat.)
At the Luther-Emery point it is possible to obtain closed form expressions for the matrix elements of the vertex operator between the vacuum and various multisoliton states and from that to compute $`Z_s`$ explicitly. We will report this calculation in a forthcoming paper, Ref. C. Here, we use a simple scaling argument, which can be generalized to the case of nonzero interchain coupling, to derive the principal features of this result, especially the dependence of $`Z_s`$ on $`\mathrm{\Delta }_s`$. In the absence of a spin gap, and at $`T=0`$, $`G_s`$ can be readily evaluated to give the scaling form
$`G_s={\displaystyle \frac{\pi (v_s\mathrm{\Lambda })^{\frac{1}{2}2\gamma _s}}{\mathrm{\Gamma }(\gamma _s)\mathrm{\Gamma }(\gamma _s+\frac{1}{2})}}(\omega +v_sk)^{\gamma _s1}`$ $`(\omega v_sk)^{\gamma _s\frac{1}{2}}`$ (31)
$`\times \mathrm{\Theta }(\omega v_s|k|).`$
Because the sine-Gordon field theory is asymptotically free, the high energy spectrum, and hence the dependence of $`G_s`$ on $`\mathrm{\Lambda }`$, is unaffected by the opening of a spin gap. With this observation, it is simply a matter of dimensional analysis to see that
$$Z_s(k)=(\mathrm{\Lambda }\xi _s)^{\frac{1}{2}2\gamma _s}f_s(k\xi _s),$$
(32)
where $`\xi _s=v_s/\mathrm{\Delta }_s`$ is the spin correlation length. $`f_s`$ is a scaling function which is independent of $`K_s`$. It can be calculated using the exact matrix elements available for $`K_s=1/2`$, with the result
$$f_s(x)=c\left(1\frac{x}{\sqrt{1+x^2}}\right),$$
(33)
where $`c`$ is a numerical constant.
The above extends the earlier results of Voit and Wiegmann. In particular, the analytic structure (as a function of $`k`$ and $`\omega `$) of the one soliton contribution to Eq. (29) reproduces that found in earlier work. The explicit discussion of the nonanalyticities due at the three soliton threshold is new, although fairly obvious; more muted singularities occur at the five and higher multisoliton thresholds, which we will not discuss explicitly. The specific expression in Eq. (32) is, to the best of our knowledge, new, and is the most important feature of this result for the purposes of the present paper.
The charge pieces of both response functions are unaffected by the opening of the spin gap. Consequently, $`𝒮`$ and $`G^<`$ have power law features (which can be a peak or a shoulder depending on $`K_c`$) at $`\omega =2E_s(k/2)+𝒪(T)`$ and $`\omega =E_s(k)+𝒪(T)`$, respectively, with a shape and temperature dependence, both readily computed, determined by the still gapless charge density fluctuations. For example, we can evaluate the spectral function explicitly at $`T=0`$ in the limit $`v_s/v_c0`$ (and for arbitrary $`\omega <3E_s(k/3)`$), or when $`|\omega \mathrm{\Delta }_s|\mathrm{\Delta }_s`$ (for arbitrary $`v_s/v_c`$ ):
$`G^<(k,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \frac{B(\gamma _c,\gamma _c+\frac{1}{2})}{\mathrm{\Gamma }(\gamma _c)\mathrm{\Gamma }(\gamma _c+\frac{1}{2})}}\left({\displaystyle \frac{\mathrm{\Lambda }v_c}{2}}\right)^{\frac{1}{2}2\gamma _c}`$ (34)
$`\times `$ $`Z_s(k)[\omega E_s(k)]^{2\gamma _c\frac{1}{2}}\mathrm{\Theta }[\omega E_s(k)].`$ (35)
Here $`B(x,y)`$ is the beta function. Again, the fact that these excitations are not quasiparticles is reflected in the fact that, even where peaks in the spectral function occur, they do not narrow indefinitely as $`T0`$.
In the presence of a spin gap, the spin contribution to the long distance behavior of the superconducting susceptibility is a constant;
$`\stackrel{~}{\chi }(x,t)`$ $``$ $`|U_s^2|^2e^{i\sqrt{2\pi }\theta _c(x,t)}e^{i\sqrt{2\pi }\theta _c(0,0)}`$ (36)
$``$ $`(\mathrm{\Lambda }\xi _s)^{K_s}e^{i\sqrt{2\pi }\theta _c(x,t)}e^{i\sqrt{2\pi }\theta _c(0,0)}.`$ (37)
From this, one sees that, within a chain, one can identify $`(\mathrm{\Lambda }\xi _s)^{K_s/2}`$ as the “amplitude” and $`\sqrt{2\pi }\theta _c`$ as the “phase” of the order parameter.
## IV Low Temperature: The 3d Superconducting State
For temperatures of order $`T_c`$ and below, interchain couplings cannot be ignored. Single particle hopping and all magnetic couplings are irrelevant by virtue of the preexisting spin gap. For $`K_c>1/2`$, the Josephson coupling is perturbatively relevant, but for $`K_c<1`$, the $`2k_F`$ CDW coupling is more relevant. For the simplest realizations of the 1DEG, $`K_c<1`$ corresponds to repulsive interactions between charges. However, we have shown that for fluctuating or meandering stripes, such as occur in the high temperature superconductors, the CDW coupling gets dephased, so that the Josephson coupling is the most relevant, even when $`1/2<K_c<1`$.
Since we are interested in the onset of superconductivity, we consider the case in which the Josephson coupling between chains is more relevant. The pair tunnelling interaction between chains, which appeared in Eq. (1), can be simply bosonized:
$$H_J=J_{SC}\underset{<i,j>}{}dx[\widehat{\mathrm{\Delta }}_i^{}\widehat{\mathrm{\Delta }}_j+\mathrm{H}.\mathrm{C}.],$$
(38)
where the pair creation operator on chain number $`j`$ is
$`\widehat{\mathrm{\Delta }}^{}(x,t)`$ $`=`$ $`\psi _{1,}^{}\psi _{1,}^{}+\psi _{1,}^{}\psi _{1,}^{}`$ (39)
$``$ $`\mathrm{cos}(\sqrt{2\pi }\varphi _s)\mathrm{exp}(i\sqrt{2\pi }\theta _c),`$ (40)
and we have left the chain index implicit.
Since the state below $`T_c`$ has long range order, and since we assume that the coupling between chains is weak, it is reasonable to treat it in mean field approximation, although we continue to treat the one-dimensional fluctuations exactly. Thus, rather than considering a full three-dimensional problem, we consider the effective single chain problem defined by the Hamiltonian
$$=_s+_c𝒥\mathrm{cos}(\sqrt{2\pi }\varphi _s)\mathrm{cos}(\sqrt{2\pi }\theta _c),$$
(41)
where $`𝒥`$ is related to the pair tunnelling amplitude by the mean field relation
$$𝒥=zJ_{SC}(\mathrm{\Lambda }/\pi )^2\mathrm{cos}(\sqrt{2\pi }\varphi _s)\mathrm{cos}(\sqrt{2\pi }\theta _c),$$
(42)
where $`z`$ is the number of nearest neighbor chains. \[Since the average of $`\mathrm{cos}(\sqrt{2\pi }\varphi _s)\mathrm{sin}(\sqrt{2\pi }\theta _c)`$ vanishes, no sine term appears in the effective Hamiltonian (41).\] Note that the pair hopping term in Eq. (41) couples charge and spin, as is characteristic of higher dimensional couplings.
The mean field approximation is exact in the limit of large $`z`$ and small $`zJ_{SC}`$. In three dimensions, this mean field approximation will produce some errors in the critical regime in the vicinity of $`T_c`$, but because of the long correlation length along the chain just above $`T_c`$, the critical region is always small for small $`J_{SC}`$, and well below $`T_c`$, this approximation is safe.
Because of the presence of relevant cosine terms, there are superselection rules which divide Hilbert space into various soliton sectors. The soliton sectors are specified by two integrals:
$`N_s=`$ $`\sqrt{2/\pi }{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑x_x\varphi _s`$ (43)
$`=`$ $`\sqrt{2/\pi }[\varphi _s(\mathrm{})\varphi _s(\mathrm{})]=2{\displaystyle 𝑑xS_z},`$ (44)
and
$$N_c=\sqrt{2/\pi }_{\mathrm{}}^{\mathrm{}}𝑑x_x\theta _c=\sqrt{2/\pi }[\theta _c(\mathrm{})\theta _c(\mathrm{})].$$
(45)
$`N_s`$ is simply the number of spin solitons minus the number of antisolitons or the total value of $`S_z`$ in units of $`\mathrm{}/2`$. The interpretation of $`N_c`$ is a bit more subtle. Since we are looking at a superconducting state, the electrostatic charge of a quasiparticle is not defined. However, $`N_c`$ is a conserved “chirality” equal to the number of right moving minus the number of left moving electrons, so that we can still interpret $`eN_c`$ as a sort of quasiparticle “charge”; it represents the coupling of the quasiparticles to a magnetic flux.
The presence of the $`\mathrm{cos}(\sqrt{8\pi }\varphi _s)`$ term in the single chain Hamiltonian results in the quantization of $`N_s`$ in integer units. The presence of the $`\mathrm{cos}(\sqrt{2\pi }\varphi _s)\mathrm{cos}(\sqrt{2\pi }\theta _c)`$ term in $``$ results in the quantization condition that $`N_s+N_c`$ be an even integer! Physically, this means that excitations can have spin $`\mathrm{}`$ and charge 0 ($`N_s=2`$ and $`N_c=0`$), spin 0 and charge $`2`$ ($`N_s=0`$ and $`N_c=2`$), spin $`\mathrm{}/2`$ and charge $`1`$ ($`N_s=1`$ and $`N_c=1`$), etc., but that all the exotic quantum numbers of the soliton excitations of the isolated 1DEG are killed. Formally, the addition of the pair hopping term to the Hamiltonian of the 1DEG leads to a confinement phenomenon. Along the entire segment of chain between two spatially separated $`\pm \sqrt{\pi /2}`$ solitons, there is a change in sign of the pair hopping term. (See Eq. (40).) This leads to an energy which grows linearly with the separation $`x`$ between solitons, $`𝒥|x|`$, regardless of whether they are charge or spin solitons or antisolitons.
The importance of this observation becomes clear when we study the operators in whose correlation functions we are interested. Since
$$e^{i\sqrt{\pi /2}\theta _s(x)}\varphi _s(y)e^{i\sqrt{\pi /2}\theta _s(x)}=\varphi _s(y)\sqrt{\pi /2}\mathrm{\Theta }(yx),$$
(46)
and
$$e^{i\sqrt{\pi /2}\varphi _c(x)}\theta _c(y)e^{i\sqrt{\pi /2}\varphi _c(x)}=\theta _c(y)+\sqrt{\pi /2}\mathrm{\Theta }(xy),$$
(47)
it is clear that the fermion annihilation operator $`\mathrm{\Psi }_{1,}`$ creates a spin antisoliton and a charge antisoliton, while the $`2k_F`$ piece of the spin-raising operator, $`S_{2k_F}^+`$, creates a pair of spin solitons and a pair of charge antisolitons. Both these combinations decay into a set of free solitons in the absence of the interchain coupling, but in its presence, the former becomes a bound state, and the latter a resonant state. Thus, $`G^<`$ develops a coherent piece with a well defined dispersion relation as superconducting phase coherence between chains occurs. $`𝒮`$ develops a resonant peak at a temperature well below $`T_c`$.
### A Zero spin soliton sector
For the case in which the spin gap $`\mathrm{\Delta }_s`$ of the isolated chain is large compared to the interchain coupling, the fluctuations of the spin field are high energy (fast) compared to any charge fluctuations, and indeed only slightly affected by the onset of superconducting order. In this limit, the eigenstates can be treated in the adiabatic approximation.
In the ground state ($`N_s=0`$) sector, the spin field fluctuations are little affected by $`_J`$; all spin correlations can thus be computed as in the previous section. Moreover, because of the spin gap, so long as $`T\mathrm{\Delta }_s`$, the spin fields can be approximated by their ground state. For computing the charge part of the wave function, we can replace the operator $`\mathrm{cos}(\sqrt{2\pi }\varphi _s)`$ in $`_J`$ by its expectation value at zero temperature in the decoupled ground state,
$$\mathrm{cos}(\sqrt{2\pi }\varphi _s)\mathrm{cos}(\sqrt{2\pi }\varphi _s)_o𝒞_s(\mathrm{\Lambda }\xi _s)^{\frac{K_s}{2}},$$
(48)
where the subscript $`\mathrm{`}\mathrm{`}o\mathrm{"}`$ refers to the expectation value in the ensemble with $`J_{SC}`$ set equal to zero, (see also Ref. C). This leaves us with a sine Gordon equation for the charge degrees of freedom, with potential
$$𝒥𝒞_s\mathrm{cos}(\sqrt{2\pi }\theta _c).$$
(49)
Again, we solve this problem by refermionizing
$$\mathrm{\Psi }_{c,\lambda }^{}=𝒩_c\mathrm{exp}[i\sqrt{\pi /2}(\theta _c2\lambda \varphi _c)].$$
(50)
The refermionized form of the Hamiltonian is
$`_c=`$ $`i\stackrel{~}{v}_c[\mathrm{\Psi }_{c,1}^{}_x\mathrm{\Psi }_{c,1}\mathrm{\Psi }_{c,1}^{}_x\mathrm{\Psi }_{c,1}]`$ (53)
$`\stackrel{~}{\mathrm{\Delta }}_c[\mathrm{\Psi }_{c,1}^{}\mathrm{\Psi }_{c,1}^{}+\mathrm{H}.\mathrm{c}.]`$
$`+g_c\mathrm{\Psi }_{c,1}^{}\mathrm{\Psi }_{c,1}^{}\mathrm{\Psi }_{c,1}\mathrm{\Psi }_{c,1},`$
where
$`\stackrel{~}{v}_c`$ $`=`$ $`v_c\left({\displaystyle \frac{1}{4K_c}}+K_c\right),`$ (54)
$`\stackrel{~}{\mathrm{\Delta }}_c`$ $`=`$ $`{\displaystyle \frac{\pi 𝒥𝒞_s}{\mathrm{\Lambda }}},`$ (55)
$`g_c`$ $`=`$ $`2\pi v_c\left({\displaystyle \frac{1}{4K_c}}K_c\right).`$ (56)
Since $`N_s=0`$, the superselection rule implies $`N_c=2m`$, which upon refermionization is simply the condition:
$$\underset{\lambda }{}\lambda 𝑑x[\mathrm{\Psi }_{c,\lambda }^{}\mathrm{\Psi }_{c,\lambda }]=N_c/2=m.$$
(57)
It is also interesting to note that the superconducting pair creation operator can be expressed in an intuitively appealing form in terms of charge soliton creation operators
$$\widehat{\mathrm{\Delta }}^{}\mathrm{cos}(\sqrt{2\pi }\varphi _s)\mathrm{\Psi }_{c,1}^{}\mathrm{\Psi }_{c,1}^{}.$$
(58)
Recall that here the charge solitons are spinless fermions. This expression emphasizes the fact that spin gap formation, which is associated with the quenching of the fluctuations of the spin density phase, $`\varphi _s`$, can also be identified with the growth of the amplitude of the superconducting order parameter. While the charge solitons clearly also make a contribution to the amplitude of the order parameter, the phase of the order parameter comes entirely from the charge.
For $`K_c=1/2`$, just as for the Luther-Emery point for the spin fields, the refermionized Hamiltonian for the charged excitations is noninteracting and massive (gapped), and $`\mathrm{\Delta }_c=\stackrel{~}{\mathrm{\Delta }}_c`$. In computing the asymptotic form of correlations we will set $`K_c=1/2`$. We can now readily compute the expectation value of the pair hopping term so as to relate two physically important quantities: the excitation energy scale, $`\mathrm{\Delta }_c`$, and the interchain portion of the internal energy
$`\mathrm{\Delta }_c\mathrm{\Psi }_{c,1}^{}\mathrm{\Psi }_{c,1}^{}+\mathrm{H}.\mathrm{c}.`$ $`=𝒥\mathrm{cos}(\sqrt{2\pi }\varphi _s)\mathrm{cos}(\sqrt{2\pi }\theta _c)`$ (60)
$`=(\mathrm{\Delta }_c/\pi \xi _c)u_0(\mathrm{\Delta }_c,T),`$
where $`\xi _c=v_c/\mathrm{\Delta }_c`$ is the charge correlation length. Equation (60) has the form of a BCS gap equation with
$$u_0(\mathrm{\Delta }_c,T)=_0^{v_c\mathrm{\Lambda }}𝑑x\frac{1}{\sqrt{x^2+\mathrm{\Delta }_c^2}}\mathrm{tanh}\left(\frac{1}{2T}\sqrt{x^2+\mathrm{\Delta }_c^2}\right),$$
(61)
where the mean field relation for $`\mathrm{\Delta }_c(T)`$ is
$$u_0(\mathrm{\Delta }_c,T)=\frac{\pi v_c}{zJ_{SC}𝒞_s^2}.$$
(62)
Consequently we find the familiar BCS relations
$$T_c=0.57\mathrm{\Delta }_c(0),$$
(63)
$$\mathrm{\Delta }_c(0)=2v_c\mathrm{\Lambda }\mathrm{exp}[\pi v_c/zJ_{SC}𝒞_s^2],$$
(64)
$$\mathrm{\Delta }_c(T)1.74\mathrm{\Delta }_c(0)\sqrt{1T/T_c}\mathrm{for}TT_c.$$
(65)
In general, the actual form of $`\mathrm{\Delta }_c(0)`$ in terms of $`J_{SC}`$ and $`𝒞_s`$ is modified according to the microscopic value of $`K_c`$.
The transverse superconducting phase stiffness $`\kappa _{}`$ (proportional to the superfluid density) is
$$\kappa _{}=2\pi aH_J,$$
(66)
where $`d`$ is the spacing between chains and $`H_J`$ is given in Eq. (60). Thus at zero temperature $`\kappa _{}T_c^2/v_c`$. As is shown in Appendix B, for a system with equal areas of domains in which the stripes run along the $`x`$ and $`y`$ directions, the macroscopic phase stiffness is equal to the geometric mean of the superfluid density in the directions parallel and perpendicular to the chains, $`\overline{\kappa }=\sqrt{\kappa _{}\kappa _{}}`$. Since the phase stiffness along the chains is simply $`\kappa _{}=v_cK_c`$, it follows that $`\overline{\kappa }(T=0)`$ is (up to logarithmic corrections coming from $`u_0`$) simply proportional to $`T_c`$. This is a microscopic realization of a more general phenomenon which occurs in systems with low superfluid density; it is phase ordering, as opposed to pairing, which determines $`T_c`$. In a future publication, we will study the effects of quantum and thermal phase fluctuations on the evolution of the superfluid density of a quasi one-dimensional superconductor.
With little additional effort, we can study the pair field susceptibility $`\chi `$ at energies small compared to $`2\mathrm{\Delta }_s`$. In this low energy limit, as in Eq. (37), we can replace the spin operators in $`\chi `$ by their ground state expectation values.
The charge part of $`\stackrel{~}{\chi }`$ can be expressed in terms of the charge fermion fields
$$\stackrel{~}{\chi }_c(x,t)=\mathrm{\Psi }_{c,1}^{}(x,t)\mathrm{\Psi }_{c,1}^{}(x,t)\mathrm{\Psi }_{c,1}(0,0)\mathrm{\Psi }_{c,1}(0,0).$$
(67)
At the free charge fermion point ($`K_c=1/2`$) the corresponding spectral function is readily evaluated, for $`T=0`$ and $`\omega 2\mathrm{\Delta }_s`$, with the result
$`\chi (k,\omega )`$ $`=`$ $`\left({\displaystyle \frac{𝒞_su_0}{\xi _c}}\right)^2\delta (k)\delta (\omega )`$ (68)
$`+`$ $`{\displaystyle \frac{𝒞_s^2[\omega ^24E_c^2(k/2)+2\mathrm{\Delta }_c^2]}{4v_c^2|q_1E_c(q_2)q_2E_c(q_1)|}}\mathrm{\Theta }[\omega 2E_c(k/2)],`$ (69)
where $`E_c(k)`$ and $`q_{1,2}`$ are the analogs of Eqs. (26,27) with $`\mathrm{\Delta }_c`$ substituted for $`\mathrm{\Delta }_s`$ and $`v_c`$ for $`v_s`$.
Away from the Luther-Emery point, if $`g_c>0`$ ($`K_c<1/2`$), the two solitons repel, and hence the effect of $`g_c`$ can be ignored, but for $`g_c<0`$ ($`K_c>1/2`$), there is an attractive interaction between the two solitons and hence, this being after all a one dimensional problem, they form a bound state. This will slightly modify the expression for $`\chi `$.
### B The one hole sector
In the one soliton sector of the spin Hamiltonian, the adiabatic approximation requires reexamination. While for the most part, the spin modes are fast compared to the charge modes, the Goldstone mode (translation mode of the spin soliton) is slow compared to all other modes, and so must be treated in the inverse adiabatic approximation. Thus, we consider the charge Hamiltonian with a spin antisoliton at fixed position, $`R_s`$. The pair tunnelling term is then
$$𝒥𝒞_s\mathrm{sign}(xR_s)\mathrm{cos}(\sqrt{2\pi }\theta _c),$$
(71)
where we have used the fact that $`\xi _s=v_s/\mathrm{\Delta }_s`$ (which characterizes the width of the spin soliton) is small compared to the charge correlation length, $`\xi _c=v_c/\mathrm{\Delta }_c`$, to approximate the profile of the spin soliton by a step function. Upon refermionization, the charge Hamiltonian is still of the same form as Eq. (53) with the term proportional to $`\stackrel{~}{\mathrm{\Delta }}_c`$ replaced by
$$\stackrel{~}{\mathrm{\Delta }}_c\stackrel{~}{\mathrm{\Delta }}_c\mathrm{sign}(xR_s).$$
(72)
For $`K_c=1/2`$, upon the canonical transformation,
$$\psi _{c,1}=\mathrm{\Psi }_{c,1}^{},\psi _{c,1}=\mathrm{\Psi }_{c,1},$$
(73)
the charge soliton Hamiltonian is of the same form as the fermionic Hamiltonian of a commensurability 2 Peierls insulator, “polyacetylene”, in the presence of a topological soliton. As is well known, there is an index theorem that implies the existence of a zero energy bound state associated with the soliton, the famous “midgap state” or “zero mode”. All other fermionic states have energies greater than or equal to $`\mathrm{\Delta }_c`$. Importantly, since in this sector $`N_s=1`$, the superselection rule $`N_c=2m+1`$, requires that the fermion number is half integer!
$$\underset{\lambda }{}dx:[\psi _{c,\lambda }^{}\psi _{c,\lambda }]:=N_c/2=m+1/2.$$
(74)
This is essential, since with the midgap state occupied the fermion number is $`+1/2`$, while with it empty the fermion number is $`1/2`$. The midgap state is associated with the bound state of the spin and charge antisolitons.
To compute the charge contribution to the soliton creation energy we need to evaluate the difference between the ground state energies of the charge Hamiltonian in the presence and absence of a kink. We have done this by taking the limit of vanishing soliton width of a general expression of Takayama, Lin-Liu, and Maki, (and dividing by 2 for the spinless case). The resulting soliton creation energy is just $`\mathrm{\Delta }_c/2`$ ; in other words, the rest energy of the electron, i.e. the bound state of a spin soliton and a charge soliton, is
$$\mathrm{\Delta }_0=\mathrm{\Delta }_s+\mathrm{\Delta }_c/2\mathrm{\Delta }_s.$$
(75)
From this discussion, we can immediately conclude that for $`TT_c\mathrm{\Delta }_s`$, the one hole spectral function has a coherent piece and a multiparticle incoherent piece,
$$G^<(k,\omega )=Z(k)\delta [\omega (k)]+G^{(multi)}(k,\omega ),$$
(76)
where
$$(k)=\sqrt{v_s^2k^2+\mathrm{\Delta }_0^2}.$$
(77)
This follows from the fact that the bound state of a spin soliton and a charge soliton has the same quantum numbers as a hole. The multiparticle piece has a threshold slightly above the single hole threshold at $`\omega =(k)+2\mathrm{\Delta }_c`$.
The overlap factor, $`Z(k)`$, contains factors from both the spin and the charge parts of the wavefunction; so long as $`k\xi _s1`$, $`Z(k)=Z_c(k)Z_s(0)`$ where $`Z_s(0)`$ depends on the spin correlation length as in Eq. (32), and $`Z_c(k)`$ contains all remaining contributions. We can obtain a scaling form for $`Z_c`$ using the same method of analysis employed previously for $`Z_s`$. Specifically, at $`T=0`$ in the absence of interchain coupling, and for $`\omega 3\mathrm{\Delta }_s`$ and $`|k\xi _s|1`$, $`G^<`$ is given by the expression in Eq. (35). Since the opening of a charge gap does not affect the high energy physics, the dependence of $`G^<`$ on $`\mathrm{\Lambda }`$ is unaffected by the interchain coupling. Indeed, so long as $`\mathrm{\Delta }_c\mathrm{\Delta }_s`$, the dependence of $`G^<`$ on $`\mathrm{\Delta }_s`$ is likewise unchanged. Thus, by dimensional analysis, it follows that
$$Z(k)=Z_s(0)(\mathrm{\Lambda }\xi _c)^{\frac{1}{2}2\gamma _c}A_{\gamma _c}\stackrel{~}{f}(k\xi _c),$$
(78)
where $`\stackrel{~}{f}`$ is a scaling function and
$$A_{\gamma _c}=\frac{B(\gamma _c,\gamma _c+1/2)}{\mathrm{\Gamma }(\gamma _c)\mathrm{\Gamma }(\gamma _c+1/2)}.$$
(79)
Unfortunately, we do not have exact results from which to compute $`\stackrel{~}{f}(x)`$ explicitly, but there is no reason to expect it to have any very interesting behavior for small $`x`$.
At temperatures between $`T=0`$ and $`T=T_c`$, the same arguments lead to a simple approximate expression for the spectral function. Specifically, the principal temperature dependence comes from $`\mathrm{\Delta }_c`$ which is a decreasing function of $`T`$. At mean field level, the temperature dependence of $`\mathrm{\Delta }_c`$ can be computed from Eq. (61). In particular it vanishes at $`T_c`$ according to Eq. (65). Since fluctuation effects produce superconducting correlations between neighboring chains at temperatures above $`T_c`$, this simple mean field behavior will be somewhat rounded, but the qualitative point that $`\mathrm{\Delta }_c`$ becomes small at temperatures above $`T_c`$ is quite robust.
Consequently, the quasiparticle weight, $`Z`$, which is proportional to $`\mathrm{\Delta }_c^{2\gamma _c+\frac{1}{2}}`$, is a strongly decreasing function of $`T`$ which vanishes in the neighborhood of $`T_c`$. The quasiparticle gap, $`\mathrm{\Delta }_0`$, on the other hand, is only weakly temperature dependent, dropping from its maximum value $`\mathrm{\Delta }_0=\mathrm{\Delta }_s+\frac{1}{2}\mathrm{\Delta }_c(0)`$ at $`T=0`$ to $`\mathrm{\Delta }_0=\mathrm{\Delta }_s`$ in the neighborhood of $`T_c`$. Scattering off thermal excitations will, of course, induce a finite lifetime for the quasiparticle at finite temperatures.
Neither a charge soliton nor a spin soliton can hop from one chain to the next, but a hole can. The problem of the transverse dispersion of the coherent peak in the single hole spectral function is addressed in Appendix A. Not surprisingly, we find that the effective interchain hopping matrix element, $`t_{}`$, is replaced by an effective interchain hopping matrix element,
$$t_{}^{eff}=Z(k)t_{}.$$
(80)
Thus, the dispersion of the coherent peak transverse to the chain direction is an independent measure of the degree of interchain coherence.
### C The two spin soliton sector
To compute $`𝒮`$, we need to study states in the $`N_s=2`$ sector. Interestingly, (in contrast to the case of an ordered CDW) in a quasi one-dimensional superconductor, the $`2k_F`$ spin density wave operator also creates two charge antisolitons: $`N_c=2`$. Again, for the most part, the spin fluctuations are fast and high energy compared to the scale of the charge fluctuations, and can thus be treated in the adiabatic approximation - indeed, they are little affected by the presence of the interchain Josephson coupling. However, there are two low frequency modes associated with the soliton translational degrees of freedom, which must be treated in the antiadiabatic approximation. Consequently, we obtain an effective Schrödinger equation governing the center of mass motion of the two spin solitons:
$$H^{eff}2\mathrm{\Delta }_s\frac{1}{2M^{}}\underset{j=1}{\overset{2}{}}\frac{^2}{x_j^2}+V(x_1x_2),$$
(81)
where $`x_j`$ is the position of soliton $`j`$,
$$M^{}=\mathrm{\Delta }_s/v_s^2,$$
(82)
and $`V`$ is the adiabatic spin soliton potential, obtained by integrating out the (relatively fast) fluctuations of the charge degrees of freedom.
To compute $`V(R)`$, we again rely on the analogy between the refermionized version of the charge part of the Hamiltonian and solitons in polyacetylene. In this case, $`V(R)`$ is recognized as the difference in the ground state energy of a massive Dirac fermion in the presence and absence of a pair of zero width solitons separated by a distance $`R`$, i.e the Hamiltonian in Eq. (53) with
$$\stackrel{~}{\mathrm{\Delta }}_c\stackrel{~}{\mathrm{\Delta }}_c\mathrm{sign}(4x^2R^2).$$
(83)
Since $`2N_c=2`$, this energy difference is to be computed in the fermion number $`1`$ sector.
From the results in the previous section, it follows that
$$V(R)\mathrm{\Delta }_c\mathrm{as}R\mathrm{},$$
(84)
since in this limit, the two solitons are noninteracting, and reduce to the solution discussed in the previous section. Similarly, since for $`R=0`$, the energy approaches that of the uniform system with fermion number $`1`$,
$$V(R)\mathrm{\Delta }_c\mathrm{as}R0.$$
(85)
Moreover, from simple scaling, it is clear that
$$V(R)=\mathrm{\Delta }_c[1+v(R/\xi _c)],$$
(86)
where $`v(x)`$ is independent of the magnitude of $`\mathrm{\Delta }_c`$ and $`v(x)0`$ for $`x0`$ and $`x\mathrm{}`$. For intermediate $`R/\xi _c`$, we have been unable to obtain an analytic expression for $`v`$, although it is easily derived numerically, as described in Appendix C, with the result shown in Fig. 1. As can be seen, $`v(x)`$ rises from 0 to a gentle maximum at $`x0.3`$ where $`v(0.3)0.2`$, and then drops exponentially back to zero at large separation.
What this means is that there is no true bound state in the spin one excitation spectrum. The spin one excitations, even in the superconducting state, are always unstable to decay into a pair of far separated spin 1/2 quasiparticles. However, near the threshold energy, $`\omega =2\mathrm{\Delta }_s+\mathrm{\Delta }_c`$, there is a nearly bound (resonant) state with a lifetime which is exponentially long. Treating Eq. (81) in the WKB approximation, we see that the decay rate of the resonant state is
$$\mathrm{\Gamma }\mathrm{exp}\left[B(v_c/v_s)\sqrt{\mathrm{\Delta }_s/\mathrm{\Delta }_c}\right],$$
(87)
where $`B`$ is a constant of order 1.
Using the fact that $`𝒮\mathrm{\Lambda }^{\frac{1}{K_c}}`$ in the absence of a charge gap and utilizing the same scaling arguments applied previously to the coherent piece of $`G^<`$, it is easy to see that the weight associated with this resonant state is proportional to $`\mathrm{\Delta }_c^{\frac{1}{K_c}1}`$. However, because the barrier height is small compared to $`\mathrm{\Delta }_c`$, the thermal decay of the resonant bound state will become large, due to activation over the barrier, at a temperature well below $`T_c`$.
### D The “BCS-like case”: no preexisting spin gap
When there is no spin gap on the isolated chain and there are repulsive interactions in the charge sector, the interchain Josephson coupling is perturbatively irrelevant. Thus, the usual case for a quasi one-dimensional superconductor is the already analyzed case with a preexisting spin gap. However, it is worthwhile considering the case (with $`K_c<1`$) in which both the spin gap and the superconducting coherence are induced by a relevant interchain Josephson coupling. This case, even though quasi one-dimensional, is much more akin to the usual BCS limit, in that there is a single gap scale in the problem, and pairing (gap formation) and superconducting coherence occur at the same temperature and with roughly the same energy scale. This case has been analyzed extensively in the literature. It should be noted, however, that here, too, since the “normal” state is a non-Fermi liquid, the coherent piece of all spectral functions will be strongly temperature dependent below $`T_c`$, and vanish in the neighborhood of $`T_c`$.
## V Summary and Implications for Experiment
In this paper, we have obtained explicit and detailed results for the properties of the superconducting state of a quasi one-dimensional superconductor. We have studied this problem as a quantum critical phenomenon, in which the quantum critical point is reached in the 1d limit of no interchain coupling, and hence we have treated the interchain Josephson coupling as a small parameter. In particular, we expect (as discussed below) that the results are pertinent to underdoped and optimally doped high temperature superconductors, where self-organized stripe structures render the system locally quasi one-dimensional.
It is often argued that, even in fairly exotic circumstances, and even when the normal state is a non-Fermi liquid, the superconducting state itself is fairly conventional and BCS-like. We have shown that there are a number of ways in which this expectation is violated. In the first place, there are two “gap” scales, $`\mathrm{\Delta }_s\mathrm{\Delta }_c`$, whereas in a BCS superconductor there is one, $`\mathrm{\Delta }_0`$, and correspondingly two correlation lengths, $`\xi _s`$ and $`\xi _c`$, in place of the one, $`\xi _0`$, of BCS theory. However, both gaps are, in a very real sense, superconducting gaps: $`\mathrm{\Delta }_s`$ is associated with spin pairing (i.e. a nonzero value of $`\mathrm{\Psi }_{s,1}^{}\mathrm{\Psi }_{s,1}`$) and the existence of a local amplitude of the order parameter. $`\mathrm{\Delta }_c`$ is a measure of interchain phase coherence. In the case in which there is a preexisting spin gap on the isolated chain, $`\xi _s`$ remains finite at the quantum critical point, whereas $`\xi _c`$ diverges. The same holds true in the superconducting phase, a bit away from the quantum critical point, where $`\xi _c`$ diverges at $`T_c`$, while $`\xi _s`$ remains finite.
It is perhaps worth noting that many of these novel aspects of the superconducting state are considerably more general than the particular model we have solved. Indeed, recently, Lee derived similar results from the gauge theory formulation of a flux phase to superconductor transition. While this derivation presupposes rather different seeming microscopic physics, it does build in the doped insulator character of the superconducting state, which is the essential feature of the results. Likewise, many features we have discussed here bear a close resemblance to the dimensional crossover from a conjectured 2d non-Fermi liquid to a 3d superconductor envisaged in the context of the inter-layer tunnelling mechanism of high temperature superconductivity.
### A Summary of Results
For the benefit of the reader who skipped the technical exposition, we begin by summarizing our most important results. We consider here the case in which there is a preexisting spin gap, $`\mathrm{\Delta }_s/2T_c`$, on an isolated chain, and we focus on the effects of the interchain Josephson coupling between stripes at lower energies.
#### 1 Thermodynamic Effects
The effect of the interchain Josephson coupling is to produce an interchain coherence scale, $`\mathrm{\Delta }_c(T)`$. At mean field level, $`\mathrm{\Delta }_c(T)`$ vanishes for any $`T`$ above $`T_c`$, and while fluctuation effects will produce a small amount of rounding to this behavior, because of the large coherence lengths along the chain the degree of rounding will always be small in the quasi-1D limit. It is the coherence scale that determines $`T_c`$, in the sense that
$$T_c\mathrm{\Delta }_c(0)/2\mathrm{\Delta }_s/2.$$
(88)
($`\mathrm{\Delta }_c`$ is expressed in terms of the strength of the Josephson tunnelling matrix elements in Eq. (64).) The superfluid densities in the directions transverse and parallel to the chain direction are, respectively,
$$\kappa _{}=2au_0\mathrm{\Delta }_c^2/v_c,\kappa _{}=v_cK_c,$$
(89)
where $`u_0`$ is a constant (see Eq. (61)) which depends weakly on parameters, $`d`$ is the spacing between chains, $`v_c`$ is the charge velocity, and $`K_c`$ is the charge Luttinger parameter. In two dimensions, if, on average, there is a 4-fold rotationally invariant mixture of domains in which the chains run along the x and y directions, respectively, the macroscopic superfluid density is isotropic and given by
$$\overline{\kappa }(T)=\sqrt{\kappa _{}\kappa _{}}\mathrm{\Delta }_c(T).$$
(90)
#### 2 Single Hole Spectral Function
The common theme in the spectral functions is that all dependence on the interchain coupling (and hence all important temperature dependences in the neighborhood of $`T_c`$) are expressible in terms of the single coherence scale, $`\mathrm{\Delta }_c`$. Moreover, it is the spectral weight of the coherent features in the spectrum, rather than their energies, which are strongly temperature dependent! This is very different from the behavior of the spectral functions near $`T_c`$ in a three dimensional BCS superconductor.
Characteristic shapes of the single hole spectral function above and below $`T_c`$ are shown in Fig. 2. Above $`T_c`$, the single hole spectral function is a broad incoherent peak. Below $`T_c`$, there is a coherent delta-function piece and a multiparticle continuum at higher energy,
$$G^<(\stackrel{}{k},\omega )=Z(k_{})\delta [\omega (\stackrel{}{k})]+G^{multi},$$
(91)
where
$$(\stackrel{}{k})=\sqrt{v_s^2k_{}^2+\mathrm{\Delta }_0^2}+2t_{}Z(k_{})\mathrm{cos}(k_{}a)+\mathrm{}.$$
(92)
Here $`k_F+k_{}`$ and $`k_{}`$ are, respectively, the components of the crystal momentum parallel and perpendicular to the chain direction.
The energy gap for the coherent peak is
$$\mathrm{\Delta }_0(T)=\mathrm{\Delta }_s+\frac{1}{2}\mathrm{\Delta }_c(T),$$
(93)
and its spectral weight is given by
$$Z(k)[\mathrm{\Delta }_c(T)]^{2\gamma _c+\frac{1}{2}}.$$
(94)
Thus, $`Z(k)`$ (and with it the transverse bandwidth) is the most strongly temperature dependent feature of the spectral function.
The multiparticle incoherent piece $`G^{multi}`$ starts at a threshold energy $`(\stackrel{}{k})+2\mathrm{\Delta }_c(T)`$. This is the origin of the gap between the coherent peak and the incoherent shoulder in Fig. 2. Various forms of damping, including phase fluctuations transverse to the stripes, will broaden this structure, leading to a peak-dip-shoulder form of the spectral function. However, the distance from the coherent peak to the dip should be proportional to $`\mathrm{\Delta }_c(T)`$ and hence, at $`T=0`$, to $`T_c`$.
#### 3 The Spin Response Function
The spin response function is entirely a multiparticle continuum; even below $`T_c`$, we find that any spin-1 mode is unstable to decay into two spin 1/2 “quasiparticles”. However, at low temperature, we find that there is a spin-1 resonant state with an exponentially long lifetime near the threshold energy $`2\mathrm{\Delta }_s+\mathrm{\Delta }_c=2\mathrm{\Delta }_0`$, with momentum $`2k_F`$, where $`k_F`$ is the Fermi momentum on a stripe. Even here, because the barrier to decay is quantitatively small compared to $`T_c`$, we expect that no sharp resonant state will appear in the spectrum in the neighborhood of $`T_c`$. Rather, it will appear as the temperature falls below $`T0.2\mathrm{\Delta }_c(0)0.4T_c`$.
### B Implications of Two Scales
The existence of two scales in the superconducting state appears in different experiments in fairly obvious ways:
1) Since an electron has spin and charge, the gap measured in single particle spectroscopies, such as ARPES or tunnelling, is $`\mathrm{\Delta }_0=\mathrm{\Delta }_s+(1/2)\mathrm{\Delta }_c(T)`$. (See Eq. (75).) Manifestly, this gap scale decreases slightly with increasing temperature, but remains large, roughly $`\mathrm{\Delta }_s`$, above $`T_c`$. The gap scale $`\mathrm{\Delta }_s`$ is unrelated to $`T_c`$, and moreover $`\mathrm{\Delta }_s(T=0)2T_c`$, which physically is the statement that the onset of phase coherence, not pairing, is what determines $`T_c`$. Consequently, the zero-temperature superfluid density is a better predictor of $`T_c`$ than $`\mathrm{\Delta }_0(T=0)`$. (See Eq. (66) and subsequent discussion.) Similarly, pure spin probes, such as NMR or neutron scattering, see a gap which is approximately $`\mathrm{\Delta }_s`$ per spin 1/2. (See Eqs. (81).)
2) Experiments involving singlet pairs of electrons, such as Andreev tunnelling, could exhibit an energy scale $`\mathrm{\Delta }_c`$; a scale, moreover, which vanishes at (or near) $`T_c`$, and is related in magnitude to $`T_c`$ in a more or less familiar manner, $`\mathrm{\Delta }_c(0)/2T_c`$. More complicated spectroscopies, such as SIS tunnelling (e.g. tunnelling across a break junction) should reveal gap-like features with both energy scales, $`\mathrm{\Delta }_s`$ and $`\mathrm{\Delta }_c`$.
3) The existence of two correlation lengths implies that different measurements will find the order parameter magnitude depressed over distinct distances: If an impurity destroys the superconducting gap locally, the single-particle density of states, as determined, for instance, with a scanning tunnelling microscope, will basically recover over a length scale $`\xi _s`$ (although, subtle effects will persist out to a scale $`\xi _c`$). By contrast, the magnetic field strength near the core of a vortex, which otherwise would diverge logarithmically at short distances, is reduced inside a “core radius” due to the fact that the superfluid density is depressed, (i.e. there is a lower current density per unit phase gradient). Since this latter effect involves only charge motion, the vortex core radius is of order $`\xi _c`$. This “magnetic” core radius is measured, in principle, in $`\mu `$SR.
4) The superconducting state reflects the non-Fermi liquid character of the normal state in many ways, but it has a complex scalar order parameter as in a conventional (BCS) superconducting state. This means that we might expect well-defined elementary excitations with the quantum numbers of the electron quasiparticle, as indeed we have found. However, in a conventional superconductor, the quasiparticle energy is shifted by the opening of the gap, and the lifetimes of all elementary excitations (as observed e.g. in ultrasonic attenuation) are strongly temperature dependent below $`T_c`$. In the present case, it is the spectral weight associated with the elementary excitations which is strongly temperature dependent, not the lifetime or the energy. Moreover, even as $`T0`$, the quasiparticle weight remains small, in proportion to a positive power of the distance from the quantum critical point. See Eq. (94).
### C Two Scales in the High Temperature Superconductors
It has been noted, in the so-called “Yamada plot,” that $`T_c`$ in underdoped high temperature superconductors is proportional to the observed incommensurability in the low energy spin structure factor. The magnetic incommensurability, $`\pi /d`$, is inversely related to the mean separation between charge “stripes,” $`d`$. Thus, the Yamada plot implies that $`T_c`$ is inversely proportional to the mean spacing between stripes; As the stripes become more separated, and the electronic structure becomes more one dimensional, $`T_c0`$. This observation strongly supports the idea that the anomalous electronic properties of these materials reflect the properties of nearby phases of the 1DEG. Indeed, many of the spectral features listed above have been observed, with various levels of confidence, in experiments on the high temperature superconductors:
1) The best single particle spectra (ARPES and tunnelling) exist for Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8</sub> because it cleaves easily. For underdoped and optimally doped materials, the single particle gap as measured by tunnelling and ARPES is found to be large: $`\mathrm{\Delta }_035`$meV, in the “flat-band” region near the $`\overline{M}`$, or $`(\pi ,0)`$ and $`(0,\pi )`$, points of the Brillouin zone. ($`\mathrm{\Delta }_0/2T_c`$ lies in the range 2 to 6. The $`\overline{M}`$ region is where the maximum of a d-wave superconducting gap is expected.) The gap persists in some form or other to temperatures well above $`T_c`$. Moreover, $`\mathrm{\Delta }_0`$ increases with underdoping while $`T_c`$ decreases. This gap is quite clearly a superconducting gap in that it has (at low $`T`$) the characteristic d-wave form expected of the superconducting gap, and it evolves smoothly with overdoping into a gap of only slightly smaller magnitude which opens, in a more conventional manner, in the neighborhood of $`T_c`$. (We focus on the gap near the $`\overline{M}`$ point, especially, because there are both theoretical and experimental reasons to think that the “flat-band” region is associated with states in “stripes” or “fluctuating stripes.”) Moreover, as discovered first by Uemura and coworkers, $`T_c`$ is roughly proportional to the zero temperature superfluid density for underdoped materials, consistent with the notion that it is phase ordering, not pairing, which determines $`T_c`$.
2) Deutscher has argued that the gap scale determined by low temperature Andreev tunnelling spectroscopy is considerably smaller than that determined from single particle tunnelling measurements in underdoped materials, while the two gap scales approach each other in overdoped materials. This issue is well worth revisiting in more detail. The single particle gap scale is strongly apparent in SIS tunnelling spectra - we do not know of any convincing analysis which reveals the smaller charge gap scale in such experiments.
3) The vortex core radius has been measured with both scanning tunnelling microscopy (STM) and $`\mu `$SR. The $`\mu `$SR study measures the magnetic field distribution in the material, and infers the core radius from the high-field cutoff of the distribution. For large applied fields ($`B6T`$), both methods are in rough agreement that the core radius is about 15Å. However, the core radius deduced from the $`\mu `$SR measurements is strongly field dependent, so that at low fields ($`B0.5T`$) it yields a core radius around 120Å. By contrast, preliminary evidence from STM experiments suggests that the core radius measured by that method is not strongly field dependent, so that, in low fields, the results of the two methods differ by almost an order of magnitude. However, there appear to be differences in the STM results of different groups. Certainly, the core radius inferred from STM studies of the gap suppression in the vicinity of an impurity at zero magnetic field are suggestive of a rather short coherence length. While the experimental results are, by no means, definitive, we would tentatively like to explain the discrepancy between the STM and $`\mu `$SR results at low field as evidence of the existence of two coherence lengths in the superconducting state.
4) It has been realized for a long time that there are no sharp quasiparticle features in the ARPES spectrum near the superconducting gap maximum (near the $`\overline{M}`$ point of the Brillouin zone) in the normal state, and it has been argued that they disappear due to a lifetime catastrophe which occurs as the temperature is raised above $`T_c`$. Recent high-resolution ARPES measurements in optimally doped Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8</sub> have revealed a new picture of the emergence of these peaks. Within experimental resolution, neither the energy nor the width of the peak changes as the temperature is raised from well below $`T_c`$ to slightly above $`T_c`$; rather, it is the intensity of the peak that is strongly temperature dependent in the neighborhood of $`T_c`$. The intensity vanishes slightly above $`T_c`$, without any apparent change in the shape of the peak itself. Indeed, the sharp temperature dependence of this intensity in the neighborhood of $`T_c`$ is consistent with its being proportional to a fractional power of the (local) superfluid density. (See Eq. (94).) Additional evidence for this comes from an old observation of Harris et al. that, as a function of underdoping, the weight in the peak at low temperatures decreases with decreasing superfluid density. Moreover, Shen and Balatsky have argued that a small dispersion of the ARPES peak in the direction perpendicular to the putative stripe direction scales more or less with $`T_c`$, consistent with our Eq. (80). The distance between the coherent peak and the dip feature in ARPES curves near the $`\overline{M}`$ point decreases with underdoping, consistent with the zero temperature distance being proportional to $`T_c`$.
A similar temperature evolution has been observed for the so called “resonant peak” in neutron scattering in YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> and Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8</sub> (although no such feature has been seen in La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub>). We would like to identify this phenomenon, as well, with a dimensional crossover of the sort discussed here. However, the spin resonance we have found in the present model is clearly not directly related to the observed resonant peak. In particular, all features we have found are peaked at a momentum $`2k_F`$, while the resonant peak is centered on the antiferromagnetic wave-vector $`(\pi ,\pi )`$. Moreover, the peak we have found disappears through a lifetime catastrophe well below $`T_c`$, while the resonant peak is sharp immediately below $`T_c`$. Clearly, at the very least, to have a theory with anything more than a very rough caricature of this observed magnetic behavior, we need to expand the considered model to include the effects of the antiferromagnetic “strips” between the “stripes.”
### D Further Implications for High Temperature Superconductors
Finally, we end with a few additional observations concerning insights into the behavior of the high temperature superconductors that can be obtained from the analysis of this paper:
In the superconducting state of a structurally quasi-1D superconductor, we have found there are two emergent length scales, $`\xi _s`$ and $`\xi _c`$. If the quasi-1D electronic structure is self organized, as it is in the high temperature superconductors, there are potentially two additional emergent length scales: the mean spacing between stripes, $`d`$, and the persistence length of the stripes, $`\xi _{stripe}`$. $`d`$ can be determined directly from the charge incommensurability (or indirectly from the spin) structure factor. $`\xi _{stripe}`$ is much harder to determine experimentally, although it is bounded below by the correlation length of the magnetic order. So long as $`\xi _{stripe}`$ is the longest length scale in the problem, i.e. so long as $`\xi _{stripe}\xi _c`$, it is possible to assume, as we have here, that the superconducting properties of the system are quasi one-dimensional. Where this inequality is violated, the correct theory of the superconducting state needs to be significantly modified. As was pointed out previously, so long as the weaker condition, $`\xi _{stripe}\xi _s`$, is satisfied, it is possible to have a one dimensional theory of spin gap formation. At present experiments are unclear about the range of doping and materials for which either of these inequalities is satisfied, which is the most important source of uncertainty in the application of these ideas to the high temperature superconductors. Certainly, with sufficient overdoping, the stripes loose their integrity and the application of these ideas becomes suspect.
To get a feeling for magnitudes, we can make rough quantitative estimates of the remaining length scales from well-established experimental data in the high temperature superconductors, although numbers vary from material to material, and as a function of doping concentration, $`x`$. The spin velocity in the undoped antiferromagnet is around $`v_s0.8`$eV-Å, and the superconducting gap is $`\mathrm{\Delta }_035`$meV, so $`\xi _s`$ 20Å. The charge coherence length $`\xi _c=\xi _s(v_c/v_s)(\mathrm{\Delta }_s/\mathrm{\Delta }_c)`$, so if we estimate $`v_c/v_st/J23`$ and $`\mathrm{\Delta }_c2T_c16`$meV, we find that characteristically $`\xi _c100150`$Å. (This is in good agreement with the $`\mu `$SR measurement of the vortex core radius cited above. ) The spacing between stripes is in the range of four or more lattice constants, $`d16`$Å.
A crossover magnetic field, which can be identified as a mean field $`B_{c2}`$, can be estimated as the field at which there is one vortex per coherence length $`\xi _c`$ between each pair of neighboring stripes; this leads to an estimate
$$B_{c2}\varphi _0/\xi _cd,$$
(95)
where $`\varphi _0=hc/2e`$ is the superconducting flux quantum. While $`B_{c2}`$ estimated in this fashion is quite large ($`\varphi _0/\xi _cd=80T`$ for $`d=16`$Åand $`\xi _c=100`$Å) it is small compared to the characteristic magnetic field,
$$B_s2\varphi _0/\xi _sw,$$
(96)
at which orbital effects lead to the destruction of the spin gap. Here, $`w`$ is the “width” of a stripe - i.e. the width of the 1d region involved in spin gap formation. In the “spin gap proximity effect” mechanism proposed previously this would imply that $`w`$ is one to two times the crystalline lattice constant. The extremely large value of $`B_s`$ rationalizes the lack of any observable reduction of the spin gap temperature in the recent NMR experiments of Gorny et al. up to fields as high as 12T in YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub>. (See, also, the discussion in Ref. C.)
One important difference between a stripe phase and the array of chains studied here, is that in the stripe phase there are additional electronic degrees of freedom which live in the antiferromagnetic strips between the stripes. The two-component nature of the electronic structure of doped antiferromagnets, is characteristic of the micro-phase separation physics that gives rise to this state. Of course, the antiferromagnetic strips are themselves quasi one-dimensional magnets, so that any magnetic ordering must be viewed, in similar spirit to that considered here, as resulting from a dimensional crossover. Indeed, it is certainly the spins in the insulating strips that make the dominant contribution to the “resonant peak” observed in neutron scattering. A detailed theory of this peak is beyond the scope of the present model, but is embodied in the spin gap proximity effect.
However, we have found a new neutron scattering resonance for a quasi one-dimensional superconductor. While the dimensional crossover causes no bound state in the spin-1 excitations, we find a resonant state of two spin-$`\frac{1}{2}`$ quasiparticles appearing below $`T0.4T_c`$. The mode appears at an energy $`2\mathrm{\Delta }_s+\mathrm{\Delta }_c=2\mathrm{\Delta }_0`$, or twice the single particle gap as measured by ARPES or tunneling, and at momentum $`2k_F`$, where $`k_F`$ defines the Fermi surface associated with a stripe. Since this is a four soliton resonance, it may be qualitatively sensitive to deviations from the limit $`\mathrm{\Delta }_c<<\mathrm{\Delta }_s`$, so that the resonance is likely to be most well defined in the underdoped region where $`T_c\mathrm{\Delta }_0`$.
Finally, we remark that the ARPES spectrum along the symmetry direction from (0,0) to $`(\pi ,\pi )`$, i.e. along the ray which is expected to pass through the node of a d-wave gap function, is very different in character from that in the $`\overline{M}`$ that we have discussed. In clean samples of optimally doped Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8</sub>, there is a peak in the spectral function both above and below $`T_c`$, and the peak reaches the Fermi surface at a well defined “nodal point”, $`\stackrel{}{k}_n=(0.44\pi ,0.44\pi )`$. This peak does not exhibit the characteristics of a quasiparticle peak, in that its width is always larger than its energy; indeed, it seems to exhibit quantum critical behavior reminiscent of a Luttinger Liquid. Moreover, there is no qualitative change in the temperature evolution of this peak as the temperature is lowered from two or three times $`T_c`$ down to temperatures as low as at least 1/2 $`T_c`$; the character of the nodal excitation seems to be remarkably insensitive to the onset of superconductivity. By contrast, in optimally doped La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub>there is apparently no observable peak along the nodal direction, and indeed little or no spectral weight within about 0.5 eV of the Fermi energy. Indeed, recent neutron scattering studies of the low energy magnetic scattering in the neighborhood of $`2\stackrel{}{k}_n`$ have revealed the existence of a clean spin gap at low temperatures, which is apparently inconsistent with the existence of any gapless nodal quasiparticle excitations.
It is clear that whatever spectral response is observed near $`\stackrel{}{k}_n`$ is not associated with the vertical and horizontal stripes studied here, because a stripe wave vector does not span the “Fermi surface” along this direction. It could be associated with diagonal stripes, which have been observed recently in various insulating materials, in which case the observed quantum critical behavior might truly be that of a Luttinger liquid. An alternative picture is backflow associated with holes that have not condensed into vertical or horizontal stripes. Both explanations are conceivable as there are strong reasons to expect that the orienting potential, which locks the stripes along a particular (vertical or horizontal) crystallographic direction to be stronger in La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> than in Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8</sub>. However, other sources of quantum critical behavior are certainly possible. We will defer further discussion of these classes of excitations to a future study.
###### Acknowledgements.
We would like to acknowledge useful discussions with E. Fradkin, L. Pryadko, J. Tranquada, Z. X. Shen, P. B. Wiegmann, A. Tsvelik, O. Zachar, and C. Nayak. This work was supported in part by NSF grants number DMR98-08685 (EWC and SAK), DMR-9814289 (SAK) and DOE grant number DE-AC02-98CH10886 (VJE). D. O. has been supported by the Rothschild Fellowship. SAK gratefully acknowledges the hospitality of the Physics and Applied Physics Departments at Stanford University where much of this work was carried out.
## A The effect of interchain single particle hopping
Until now, we have ignored single particle hopping between chains. This is because, especially in the presence of a spin gap, it is irrelevant in the renormalization group sense. However, in the superconducting state, we expect the quasiparticles to be able to propagate coherently between chains. Because these terms are irrelevant, their effects on the spectrum can be computed in ordinary degenerate perturbation theory. It is easy to see that to first order in the interchain hopping, the quasiparticle energy is
$$(\stackrel{}{k})=\sqrt{v_s^2k_{}^2+\mathrm{\Delta }_0^2}+Z(k_{})ϵ^{()}(\stackrel{}{k}_{})+𝒪\left(ϵ^{()}\right)^2,$$
(A1)
where $`ϵ^{()}(\stackrel{}{k}_{})=2t_{}\mathrm{cos}(k_{}a)`$ is the interchain contribution to the quasiparticle dispersion, and $`k_{}+k_F`$ and $`\stackrel{}{k}_{}`$ are, respectively, the components of the crystal momentum parallel and perpendicular to the chain direction.
This is highly reminiscent of the spectrum we would have obtained were we to compute the spectrum of a quasi one-dimensional superconductor using BCS mean field theory
$`_\stackrel{}{k}^{(BCS)}=`$ $`\sqrt{[v_Fk+ϵ^{()}(\stackrel{}{k}_{})]^2+\mathrm{\Delta }^2}`$ (A2)
$`=`$ $`_k^{(BCS)}+[v_sk/_k^{(BCS)}]ϵ^{()}(\stackrel{}{k}_{})+𝒪\left(ϵ^{()}\right)^2,`$ (A3)
with the differences that the Fermi velocity is replaced by the (slower) spin velocity, the superconducting gap is the sum of the (single-chain) spin gap and the (interchain) charge gap, and the interchain bandwidth is reduced by the quasiparticle weight factor $`Z`$.
## B Macroscopic Superfluid Density
In this appendix, we compute the macroscopic phase stiffness (superfluid density) tensor $`K_{ab}[\kappa ]`$ in two dimensions $`(a=x,y)`$ given a microscopic distribution of the (in general anisotropic) local phase stiffness tensor, $`\kappa _{ab}(\stackrel{}{r})`$. We include the derivation here for pedagogical purposes, although the results exist elsewhere in the literature .
$`\kappa `$ determines the relation between the local current density, $`\stackrel{}{j}(\stackrel{}{r})`$ and the gradient of the phase according to
$$j_a(\stackrel{}{r})=\kappa _{ab}_b\theta (\stackrel{}{r}).$$
(B1)
From the equation of continuity, it follows that $`\stackrel{}{}\stackrel{}{j}=0`$, so we can express $`\stackrel{}{j}`$ in terms of a potential, $`j_a(\stackrel{}{r})=ϵ_{ab}_b\varphi (\stackrel{}{r})`$, so that
$$ϵ_{ab}_b\varphi (\stackrel{}{r})=\kappa _{ab}_b\theta (\stackrel{}{r}).$$
(B2)
To compute $`K_{xx}`$ in a rectangular geometry, this equation is to be solved subject to the boundary conditions that $`\theta =0`$ for $`x=0`$ and $`\theta =\mathrm{\Delta }\theta `$ for $`x=L_x`$ (independent of $`y`$) and (from the condition that no current can flow out of the sample in the $`y`$ direction) $`\varphi =0`$ for $`y=0`$ and $`\varphi =\mathrm{\Delta }\mathrm{\Phi }`$ for $`y=L_y`$. For a given distribution of $`\kappa `$, we solve this equation for given $`\mathrm{\Delta }\theta `$ to determine $`\mathrm{\Delta }\varphi `$, from which we determine $`K`$ according to
$$K_{xx}[\kappa ]=\mathrm{\Delta }\varphi /\mathrm{\Delta }\theta .$$
(B3)
The key observation is the same potential and phase that satisfy Eq. (B3), also satisfy the dual equation
$$ϵ_{ab}_b\theta (\stackrel{}{r})=\kappa _{ab}^D_b\varphi (\stackrel{}{r}),$$
(B4)
where
$$\kappa _{ab}^D(\stackrel{}{r})ϵ_{ac}\kappa _{cd}^1(\stackrel{}{r})ϵ_{db}.$$
(B5)
Therefore
$$K_{xx}[\kappa ]K_{yy}[\kappa ^D]=1.$$
(B6)
We can apply this general result to the problem of interest here. Consider the case of a square geometry in which, because of some assumed domain structure, the system is macroscopically isotropic ($`\overline{\kappa }K_{xx}=K_{yy}`$) despite the existence of microscopic anisotropy in each ”stripe” domain. It follows that
$$\overline{\kappa }(\kappa _{},\kappa _{})\overline{\kappa }(1/\kappa _{},1/\kappa _{})=1.$$
(B7)
It follows that $`\overline{\kappa }(\kappa _{},\kappa _{})=\sqrt{\kappa _{}\kappa _{}}`$. Other solutions to Eq. (B7) exist, but are not homogeneous functions.
## C The Effective Potential, $`v(x)`$
To compute the effective potential which appears in Eq. (81), we consider the discrete version of the refermionized Hamiltonian,
$$H=\underset{n}{}[t_0+(1)^n\mathrm{\Delta }(n)/2][c_n^{}c_{n+1}+\mathrm{H}.\mathrm{C}.],$$
(C1)
where $`v_c=2t_0`$ and
$$\mathrm{\Delta }(n)=\mathrm{\Delta }_c\mathrm{sign}(R^24n^2),$$
(C2)
corresponding to a pair of solitons separated by a distance $`R`$. (We have set the lattice constant equal to 1.)
We compute the ground state energy on a system of $`2N`$ sites by computing the single particle eigenvalues, and then summing over the lowest lying $`N1`$ of them to get the total energy as a function of $`R`$. This is precisely the program carried out previously to study various solitonic states in the SSH model of polyacetylene. With open boundary conditions, the Hamiltonian matrix is tridiagonal, and so particularly simple to study numerically for large system sizes. We have carried out this program numerically for system sizes up to $`2N=3000`$, and for $`\mathrm{\Delta }_c=`$ 0.2, 0.1, 0.05, and 0.02; the continuum limit is obtained when $`\mathrm{\Delta }_c0`$. Even for the smallest values of $`\mathrm{\Delta }_c`$, we find no significant finite size effects at these large system sizes. The results for $`v(x)`$ computed in this way are summarized in Fig. 1. The fact that the asymptotic value of $`v`$ is always slightly negative is a reflection of the fact that in the limit of small soliton width (which is equal to 1, the lattice constant, in the present calculation), the soliton creation energy is a very strongly varying function of the width, as found previously by Takayama, Lin-Liu, and Maki, and only approaches its true asymptotic limit $`\mathrm{\Delta }_c/2`$, when $`1/\xi _c`$ is extremely small. |
warning/0001/cond-mat0001215.html | ar5iv | text | # Amorphization of Vortex Matter and Reentrant Peak Effect in YBa2Cu3O7-δ
## I Introduction
The advent of high T<sub>c</sub> superconductors \[HTSC\] provided a new impetus to studies on phase transformations in the flux line lattices (FLL) of type-II superconductors in general and they remain a subject of considerable current interest. The competition amongst the elasticity of vortex medium, the pinning of vortex lines and the fluctuation effects of the thermal energy is expected to yield an ordered vortex solid (presumably a dislocation free Bragg glass phase ), a plastically deformed vortex solid ( presumably a vortex glass phase with proliferation of dislocations) and a pinned as well as an unpinned liquid phases in different parts of the field-temperature (H,T) phase space . In the HTSC cuprate systems, thermal fluctuations are large and in the high temperature part of the (H,T) space (i.e., for $`t=\frac{T}{T_c(0)}`$ $``$ 0.9, where T<sub>c</sub>(0) is the superconducting transition temperature in zero field), a first order melting transition occurs between an ordered solid and a nearly pinning free vortex liquid state . This transition in the most investigated YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> (YBCO) system is located at the upper edge of the peak effect (PE) phenomenon in the critical current density J<sub>c</sub> for fields of the order of a few tesla. At lower temperatures (i.e., $`t<0.9`$) and higher fields, where thermal fluctuations are weaker, a ‘fishtail’ or ‘second magnetization peak’ anomaly occurs presumably across a $`\mathrm{𝑡𝑟𝑎𝑛𝑠𝑓𝑜𝑟𝑚𝑎𝑡𝑖𝑜𝑛}`$ between a weakly pinned solid and a stronger pinned solid .
The detailed characteristics of the evolution of the anomalous variation in J<sub>c</sub> across the (H,T) region of the second magnetization peak and that of the conventional peak effect are far from being understood in the context of a variety of HTSC systems . In contrast, in the conventional low T<sub>c</sub> superconductors (LTSC) with considerably smaller (but often not negligible) thermal fluctuations, a PE near the normal state boundary (H<sub>c2</sub> line) is a common occurrence over the entire field-temperature regime . Various attempts to gain an understanding of the anomalous variations in J<sub>c</sub> attribute this phenomenon to a softening of the (weakly pinned) ordered lattice and a consequent transformation into a more strongly pinned amorphous solid or a pinned liquid phase . Disentangling the effects of the various possible sources of pinning and that of the thermal fluctuations has however remained difficult . Despite this, there is a widespread acceptance ( supported by evidence from microscopic $`\mu `$SR studies as well ) that the anomalous (sharp) variations in J<sub>c</sub> signal a change in the state (i.e., in terms of spatial and temporal correlations) of the vortex matter. In the collective pinning framework due to Larkin and Ovchinnikov , J<sub>c</sub> relates inversely to the correlation volume V<sub>c</sub> of the Larkin domain as,
$$J_cH=\sqrt{\frac{n_p<f_p^2>}{V_c}},$$
(1)
where n<sub>p</sub> is the density of pinning centers and f<sub>p</sub> is the elementary pinning interaction. Furthermore, in the context of LTSC systems, it is now known that upon increasing the effective pinning either externally by enhancing the quenched random inhomogeneities in the atomic lattice or by changing the magnetic field in a given sample, the PE often develops internal structure. This has led to suggestions of a richer and more complex transformations occurring in a stepwise manner between an ordered ( quasi-lattice) phase and a fully amorphous one. Even more interestingly, the low field part of the phase diagrams in the weakly pinned samples of an archetypal LTSC system 2H-NbSe<sub>2</sub> show a variety of novel features, including a reentrance into a disordered phase which is qualitatively similar to a reentrant glass / reentrant pinned liquid, postulated theoretically .
In this paper, we report on the search and the identification of the anomalous variations in J<sub>c</sub> via the ac magnetization measurements in a twinned crystal (T<sub>c</sub>(0) $``$ 93.3 K) of the YBCO system in the rarely reported low field-high temperature part of the vortex phase diagram. In the YBCO system, the presence of twin boundaries is considered to facilitate the occurrence (and detection) of the PE phenomenon. The twinned crystal piece (dimensions $``$ 0.5$`\times `$0.5$`\times `$0.04 mm<sup>3</sup> and mass $``$ 700 $`\mu `$g) chosen for the measurements being reported here is obtained by flux growth technique, along with detwinned YBCO crystals on which the first order melting transition was observed via differential calorimetry . In this twinned crystal, we are able to detect the PE phenomenon over a wide a field range (20 Oe to 10 kOe), the lower limit being much lower than any previous study. The evolution in the characteristic of the PE is similar to that in the weakly pinned samples of a variety of LTSC systems , but some significant differences exist as well.
## II EXPERIMENTAL DETAILS
The ac magnetization measurements were performed using a well shielded home built ac susceptometer . The dc field (co-axial with ac field) was kept parallel to the c-axis of the YBCO crystals. Both twinned and detwinned samples were investigated but PE was detected only in the former. Data were taken in both the field-cooled and the zero field-cooled (ZFC) modes. Most of the data shown are in the former case, where the PE is more conspicuous.
## III RESULTS AND ANALYSIS
### A Location of the Peak Effect and elucidation of its characteristic features
Figures 1 and 2 present a collation of the temperature dependent in-phase ac susceptibility data recorded with h<sub>ac</sub> of 0.5 Oe (r.m.s.) and at the frequency 211 Hz in the field ranges where anomalous variations in J<sub>c</sub> could be identified. In zero field, the normalized 4$`\pi \chi `$ has a value -1 (c.g.s. units) in the superconducting state and it crashes sharply towards the zero value (the transition width (10-90$`\%`$), $`\mathrm{\Delta }`$T<sub>c</sub>(0)$``$ 0.8 K ) in a featureless manner (see inset of Fig.1(a)). We now focus on the $`\chi (T)`$ responses shown in the Figs.1(a) through 1(f) as H increases from 20 Oe to 250 Oe. In H=20 Oe (where FLL a<sub>0</sub> $``$ 1.1$`\times `$10<sup>4</sup>Å), $`\chi (T)`$ is still featureless, though the superconducting transition is broader. Figures 1(b) and 1(c) show that as H increases from 40 Oe ( a<sub>0</sub> $``$ 8$`\times `$10<sup>3</sup>Å) to 170 Oe (a<sub>0</sub> $``$ 4$`\times `$10<sup>3</sup>Å), a characteristic feature reminiscent of the PE develops across the temperature region marked by a pair of arrows at T<sub>pl</sub> and T<sub>p</sub> in the respective panels. The shape of the $`\chi (T)`$ curve at H=250 Oe (a<sub>0</sub> $``$ 3.3$`\times `$10<sup>3</sup>Å) in the inset of Fig.1(f) identifies the PE phenomenon.
Within the Bean’s critical state model prescription of the magnetization of irreversible superconductors , the shielding response $`\chi (T)`$ can be approximated as :
$$\chi (T)1+\alpha \frac{h_{ac}}{J_c};forh_{ac}<<H^{},$$
(2)
$$\chi (T)\beta \frac{J_c}{h_{ac}};forh_{ac}>H^{},$$
(3)
where $`\alpha `$ and $`\beta `$ are geometry and size dependent factors, H is the (threshold) parametric field at which the magnetic field penetrates the center of the sample. Equations 2 and 3 imply that the temperature variation in $`\chi `$ is governed by the temperature variation of J<sub>c</sub> in a given H. Thus, the observed minimum in $`\chi (T)`$ in the inset of Fig.1(f) marks a peak in J<sub>c</sub>, i.e., the PE, which is characterized by two temperatures, the onset T<sub>pl</sub> and the peak at T<sub>p</sub>. The vanishing of J<sub>c</sub> above T<sub>p</sub>, i.e., the depinning of the vortex state is characterized by the sharp increase in $`\chi `$ to its normal state value. A depinning temperature T<sub>dp</sub> (notionally the mid point of collapse in J<sub>c</sub>) is obtained from the peak in d$`\chi (T)`$/dT, examples of which are shown in the insets of Fig.1(b) and Fig.2(d).
The other noteworthy features contained in Fig.1 and Fig. 2 could be summarized as follows.
(1) The PE is sharpest around an intermediate field of $``$1000 Oe, it becomes a broad anomaly for both higher and lower fields, and cannot be detected above 10 kOe and below 40 Oe.
(2) Characteristic structure (e.g., two well resolved minima instead of a single composite minimum in $`\chi `$), very similar to what has been seen in LTSC systems, appears in the intermediate field regime as evident in Figs.1(f), 2(a) and 2(b). Such a structure has been reported in an other study , but in higher fields than is shown here.
(3) Depinning phenomenon across T<sub>dp</sub> is a sharp anomaly. A simple fit of the form: $`4\pi \chi =a+T/\mathrm{\Delta }T_{dp},`$ yields a characteristic width $`\mathrm{\Delta }`$T<sub>dp</sub>, whose field dependence is shown in the inset of Fig.2(b). Interestingly, $`\mathrm{\Delta }`$T<sub>dp</sub>(H) is smallest for H$``$1000 Oe and increases for both increasing and decreasing fields, as shown in the inset of Fig.2(b). The $`\mathrm{\Delta }`$T<sub>dp</sub>(H) at 1 kOe is smaller than its value even at zero field (see the encircled data point), the latter is very close to the $`\mathrm{\Delta }`$T<sub>c</sub>(0), which is shown in the inset of Fig.1(a).
### B Observation of the Differential Paramagnetic Effect
Figures 3(a) and 3(b) show the $`\chi (T)`$ values measured with a higher h<sub>ac</sub> of 2 Oe (r.m.s.) at a frequency of 211Hz for the vortex states obtained in 320 Oe and 20 Oe, respectively. The PE in the main panel of Fig.3(a) is pronounced . A significant new feature to note, however, is that $`\chi `$ makes a sharp transition from diamagnetic values to paramagnetic values across the depinning temperature T<sub>dp</sub>. The inset in Fig.3(a) reveals a small paramagnetic peak above the temperature T<sub>irr</sub> (only slightly greater than the so designated T<sub>dp</sub>), the ubiquitous irreversibility temperature. The paramagnetic $`\chi `$ response returns to the background level on crossing over to the normal state at T<sub>c</sub>(H). An identical paramagnetic peak located at the edge of the depinning transition has earlier been reported in LTSC systems, such as, CeRu<sub>2</sub> and 2H-NbSe<sub>2</sub> and identified with the differential paramagnetic effect (DPE) . Between T<sub>irr</sub>(H) and T<sub>c</sub>(H), the vortex state is reversible, for which dM/dH is positive, i.e., diamagnetic dc magnetization decreases as H increases, implying, a DPE. The paramagnetic peaks in the $`\chi (T)`$ data in single crystals of YBCO and BSCCO had been reported in other studies as well . Morozov et al had reported a sharp paramagnetic peak in a single crystal of Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8</sub> (BSCCO) via local micro-Hall sensors at temperatures where a step change in equilibrium magnetization (M<sub>eq</sub>) due to FLL melting is expected. However, no peak effect has been reported in these studies. The question, as to whether the sharp onset of the DPE peak marks the step change in M<sub>eq</sub> at the FLL melting transition and/or it reflects the depinning transition located at the upper edge of the PE phenomenon, requires more detailed investigations. We note, however, that Ishida et al and Ravikumar et al have presented evidence of a step change in M<sub>eq</sub> across the PE region in the low field-high temperature part of the (H,T) phase space in single crystals of YBCO and 2H-NbSe<sub>2</sub>, respectively. Thus, the sharp transition from the PE peak to a DPE peak in weakly pinned systems (as in Fig.3(a)) may indeed mark the transition from a pinned to an unpinned state of the vortex matter.
Furthermore, the Fig.3(b) shows yet another effect of a larger ac amplitude: a residual PE now visible in contrast with the low amplitude data in Fig.1(a). It is possible that a larger ac signal anneals the dilute lattice thereby making the amorphization more easily detectable. The possible annealing effect of a larger ac field needs further investigations. Studies on samples of LTSC had earlier revealed that specific details of the structure in the $`\chi (T)`$ curve across the PE depend on the amplitude of the ac field.
### C History effects in $`\chi (T)`$ behavior
The history effects pertain to the dependence of $`\chi (H,T)`$ on the path followed in reaching a given H and T. A simple way to explore them is to examine the difference in $`\chi (T)`$ between ZFC and FC modes. Fig.4 displays the $`\chi (T)`$ data in h<sub>ac</sub> of 2 Oe (r.m.s.) at H = 500 Oe ($``$c) (i.e., when the PE peak is well recognizable) for both ZFC and FC modes. The inset in Fig.4 shows the data across the PE region on an expanded scale. Note first that above T<sub>p</sub>, $`\chi (T)`$ behavior is the same for both ZFC and FC modes. However, prior to the entry into the PE region (i.e., for $`T<T_{pl}`$), $`|\chi _{FC}|>|\chi _{ZFC}|`$ . As per eqn.(2) and eqn.(3), this inequality would translate as J$`{}_{}{}^{FC}{}_{c}{}^{}`$ $``$ J$`{}_{}{}^{ZFC}{}_{c}{}^{}`$ for $`T<T_p`$. In terms of Larkin-Ovchinnikov scenario (cf. eqn.(1)), this implies that the correlation volume V<sub>c</sub> in the FC state in larger than that in the ZFC state. This inference agrees with the direct observation of better order for the vortex states (mostly at low fields) prepared in the FC manner as compared to those prepared in the ZFC manner in crystals of YBCO as well as in other high T<sub>c</sub> cuprates . However, the above inequality is in complete contrast to the situation in weakly pinned samples of LTSC systems on which history dependent J<sub>c</sub> data have been reported for long .
In transport studies of single crystals of niobium, Steingart, Putz and Kramer had noted the inequality,
$$J_c^{FC}(H)>J_c^{rev}(H)>J_c^{ZFC}(H),$$
(4)
where $`J_c^{rev}(H)`$ is the current density while reversing the field from above H<sub>c2</sub> to a given H in an isothermal scan. In recent years, the history effects and metastability (supercooling, etc.) has received new impetus . The observations that (i)$`J_c^{FC}(H)>J_c^{ZFC}(H)`$ for $`H<H_p`$ and (ii)$`J_c^{FC}(H)J_c^{ZFC}(H)`$ for $`H>H_p`$ in LTSC systems have led to the suggestion that FC states attempt to freeze in (i.e., supercool) the amorphous correlations present at (and above) the peak position of the PE. A FC state is therefore more strongly pinned than the corresponding ZFC state. The ZFC phase, in contrast, is considered to be prepared by exposing the weakly pinned superconducting sample to an applied field which generates vortices that enter the sample at high velocities and are able to overcome the effects of pinning centers to eventually explore the well ordered stable state of the system. It may be speculated here that while field cooling in the YBCO case, the temperature interval between the irreversibility temperature T<sub>irr</sub> and the peak temperature T<sub>p</sub> is so narrow (cf. Fig.3(a)) that the sample fails to explore the (equilibrium) amorphous phase during a very fast cool down. The vortex density in the sample is uniform in the FC state, whereas in the ZFC case, in particular at low fields, the strong pinning effects near the edges in the platlet shaped samples of YBCO result in the vortex density being non-uniform across the cross-section of the sample. Large thermal energy would help overcome the effects of pinning centers in both cases, however, the FC state being more (macroscopically) uniform eventually produces a better ordered FLL with a larger correlation volume V<sub>c</sub> of the Larkin domain (than that in the ZFC case). We note that Kokkaliaris et al have recently reported that in an untwinned YBCO crystal, which displays the phenomenon of a broad second magnetization peak (to be designated as H$`{}_{}{}^{s}{}_{p}{}^{}`$) at high fields and lower temperatures, $`J_c^{ZFC}(H)<J_c^{rev}(H)`$ for $`H<H_p^s`$. Such an inequality is indeed consistent with the behavior in LTSC systems (cf. eqn.(4)). The role that the width of the PE region could play in deciding the nature of the inequality between J$`{}_{}{}^{FC}{}_{c}{}^{}`$ and J$`{}_{}{}^{ZFC}{}_{c}{}^{}`$ in LTSC/HTSC systems remains to be further investigated.
It is to be noted that regardless of the sign of the history dependence, in all cases reported so far, the T<sub>p</sub> provides an upper limit of the history dependence. Thus, in analogy with disordered magnets such as spin glasses, T<sub>p</sub> marks a characteristic “transition" temperature above which the system is disordered in “equilibrium".
### D Construction of the vortex phase diagram and its comparison with earlier reports
Figure 5 shows a collation of the data corresponding to the PE region marked by the onset (T<sub>pl</sub>) and the peak (T<sub>p</sub>) positions along with the depinning (T<sub>dp</sub>) and the superconducting transition (T<sub>c</sub>) temperatures over the (H,T) region, where the PE phenomenon could be identified. In view of the observation that the sharp depinning transition (i.e., the irreversibility temperature) immediately succeeds the peak of the PE, the T$`{}_{p}{}^{}(H)`$ line in Fig.5 runs parallel to the T$`{}_{dp}{}^{}(H)`$ line. Several reports in the twinned and untwinned crystals of YBCO provide evidence that the T$`{}_{p}{}^{}(H)`$ values are located either in close proximity to the melting temperatures (as determined from the transport studies ) of the underlying pristine FLL or they coincide with the temperatures of a step change in the equilibrium magnetization . One could draw the loci of different features of the PE phenomenon seen in transport and ac susceptibility measurements , however, they all seem to indicate that H-T lines so drawn would conform to the power-law FLL melting relationship :
$$H_m=H_0(1\frac{T}{T_c(0)})^n,$$
(5)
where the prefactor H<sub>0</sub> $``$ 10<sup>6</sup> Oe and n $``$ 1 to 2 . Most of the data on the PE curve or the melting line or the irreversibility line ( in YBCO ) in the literature, which have been found to conform to the power law relationship, are at fields larger than 1 kOe. A broken line satisfying the T$`{}_{dp}{}^{}(H)`$ data points for H $`>`$ 0.25 kOe in Fig.5 attests to the efficacy of the power law behavior. We, however, find that in the low field region (H$`<`$0.5 kOe), the deviations from the power law relationship (H$``$(1-t)<sup>n</sup>) set in a significant manner for both the curves T$`{}_{pl}{}^{}(H)`$ and T$`{}_{p}{}^{}(H)`$.
We recall that a sudden collapse of the irreversibility line between 1 kOe and 0.2 kOe in the twinned and the untwinned crystals of YBCO was reported by Krusin-Elbaum et al from high frequency ac susceptibility measurements. However, in their data, there is no evidence of the occurrence of a PE. Below 200 Oe, the irreversibility line of Krusin-Elbaum et al can be seen to proceed towards the T$`{}_{c}{}^{}(0)`$ value in a power-law manner once again, but with a different value of the exponent n. Following Ref.60, we show in Fig.6, the plot of log H versus log (1-t) for the T<sub>dp</sub> line in our sample of YBCO. Note that the data for 6000 Oe $``$ H $``$ 250 Oe and for 70 Oe $``$ H $``$ 20 Oe could be considered to conform to the power law relationship albeit with different values of the exponent n in the two intervals . In the high field region ( H $`>`$ 0.25 kOe), both the exponent n$``$2 and the prefactor $`H_07.9\times 10^6Oe`$ are comparable to similar values (n$``$2 and $`H_06\times 10^6Oe`$) found by Nishizaki et al in a high quality twinned crystal of YBCO. In between the high field and low field regions, i.e., between 250 Oe and 70 Oe, T<sub>dp</sub> curve shows a nose like turnaround feature. Such a feature appears accentuated in T$`{}_{p}{}^{}(H)`$ and T$`{}_{pl}{}^{}(H)`$ curves, as can be clearly viewed via a replot of the vortex phase diagram in a semi-log manner as shown in the Fig.7. In order to comprehend the genesis of the turnaround in T$`{}_{p}{}^{}(H)`$ curve around 100 Oe, we refer to the $`\chi (T)`$ curves displayed in Figs.1(b) to 1(f). These figures show that the PE develops gradually, followed by a progressively rapid depinning immediately thereafter (cf. data in the inset of Fig.2(b)). A careful examination of the $`\chi (T)`$ curves in Figs.1(b), 1(c) and 1(d) at H = 40 Oe, 70 Oe and 120 Oe shows how the sharpening of the PE features results in an increase in T$`{}_{p}{}^{}(H)`$ ( and T$`{}_{dp}{}^{}(H)`$) values with increasing H. Above about 150 Oe, T$`{}_{p}{}^{}(H)`$ values start to show the usual decrease with increase in H. The shape of the T$`{}_{p}{}^{}(H)`$ curve in Fig.7 is reminiscent of the reentrant characteristic in the PE curve noted first by Ghosh et al in NbSe<sub>2</sub>.
We further note that the broadening and eventual disappearance of the PE phenomenon at fields less than 30 Oe is also similar to the behavior in NbSe<sub>2</sub>. If we assume that the PE curve(s) separates the ordered and disordered phases of vortex matter then the absence of the PE at low fields ($`<`$30 Oe, where a<sub>0</sub> $`>`$ 1$`\mu `$) amounts to the absence of an ordered phase, when the dilute vortices are individually pinned by the quenched random disorder in the atomic lattice. In analogy with the results in NbSe<sub>2</sub> and other LTSC systems , we have chosen to label the vortex phase between the higher field upper portion of the T$`{}_{p}{}^{}(H)`$ curve and the T<sub>dp</sub> line as the pinned amorphous (analogous to the pinned liquid) phase. The (H,T) region below the reentrant leg of the T$`{}_{p}{}^{}(H)`$ curve at low fields has been termed as the ‘reentrant disordered’ ( and filled with dotted lines) in the Fig.7. The reentrant disordered and pinned amorphous regions would overlap near the turnaround in the T$`{}_{p}{}^{}(H)`$ curve. At temperatures above the turnaround feature (t $`>`$ 0.99), the vortex state remains disordered at all fields. However, below this temperature (i.e., t$`<`$0.98) an ordered vortex phase (presumably akin to an elastic/Bragg glass) exists between the high field upper portion of T$`{}_{pl}{}^{}(H)`$ curve and the ‘reentrant disordered’ phase . The vortex phase between T$`{}_{pl}{}^{}(H)`$ and T$`{}_{p}{}^{}(H)`$ curves is expected to be plastically deformed and has therefore been termed as plastic glass , detailed theoretical understanding of which is still lacking. The vortex phase diagram in YBCO (in Fig.5/Fig.7) thus shows a close resemblance with the generic phase diagram for a weakly pinned superconductor drawn by Banerjee et al on the basis of characteristics of the peak effect in the NbSe<sub>2</sub> system.
It is instructive to dwell a little more on the similarities in the present YBCO crystal and the 2H-NbSe<sub>2</sub> sample studied by Ghosh et al . We note that the nose feature in t<sub>p</sub> (=$`\frac{T_p(H)}{T_c(0)}`$) curve occurs at nearly the same (H,t) value of ($``$ 100 Oe, $``$0.99) and the PE also disappears in both the samples at H$`<`$30 Oe. Further, the FLL spacing a<sub>0</sub> near the nose feature in NbSe<sub>2</sub> was comparable to its in-plane penetration depth $`\lambda _{ab}`$ . We find the similar correspondence in YBCO; a<sub>0</sub> at 100 Oe is $``$5.0$`\times `$10<sup>3</sup>Å and the measured value of $`\lambda `$ is $``$ 5.5$`\times `$10<sup>3</sup>Å at t$``$0.99 in another crystal of YBCO grown by flux growth technique. The upper portion of the t<sub>p</sub> curve in 2H-NbSe<sub>2</sub> was fitted (see inset in Fig.2 of Ghosh et al ) to the FLL melting relation with the exponent n=2 and the prefactor $`H_0=\beta _m(c_l^4/G_i)H_{c2}(0)`$, where the various symbols have their usual meaning . It yielded the Lindemann number c<sub>L</sub> in NbSe<sub>2</sub> as 0.17. Following Ghosh et al , we plot H versus (1-t<sub>p</sub>)<sup>2</sup> in YBCO in the inset of Fig.6. A linear fit to the data above 500 Oe in this inset yields $`c_L0.25`$ (where H$`{}_{c2}{}^{}(0)`$ $``$ 150T, $`\beta _m=5.6`$, $`G_i=10^2`$ ) The disappearance of the PE at low fields is rationalized by invoking the notion that the pinning length L<sub>c</sub> and the entanglement length L<sub>E</sub> become comparable (cf. eqn.(6.47) of ) at H $``$ 30 Oe. Recalling once again the same relation, the L<sub>c</sub>/L<sub>E</sub> is given as,
$$\frac{L_c}{L_E}\frac{\pi \kappa ^2ln\kappa }{\sqrt{2}}(\frac{a_0}{2\pi \lambda (0)})^2(\frac{j_c}{G_ij_0})^{\frac{1}{2}}\frac{(1t)^{\frac{4}{3}}}{t},$$
$`(7)`$
where the various symbols have their usual meaning . Note, how the L<sub>c</sub>/L<sub>E</sub> scales with the crucial parameters, like, the Ginzburg number G<sub>i</sub>($`=(1/2)(k_BT_c/H_c^2\xi ^3ϵ)^2`$) and the ratio j<sub>c</sub>/j<sub>0</sub>, where j<sub>c</sub> and j<sub>0</sub> represent the critical current density and the depairing current density, respectively. The G<sub>i</sub> in YBCO ($``$ 10<sup>-2</sup>) is expected to be two to three orders larger than that in 2H-NbSe<sub>2</sub> , but it is compensated by the extreme smallness of j<sub>c</sub>/j<sub>0</sub> value ( $``$ 10<sup>-5</sup> to 10<sup>-6</sup>) in typically clean crystals of 2H-NbSe<sub>2</sub> as compared to that ($``$ 10<sup>-2</sup> \- 10<sup>-3</sup>) in a clean YBCO sample . At fields below 30 Oe, where L<sub>c</sub>/L<sub>E</sub>$`<`$1, the dilute vortex array is expected to be so disordered that it does not display any anomalous variation in J<sub>c</sub>(T) corresponding to an order-disorder transformation in an isofield scan. This order of magnitude estimate for the field value of disorder-order transition in FLL correlations appears consistent with isothermal Bitter decoration data in samples of YBCO and BSSCO .
## IV Conclusion
In summary, we have reported the occurrence of the peak effect phenomenon for sparse vortex arrays at very low fields as well as for the dense vortex arrays at moderately high fields in a twinned crystal of YBCO. The data delineate the order-disorder transformation regions at low fields, where the interaction is weaker and the disorder effects of the thermal fluctuations and/or small bundle (or individual) pinning dominate. The new observations at low fields agree with an earlier report of a collapse of the irreversibility line in the (H,T) region close to T$`{}_{c}{}^{}(0)`$ in the context of the YBCO system. But, in addition, it shows a close connection with the “nose" like turnaround feature in the locus of the peak effect in a LTSC 2H-NbSe<sub>2</sub> system . The notion of the thermal softening of the vortex cores has been invoked for YBCO, whereas a connection with the reentrant melting/reentrant disordered phase was proposed for 2H-NbSe<sub>2</sub>. Both ideas imply a transformation/crossover from an interaction dominated collective pinning behavior to a different regime, where pinning/depinning of vortex lines have to be considered in isolation. The observation of the complexity in the PE regime (at intermediate fields), which comprises the onset of a sudden shrinkage in the correlation volume and the internal structure, reinforces the notion of a stepwise loss of vortex correlations. Carruzzo and Yu have recently theoretically proposed a two step softening process of the FLL, in which a first order transition initially transforms an ordered vortex solid to a (defective) soft vortex solid that has smaller, but, still finite shear modulus. The latter solid in turn could undergo another first order transition to the vortex liquid state. Our present results call for more detailed study on the PE in different systems to explore the loss in order of the FLL. The presence of twin boundaries in YBCO is usually considered to be detrimental to the observation of a first order melting transition. However, in our case, the presence of twin boundaries has facilitated the observation of the peak effect at low fields. Further, a recent study of the first order FLL melting transition on the twinned crystals of NdBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> also shows the usefulness of the twinned crystals to explore basic issues in vortex state physics. As a final remark, we recall that Grier et al and Horiuchi et al have directly demonstrated how the disordered dilute vortex arrays undergo transformation to an ordered phase as the interaction effect increases due to an increase in the vortex density either by an increase in the applied field or, more interestingly, by the collation of the vortices around a strong pinning center . Such studies in samples of YBCO and NbSe<sub>2</sub>, which show reentrant characteristic in the peak effect curve are highly desirable.
Acknowledgements
It is a pleasure to acknowledge numerous discussions with Satyajit Banerjee, Shampa Sarkar, Mahesh Chandran, Nandini Trivedi, G. Ravi Kumar, Prasant Mishra, Vinod Sahni and Gautam Menon.
FIGURE CAPTIONS
Fig. 1: Temperature dependence of the normalized values of the in-phase ac susceptibility (4$`\pi \chi (T)`$) in a twinned single crystal of YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> measured with an amplitude (h<sub>ac</sub>) of 0.5 Oe (r.m.s.) and at a frequency of 211 Hz in the dc fields ($``$c) as indicated in the different panels. The insets in the Figs. 1(a) and 1(f) show the $`\chi (T)`$ response in (nominal) zero field and 250 Oe, respectively. The onset (T<sub>pl</sub>) and peak (T<sub>p</sub>) temperatures of the PE have been marked in different panels. The inset panel in Fig.1(b) shows the plot of d$`\chi `$/dT vs T in H = 40 Oe. The temperature of maximum in d$`\chi `$/dT marks the mid point of the depinning transition (T<sub>dp</sub>). An unlabelled arrow in the Fig. 1(a) identifies the temperature across which a remanence of the PE phenomenon can be identified in the $`\chi (T)`$ curve at H = 20 Oe measured with a higher h<sub>ac</sub> of 2 Oe (r.m.s.), as shown in the Fig.3(b).
Fig. 2: Plots of 4$`\pi \chi `$ vs T in higher dc fields. The inset panels in Fig.2(a) and Fig.2(c) help to view the occurrence of the PE at H = 500 Oe and 2500 Oe, respectively. The inset in Fig.2(d) shows the location of the T<sub>dp</sub> via a plot of d$`\chi `$/dT versus T at H = 6000 Oe. The inset in Fig.2(b) shows a plot of the parameter $`\mathrm{\Delta }`$T<sub>dp</sub> versus field (see text for details).
Fig. 3: The $`\chi (T)`$ measured in a h<sub>ac</sub> of 2 Oe(r.m.s.) and at a frequency of 211 Hz in (a) 320 Oe and (b) 20 Oe. An enlarged plot in the inset of panel (a) shows the presence of the differential paramagnetic effect (DPE), which follows the occurrence of the peak effect. The inset of panel (b) shows an inflection feature in the $`\chi (T)`$ curve at H = 20 Oe.
Fig. 4: Comparison of the $`\chi (T)`$ behavior (h<sub>ac</sub> = 2 Oe (r.m.s.) and f = 211 Hz) for the vortex states prepared at 500 Oe in the zero field cooled (ZFC) and the field cooled (FC) modes. The inset shows an enlarged view of the data across T<sub>p</sub>.
Fig. 5: The vortex phase diagram in the twinned YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> crystal (for H$``$c) constructed on the basis of the data on the onset (T<sub>pl</sub>), the peak (T<sub>p</sub>) and the end (i.e., the depinning temperature T<sub>dp</sub>) of the peak effect phenomenon. The schematically drawn normal state boundary (i.e., the dotted H<sub>c2</sub> line) corresponds to the (nominal) zero field transition temperature of 93.3 K and the T$`{}_{c}{}^{}(H)`$ values ascertained at high fields (H $`>`$ 0.25 kOe). The dashed line notionally marks the upper end of the pinned amorphous region and it corresponds to the relationship, H = 7.9$`\times `$10<sup>6</sup>(1-$`\frac{T}{T_c(0)}`$)<sup>2</sup> Oe. For a justification of the nomenclature of the different vortex phases, see text.
Fig. 6: Log-log plot of H vs (1-t), where $`t=\frac{T_{dp}}{T_c(0)}`$. The solid lines correspond to the power law \[ $`HH_0(1t)^n`$ \], above and below the field-temperature region of a collapse (or turn around) in the depinning line. The high field ( H $`>`$ 0.3 kOe) exponent is n $``$ 2 and the prefactor $`H_07.9\times 10^6Oe`$. The inset shows a plot of H vs (1-t<sub>p</sub>)<sup>2</sup>, where $`t_p=\frac{T_p}{T_c(0)}`$. T$`{}_{p}{}^{}(H)`$ data yield prefactor H<sub>0</sub> $``$ 3.7$`\times `$10<sup>6</sup> Oe.
Fig. 7: The same data as in Fig.5 on a semilog plot illustrating the shape of the t$`{}_{p}{}^{}(H)`$ curve at low fields ( H $`<`$ 300 Oe). As in Fig.5, the dashed line conforms to the power law relation with exponent n $``$ 2. The solid line satisfying the t$`{}_{p}{}^{}(H)`$ data points has been drawn in a free hand manner, it can be seen to run parallel to the dashed line for H $`>`$ 300 Oe. The dotted line joining $`t_{pl}(H)`$ data points is a guide to the eye. |
warning/0001/hep-th0001063.html | ar5iv | text | # References
January 2000
UTMS 2000-4
UTHEP-416
hep-th/0001063
Modular invariance of string theory on $`AdS_3`$
Akishi Kato <sup>2</sup><sup>2</sup>2akishi@ms.u-tokyo.ac.jp and Yuji Satoh <sup>3</sup><sup>3</sup>3ysatoh@het.ph.tsukuba.ac.jp
Graduate School of Mathematical Science, University of Tokyo
Komaba, Meguro-ku, Tokyo 153-8914, Japan
Institute of Physics, University of Tsukuba
Tsukuba, Ibaraki 305-8571, Japan
Abstract
We discuss the modular invariance of the $`SL(2,R)`$ WZW model. In particular, we discuss in detail the modular invariants using the $`\widehat{sl}(2,R)`$ characters based on the discrete unitary series of the $`SL(2,R)`$ representations. The explicit forms of the corresponding characters are known when no singular vectors appear. We show, for example, that from such characters modular invariants can be obtained only when the level $`k<2`$ and infinitely large spins are included. In fact, we give a modular invariant with three variables $`Z(z,\tau ,u)`$ in this case. We also argue that the discrete series characters are not sufficient to construct a modular invariant compatible with the unitarity bound, which was proposed to resolve the ghost problem of the $`SL(2,R)`$ strings.
1 Introduction
The string theory on $`AdS_3`$, namely, $`SL(2,R)`$ is important in various respects and it has been investigated for more than a decade -. It gives us the simplest string model in backgrounds with curved time. Without R-R charges, the system is described by the $`SL(2,R)`$ WZW model. This WZW model is one of the simplest models of the non-compact CFT. In addition, it is closely related to the string theory in some two ($`SL(2,R)/U(1)`$) and three (BTZ) dimensional black hole backgrounds. The appearance of the AdS/CFT correspondence aroused the renewed interest in the $`SL(2,R)`$ WZW model and its Euclidean version, the $`SL(2,C)/SU(2)`$ WZW model, e.g., -.
However, in spite of recent intensive studies, there still remain open questions for these string models at the fundamental level. Such a state of the problem was discussed in . In fact, soon after the study of the $`SL(2,R)`$ WZW model was initiated, it was found that the model contains negative-norm physical states (ghosts) . So far, there are two types of the proposals for the resolution. In one proposal , the discrete unitary series of the $`SL(2,R)`$ representations is used and it is claimed that ghosts can be removed if the spectrum is truncated so that the spin $`j`$ and the level $`k`$ of the current algebra $`\widehat{sl}(2,R)`$ satisfy
$`{\displaystyle \frac{1}{2}}j<{\displaystyle \frac{k}{2}},k>2.`$ (1.1)
(For details of $`\widehat{sl}(2,R)`$ and our conventions, see the next section.) The condition (1.1) is called the unitarity bound. In the other proposal , the principal continuous series and the free field representations of the current algebra are used. However, in both cases, it is still an open question if such proposals are compatible with other consistency conditions of string theory such as the modular invariance and the closure of the operator product expansion. This is because such consistency conditions are not well understood either. Thus it seems important to accumulate precise knowledge about them.
In the literature, there were several arguments about this issue. Regarding the OPE, see -. On the modular invariance, for instance, the modular invariants were constructed for the so-called admissible representations in . In , it was argued that modular invariants can be obtained for the discrete unitary series by adding the sectors corresponding to some winding modes. In , the modular invariants of the $`N=2`$ SCFT were expressed in terms of the $`\widehat{sl}(2,R)`$ characters based on the relationship between the $`N=2`$ superconformal algebra and $`\widehat{sl}(2,R)`$ . In , modular invariants for the principal continuous series were discussed along the line of .
In this paper, we discuss in detail the possibility of constructing modular invariants using the discrete series characters, which was also discussed in . We will not include additional sectors as in . In the next section, we review the basics of $`SL(2,R)`$ and $`\widehat{sl}(2,R)`$ and give the modular transformations of the discrete series characters when no singular vectors appear. In section 3, we discuss the case in which finite number of the characters are used. In section 4, we discuss the case in which infinite number of the characters without singular vectors are used. In section 5, we give a summary of our results and discuss the implication to the unitarity bound (1.1).
2 Setup
The $`SL(2,R)`$ current algebra is defined by the commutation relations
$`[J_n^0,J_m^0]`$ $`=`$ $`{\displaystyle \frac{1}{2}}kn\delta _{n+m},[J_n^0,J_m^\pm ]=\pm J_{n+m}^\pm ,`$
$`[J_n^+,J_m^{}]`$ $`=`$ $`2J_{n+m}^0+kn\delta _{n+m}.`$ (2.1)
The generators of the associated Virasoro algebra are given by
$`L_n`$ $`=`$ $`{\displaystyle \frac{1}{k2}}{\displaystyle \underset{m𝐙}{}}:{\displaystyle \frac{1}{2}}J_{nm}^+J_m^{}+{\displaystyle \frac{1}{2}}J_{nm}^{}J_m^+J_{nm}^0J_m^0:,`$ (2.2)
and its central charge is
$`c`$ $`=`$ $`{\displaystyle \frac{3k}{k2}}.`$ (2.3)
The current algebra $`\widehat{sl}(2,R)`$ contains zero mode subalgebra $`SL(2,R)`$ generated by $`J_0^0`$ and $`J_0^\pm `$. In particular, for a given $`\widehat{sl}(2,R)`$ representation $`V`$, its ground state subspace $`V^0\{vV|L_0v=\mathrm{\Delta }_jv\}`$ with $`\mathrm{\Delta }_j=\frac{j(j+1)}{k2}`$ naturally provides a representation of $`SL(2,R)`$. The operators $`\stackrel{}{J}^2`$ and $`J_0^0`$ act on the states in $`V^0`$ as
$`\stackrel{}{J}^2|j,m=j(j+1)|j,m,`$ $`J_0^0|j,m=m|j,m.`$ (2.4)
The other states in $`V`$ are obtained by acting on $`|j,mV^0`$ with $`J_n^a`$ $`(n0)`$.
On physical grounds, we are interested in the $`\widehat{sl}(2,R)`$ representations $`V`$ where $`V^0`$ are irreducible and unitary as $`SL(2,R)`$ modules. In such cases, since $`j(j+1)`$ is real and invariant under $`jj1`$, one can assume either $`j1/2`$ or $`j=1/2+i\lambda ,\lambda 0`$ without loss of generality.
For the universal covering group of $`SL(2,R)`$, there are five classes of such representations:
(1) Identity representation: the trivial representation with $`\stackrel{}{J}^2=J_0^0=0`$.
(2) Principal continuous series: representations with $`m=m_0+n,\mathrm{\hspace{0.17em}0}m_0<1,n𝐙`$
and $`j=1/2+i\lambda ,\lambda >0`$.
(3) Supplementary series: representations with $`m=m_0+n,\mathrm{\hspace{0.17em}0}m_0<1,n𝐙`$ and
min$`\{m_0,m_01\}<j1/2`$.
(4) Highest weight discrete series ($`𝒟_{\mathrm{hw}}`$): representations with $`m=M_{\mathrm{max}}n,n=0,1,2,\mathrm{},j=M_{\mathrm{max}}1/2`$ and the highest weight state satisfying $`J_0^+|j,j=0`$.
(5) Lowest weight discrete series ($`𝒟_{\mathrm{lw}}`$) representations with $`m=M_{\mathrm{min}}+n,n=0,1,2,\mathrm{},j=M_{\mathrm{min}}1/2`$ and the lowest weight state satisfying $`J_0^{}|j,j=0`$.
If we do not take the universal covering group, the parameters are restricted to $`m_0=0,1/2`$ in (2), $`m_0=0`$ in (3) and $`j=`$ (half integers) in (4) and (5).
In the following we will focus on the $`\widehat{sl}(2,R)`$ representations whose ground state subspace corresponds to the discrete series of the type ($`𝒟_{\mathrm{hw}}`$) or ($`𝒟_{\mathrm{lw}}`$); they will be denoted by $`V_j^{\mathrm{hw}}`$ or $`V_j^{\mathrm{lw}}`$, respectively.
For, e.g., the $`SU(2)`$ current algebra, the characters are naturally defined using three variables (see, for example, ). With this in mind, we define the irreducible characters by
$`\mathrm{ch}_j(z,\tau ,u)`$ $``$ $`e^{2\pi iku}{\displaystyle \underset{V_j}{}}e^{2\pi iJ_0^0z}e^{2\pi i\tau (L_0\frac{c}{24})}.`$ (2.5)
where the summation is taken over $`V_j`$ which is the irreducible representation of the type of either $`V_j^{\mathrm{hw}}`$ or $`V_j^{\mathrm{lw}}`$. The plus sign in the first factor $`e^{+2\pi iku}`$ is due to the change $`kk`$ compared with the compact case. Since the $`SL(2,R)`$ current module has infinite degeneracy with respect to $`L_0`$, it is inevitable to keep $`z0`$ for the convergence of the charcaters. In addition, it turn out that, as functions of the only two variables $`(\tau ,z)`$, one cannot obtain proper modular transformations of the characters. Thus the use of the three variables is essential.
To calculate these characters, one needs to know about singular vectors. It is known that the $`\widehat{sl}(2,R)`$ Verma module at level $`k`$ with the highest weight $`|j,j`$ has singular vectors if and only if at least one of the following conditions is satisfied :
$`(1)`$ $`2j+1=s+(k2)(r1),`$
$`(2)`$ $`2j+1=sr(k2),`$ (2.6)
$`(3)`$ $`k2=0,`$
where $`r,s`$ are positive integers. For example, in the case of $`k2>0`$ and $`j1/2`$, the singular vectors appear when $`2j+1=sr(k2)`$.
Thus, for generic values of $`j`$, there are no singular vectors in $`V_j^{\mathrm{hw}}`$; the irreducible character $`\mathrm{ch}_j`$ coincides with that of the Verma module:
$`\chi _\mu ^{\mathrm{hw}}(z,\tau ,u)`$ $``$ $`e^{2\pi iku}e^{2\pi i\mu z}q^{\frac{\mu ^2}{k2}}i\vartheta _1^1(z|\tau ).`$ (2.7)
Here $`q=e^{2\pi i\tau }`$, $`\mu =j+1/2`$ and
$`\vartheta _1(z|\tau )`$ $`=`$ $`2q^{1/8}\mathrm{sin}(\pi z){\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1q^n)(1q^ne^{2\pi iz})(1q^ne^{2\pi iz}).`$ (2.8)
Similarly, for $`V_j^{\mathrm{lw}}`$ we have $`\chi _\mu ^{\mathrm{lw}}(z,\tau ,u)=\chi _\mu ^{\mathrm{hw}}(z,\tau ,u)`$.
So far we have considered $`𝒟_{\mathrm{hw}}`$ and $`𝒟_{\mathrm{lw}}`$ with $`\mu j+1/20`$ separately. However with the help of the symmetry $`\chi _\mu ^{\mathrm{lw}}(z,\tau ,u)`$ $`=`$ $`\chi _\mu ^{\mathrm{hw}}(z,\tau ,u)=\chi _\mu ^{\mathrm{hw}}(z,\tau ,u)`$, we extend the range of $`\mu `$ as $`\mathrm{}<\mu <+\mathrm{}`$ and drop the superscript $`\mathrm{hw}`$ with the following convention:
$`\chi _\mu (z,\tau ,u)\text{ r.h.s. of (}\text{2.7}\text{}=\{\begin{array}{cc}+\chi _\mu ^{\mathrm{hw}}(z,\tau ,u),\hfill & (\mu 0)\hfill \\ \chi _\mu ^{\mathrm{lw}}(z,\tau ,u).\hfill & (\mu 0)\hfill \end{array}`$ (2.11)
We remark that one cannot consider the specialized characters $`\chi _\mu (0,\tau ,0)`$ since they diverge in the limit $`z0`$ because of the infinite degeneracy with respect to $`L_0`$.
For a special values of $`\mu =j+1/2`$ for which the Verma module has singular vectors, the irreducible character $`\mathrm{ch}_j`$ is different from $`\chi _\mu `$. In such a case, the explicit form of $`\mathrm{ch}_j`$ seems unknown except for several cases.
In our normalization of $`(z,\tau ,u)`$, the modular transformations are generated by
$`S:`$ $`(z,\tau ,u)({\displaystyle \frac{z}{\tau }},{\displaystyle \frac{1}{\tau }},u+{\displaystyle \frac{z^2}{4\tau }}),`$
$`T:`$ $`(z,\tau ,u)(z,\tau +1,u).`$ (2.12)
Under $`T`$-transformation, the characters $`\chi _\mu `$ just get phases,
$`\chi _\mu (z,\tau +1,u)`$ $`=`$ $`e^{2\pi i\left(\frac{\mu ^2}{k2}+\frac{1}{8}\right)}\chi _\mu (z,\tau ,u).`$ (2.13)
For $`k2<0`$, the $`S`$-transformation of $`\chi _\mu (z,\tau ,0)`$ is given in . In our case with three variables, it reads as
$`\chi _\mu ({\displaystyle \frac{z}{\tau }},{\displaystyle \frac{1}{\tau }},u+{\displaystyle \frac{z^2}{4\tau }})`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{2k}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\nu e^{4\pi i\frac{\mu \nu }{k2}}\chi _\nu (z,\tau ,u).`$ (2.14)
For $`k2>0`$, the right-hand side of (2.14) does not make sense since $`\mathrm{\Delta }_j\mathrm{}`$ as $`\nu \pm \mathrm{}`$ and the integral diverges. Instead, after some calculation, we get a slightly different formula which does converge:
$`\chi _\mu ({\displaystyle \frac{z}{\tau }},{\displaystyle \frac{1}{\tau }},u+{\displaystyle \frac{z^2}{4\tau }})`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{k2}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\nu e^{4\pi \frac{\mu \nu }{k2}}\chi _{i\nu }(z,\tau ,u).`$ (2.15)
Note that an imaginary $`\mu =j+1/2`$ corresponds to a spin of the principal continuous series but $`\chi _{i\mu }`$ are not the characters for those representations any more.
3 Modular invariants from finite number of discrete series characters
Now let us start the discussion of the modular invariance. In this section, we will discuss the possibility of constructing the modular invariants using finite number of the discrete series characters.
First, we would like to discuss the modular invariants using the characters in generic cases, $`\chi _\mu `$. We then show that it is impossible to make modular invariants from finite number of $`\chi _\mu `$; more precisely, a finite sum of the left and right characters
$`Z(z,\tau ,u)`$ $``$ $`{\displaystyle \underset{\mu ,\mu ^{}}{}}N_{\mu ,\mu ^{}}\chi _\mu (z,\tau ,u)\chi _\mu ^{}^{}(z,\tau ,u),`$ (3.1)
with non-zero coefficients $`N_{\mu ,\mu ^{}}`$ cannot be modular invariant. The argument is a simple application of Cardy’s for $`c>1`$ CFT . We compute in two different ways the leading behavior of $`Z`$ in the limit $`\tau +i0`$ with $`|z/\tau |`$ fixed. On one hand, using the $`S`$-transform of $`\vartheta _1`$, one finds
$`Z(z,\tau ,u)`$ $`=`$ $`e^{4\pi k\mathrm{Im}u}e^{2\pi \mathrm{Im}\frac{z^2}{\tau }}|\tau \vartheta _1^2({\displaystyle \frac{z}{\tau }}|{\displaystyle \frac{1}{\tau }})|{\displaystyle }N_{\mu ,\mu ^{}}e^{2\pi i(\mu z\mu ^{}z^{})}q^{\frac{\mu ^2}{k2}}(q^{})^{\frac{\mu ^2}{k2}},`$ (3.2)
$``$ $`e^{4\pi k\mathrm{Im}u}\left({\displaystyle N_{\mu ,\mu ^{}}}\right)|{\displaystyle \frac{\tau }{4}}\mathrm{sin}^2(\pi z/\tau )|\stackrel{~}{q}^{1/4},`$
where $`\stackrel{~}{q}=e^{2\pi i/\tau }`$. On the other hand, if the partition function is modular invariant, one should get
$`Z(z,\tau ,u)`$ $`=`$ $`Z({\displaystyle \frac{z}{\tau }},{\displaystyle \frac{1}{\tau }},u+{\displaystyle \frac{z^2}{4\tau }})`$
$``$ $`e^{4\pi k\mathrm{Im}u}e^{2\pi i(\mu _0\frac{z}{\tau }\mu _0^{}(\frac{z}{\tau })^{})}N_{\mu _0,\mu _0^{}}{\displaystyle \frac{1}{4}}|\mathrm{sin}^2(\pi z/\tau )|\stackrel{~}{q}^{\frac{\mu _0^2+\mu _0^2}{k2}\frac{1}{4}},`$
where the pair $`(\mu _0,\mu _0^{})`$ is chosen so that $`(\mu _0^2+\mu _0^2)/(k2)`$ takes the maximum value there. Clearly the two behaviors (3.2) and (S3.Ex6) cannot be compatible and hence the statement was shown.
Next, we would like to consider more general cases including the representations with singular vectors. In such cases, the characters are not always given by $`\chi _\mu `$ and we need to consider the modular invariant partition functions in terms of the irreducible characters $`\mathrm{ch}_j`$ rather than $`\chi _\mu `$:
$`\widehat{Z}(z,\tau ,u)`$ $``$ $`{\displaystyle \underset{\mu ,\mu ^{}}{}}N_{\mu ,\mu ^{}}\mathrm{ch}_j(z,\tau ,u)\mathrm{ch}_j^{}^{}(z,\tau ,u),`$ (3.4)
where $`N_{\mu ,\mu ^{}}>0`$.
We now show that, for $`k>2`$, it is impossible to construct modular invariants from finite number of the characters for $`V_j^{\mathrm{hw}}`$ (or $`V_j^{\mathrm{lw}}`$) only. We prove only the case of $`V_j^{\mathrm{hw}}`$ since the argument for $`V_j^{\mathrm{lw}}`$ is similar. To this end, we decompose the character $`\mathrm{ch}_j^{\mathrm{hw}}(z,\tau ,u)`$ in terms of the variable $`z`$ :
$`\mathrm{ch}_j^{\mathrm{hw}}(z,\tau ,u)`$ $`=`$ $`e^{2\pi iku}e^{2\pi ijz}q^{\frac{\mu ^2}{k2}\frac{1}{8}}{\displaystyle \underset{p𝐙}{}}e^{2\pi ipz}\mathrm{ch}_{j,p}(q).`$ (3.5)
Note that $`\mathrm{ch}_{j,p}(q)`$ is a power series in $`q`$ with non-negative integer coefficients. In the generic case when there are no singular vectors, i.e., $`\mathrm{ch}_j^{\mathrm{hw}}=\chi _\mu `$, we denote the above expansion by
$`\chi _\mu (z,\tau ,u)`$ $`=`$ $`e^{2\pi iku}e^{2\pi ijz}q^{\frac{\mu ^2}{k2}\frac{1}{8}}{\displaystyle \underset{p𝐙}{}}e^{2\pi ipz}\chi _p(q).`$ (3.6)
Comparing this with (2.7) gives
$`{\displaystyle \underset{p𝐙}{}}e^{2\pi ipz}\chi _p(q)`$ $`=`$ $`e^{\pi iz}q^{1/8}i\vartheta _1^1(z|\tau ).`$ (3.7)
From the explicit form of the theta function (2.8), one finds that the expansion (3.6) converges absolutely (at least) for
$`0<\mathrm{Im}z<\mathrm{Im}\tau ,\mathrm{Re}z=\mathrm{Re}\tau =\mathrm{\hspace{0.17em}0}.`$ (3.8)
Since any irreducible character $`\mathrm{ch}_j`$ is obtained by subtracting from $`\chi _\mu `$ the contribution from singular vectors, we obtain the inequality
$`0<\mathrm{ch}_{j,p}(q)`$ $``$ $`\chi _p(q),`$ (3.9)
if $`\mathrm{Re}\tau =0,\mathrm{Im}\tau >0`$. Therefore one can also evaluate the expansion (3.5) in the region (3.8) and finds that
$`0<e^{2\pi iku}\mathrm{ch}_j^{\mathrm{hw}}(z,\tau ,u)`$ $``$ $`e^{2\pi iku}\chi _\mu (z,\tau ,u).`$ (3.10)
From the inequality (3.10), it follows that in the region (3.8)
$`\widehat{Z}(z,\tau ,u)`$ $``$ $`Z(z,\tau ,u).`$ (3.11)
If the partition function is modular invariant, i.e., $`\widehat{Z}(z,\tau ,u)=\widehat{Z}(\frac{z}{\tau },\frac{1}{\tau },u+\frac{z^2}{4\tau })`$, in the limit $`\tau +i0`$ with $`|z/\tau |`$ fixed, the above inequality gives
$`e^{2\pi i(\mu _0\mu _0^{})z/\tau }F\left({\displaystyle \frac{z}{\tau }}\right)N_{\mu _0,\mu _0^{}}\stackrel{~}{q}^{\frac{\mu _0^2+\mu _0^2}{k2}\frac{1}{4}}`$ $``$ $`{\displaystyle \frac{1}{4}}\left({\displaystyle N_{\mu ,\mu ^{}}}\right)\mathrm{sin}^2(\pi z/\tau )|\tau |\stackrel{~}{q}^{1/4},`$ (3.12)
where $`F(z/\tau )`$ is a function of $`e^{2\pi iz/\tau }`$ which remains finite in the limit. The inequality (3.12) cannot be satisfied for $`k>2`$ and hence the statement was shown.
The above argument cannot be generalized to the case in which the partition function includes the characters for both $`V_j^{\mathrm{hw}}`$ and $`V_j^{\mathrm{lw}}`$. This is because the series in (2.5) can be defined in $`\mathrm{Im}z0`$ for $`V_j^{\mathrm{hw}}`$ whereas it can be defined in $`\mathrm{Im}z0`$ for $`V_j^{\mathrm{lw}}`$ and hence we cannot find the region where the partition function becomes real such as (3.8). In this case, the partition function is not analytic in $`z`$ unless the analytic continuation is possible. To further discuss this case, we may need to find the explicit expressions of the irreducible characters $`\mathrm{ch}_j^{\mathrm{hw}(\mathrm{lw})}`$.
One might wonder also if a similar statement holds for $`k<2`$. To check this, let us note that the arguments in this section hold for more generic highest and lowest weight representations besides $`V_j^{\mathrm{hw}}`$ and $`V_j^{\mathrm{lw}}`$ since we did not use any specific properties of the unitary representations $`𝒟_{\mathrm{hw}}`$ and $`𝒟_{\mathrm{lw}}`$. However, for $`k<2`$, modular invariants using finite number of the characters are actually known for the admissible representations. Thus a simple inequality (3.9) should not exclude the possibility of modular invariants for $`k<2`$.
4 Modular invariants from infinite number of discrete series characters
We saw that the possibility of constructing modular invariants from finite number of the discrete series characters is quite limited. In this section we consider whether the infinite sum or the integral of the characters $`\chi _\mu `$ leave the possibility of constructing modular invariants.
As we saw in section 2, modular $`S`$-transformation of $`\chi _\mu `$ is expressed as a superposition of infinitely many characters. However unlike the momentum eigenstates (plane waves) in flat space, it is not clear in what sense $`\widehat{sl}(2,R)`$ characters $`\{\chi _\mu \}_{\mu R}`$ are orthogonal or linearly independent. In particular, the formula for the modular $`S`$-transformation of $`\chi _\mu `$ might not be unique. In fact, it may be possible to get different expressions by deforming the integration contours in (2.14) and (2.15). Thus let us discuss the possibility of the modular invariants in a definite manner.
Here we would like to show that the $`S`$-transformation of $`\chi _\mu `$ cannot be expressed by using the values of $`\mu `$ belonging only to a finite interval, namely, by using $`\chi _\mu `$ with $`\mu [\mu _1,\mu _2]`$ . First, let us suppose that
$`\chi _\mu ({\displaystyle \frac{z}{\tau }},{\displaystyle \frac{1}{\tau }},u+{\displaystyle \frac{z^2}{4\tau }})`$ $`=`$ $`{\displaystyle _{\nu _1}^{\nu _2}}𝑑\nu f(\nu ;\mu ,k)\chi _\nu (z,\tau ,u),`$ (4.1)
where $`f(\nu ;\mu ,k)`$ is some function which is continuous in $`\nu [\nu _1,\nu _2]`$ and independent of $`\tau `$. Using the $`S`$-transform of $`\vartheta _1`$, one finds that (4.1) is modular invariant only if
$`ie^{\pi i(k2)z^2/(2\tau )}e^{2\pi i\mu z/\tau }\stackrel{~}{q}^{\frac{\mu ^2}{k2}}`$ $`=`$ $`\sqrt{i\tau }{\displaystyle _{\nu _1}^{\nu _2}}𝑑\nu f(\nu ;\mu ,k)e^{2\pi i\nu z}q^{\frac{\nu ^2}{k2}}.`$ (4.2)
This can never happen. Indeed, put $`z=0`$ and take the limit $`\mathrm{Im}\tau +\mathrm{}`$. The left hand side tends to 1 whereas the right hand side diverges if $`k>2`$ and tends to 0 if $`k<2`$.
Similar arguments hold in the above and the following when the integral $`_{\nu _1}^{\nu _2}𝑑\nu f(\nu ;\mu ,k)`$ is replaced with an infinite sum $`_if(\nu _i;\mu ,k)`$ with $`\nu _i[\nu _1,\nu _2]`$ as long as the sum makes sense. We will omit discussions in such cases.
Now let us consider a partition function of the following form:
$`Z(z,\tau ,u)`$ $``$ $`{\displaystyle _I}𝑑\mu 𝑑\mu ^{}g(\mu ,\mu ^{})\chi _\mu (z,\tau ,u)\chi _\mu ^{}^{}(z,\tau ,u),`$ (4.3)
where $`g(\mu ,\mu ^{})`$ is a weight function continuous on the domain $`I=[\mu _1,\mu _2]\times [\mu _1^{},\mu _2^{}]`$. It is understood that the measure is zero for the discrete values of $`\mu `$ corresponding to the representations with singular vectors. If such a partition function is modular invariant, it follows that
$`|\tau |{\displaystyle _I}𝑑\mu 𝑑\mu ^{}g(\mu ,\mu ^{})e^{2\pi i(\mu z\mu ^{}z^{})}q^{\frac{\mu ^2}{k2}}(q^{})^{\frac{\mu ^2}{k2}}`$ (4.4)
$`=`$ $`e^{\pi (k2)\mathrm{Im}\frac{z^2}{\tau }}{\displaystyle _I}𝑑\mu 𝑑\mu ^{}g(\mu ,\mu ^{})e^{2\pi i\left(\mu \frac{z}{\tau }\mu ^{}(\frac{z}{\tau })^{}\right)}\stackrel{~}{q}^{\frac{\mu ^2}{k2}}(\stackrel{~}{q}^{})^{\frac{\mu ^2}{k2}}.`$
However, since the interval $`I`$ is finite, a similar argument to the above shows that the equality (4.4) cannot be true. Moreover, a similar argument holds also in the case where $`g(\mu ,\mu ^{})`$ includes distributions (as long as the expression makes sense). Therefore, we have found that it is impossible to construct modular invariants only from $`\chi _\mu `$ with $`\mu `$ belonging to a finite interval $`\mu [\mu _1,\mu _2]`$ even if infinitely many $`\chi _\mu `$ are used.
Thus the modular invariant partition function as a superposition of $`\chi _\mu \chi _\mu ^{}^{}`$ is possible only if the partition function contains arbitrarily high spin $`|\mu |`$. Nevertheless, for $`k>2`$, the partition function becomes ill-defined since $`L_0`$ spectrum is not bounded from below: $`\mathrm{\Delta }_j\mathrm{}`$ as $`|\mu |\mathrm{}`$. Hence we have reached a conclusion that for $`k>2`$ it is impossible to construct modular invariants only from $`\chi _\mu `$, namely, from the discrete series characters without singular vectors (if any procedure such as ‘Wick rotation’ is not consistently implemented).
Finally, we consider the partition function for $`k<2`$ with no upper bound on the spin $`|\mu |`$:
$`Z(z,\tau ,u)`$ $``$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\mu 𝑑\mu ^{}g(\mu ,\mu ^{})\chi _\mu (z,\tau ,u)\chi _\mu ^{}^{}(z,\tau ,u),`$ (4.5)
To analyze this, we introduce the Fourier transform
$`g(\mu ,\mu ^{})`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\xi 𝑑\xi ^{}e^{i(\mu \xi \mu ^{}\xi ^{})}\widehat{g}(\xi ,\xi ^{}).`$ (4.6)
Then for $`k2<0`$ one finds that
$`Z({\displaystyle \frac{z}{\tau }},{\displaystyle \frac{1}{\tau }},u+{\displaystyle \frac{z^2}{4\tau }})`$ $`=`$ $`{\displaystyle \frac{2}{2k}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\nu 𝑑\nu ^{}\widehat{g}({\displaystyle \frac{4\pi \nu }{k2}},{\displaystyle \frac{4\pi \nu ^{}}{k2}})\chi _\nu (z,\tau ,u)\chi _\nu ^{}^{}(z,\tau ,u).`$ (4.7)
Thus a sufficient condition of the modular invariance is
$`g(\mu ,\mu ^{})`$ $`=`$ $`{\displaystyle \frac{2}{2k}}\widehat{g}({\displaystyle \frac{4\pi \mu }{k2}},{\displaystyle \frac{4\pi \mu ^{}}{k2}}).`$ (4.8)
It is easy to find a solution to this condition. It is just the delta-function,
$`g(\mu ,\mu ^{})`$ $`=`$ $`\delta (\mu \mu ^{}),`$ (4.9)
and this gives the diagonal partition function,
$`Z_{\mathrm{diag}}(z,\tau ,u)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\mu |\chi _\mu |^2={\displaystyle _{\mathrm{}}^0}𝑑\mu \left(|\chi _\mu ^{\mathrm{hw}}|^2+|\chi _\mu ^{\mathrm{lw}}|^2\right)`$ (4.10)
$`=`$ $`{\displaystyle \frac{1}{2}}e^{4\pi k\mathrm{Im}u}e^{(2k)\pi \frac{(\mathrm{Im}z)^2}{\mathrm{Im}\tau }}\sqrt{{\displaystyle \frac{2k}{\mathrm{Im}\tau }}}|\vartheta ^2(z|\tau )|.`$
Here we have used the fact that $`\chi _\mu `$ with $`\mu >0`$ are regarded as $`\chi _\mu ^{\mathrm{lw}}`$ with $`\mu <0`$ and recovered the superscripts hw and lw. The diagonal partition function without the variable $`u`$ was discussed in . In our case, it is straightforward to check that $`Z_{\mathrm{diag}}(z,\tau ,u)`$ is actually modular invariant because of the presence of $`u`$.
As pointed out also in , $`Z_{\mathrm{diag}}(z,\tau ,0)`$ was discussed in in the context of a path-integral approach to the $`SL(2,C)/SU(2)`$ WZW model corresponding to Euclidean $`AdS_3`$. This model has an $`\widehat{sl}_2\times \widehat{sl}_2^{}`$ current algebra symmetry and the diagonal partition function may be understood also as the partition function of this model. However, in different spectrum seems to be summed up. It is interesting to consider the precise relationship between the approach here and the one in .
5 Discussion
In this paper, we discussed the modular invariants using the discrete series characters. The arguments hold, except for the last one below Eq.(4.5), even if we set $`u=0`$ and consider $`\mathrm{ch}_j(z,\tau ,0)`$ and $`\chi _\mu (z,\tau ,0)`$. Our arguments indicate that the possibility of constructing modular invariants is very limited.
If we use only the characters without singular vectors, the possibility is only in the case where $`k<2`$ and $`\chi _\mu `$ with $`|\mu |\mathrm{}`$ are included. We gave such an example. The resulting modular invariant $`Z_{\mathrm{diag}}(z,\tau ,u)`$ may be regarded as a kind of twisted partition function. However, its physical interpretation is still unclear, in particular, regarding the factor $`e^{2\pi iku}`$.
In the case in which the characters with singular vectors are allowed, we showed that one cannot construct modular invariants from finite number of the characters for $`V_j^{\mathrm{hw}}`$ or $`V_j^{\mathrm{lw}}`$. To complete the analysis, we may need to obtain the explicit forms of the characters with singular vectors.
Nevertheless, it turns out that the case without singular vectors covers physically interesting cases and gives important implication to the unitarity bound (1.1). This is because the condition of the singular vectors (2.6) implies that there are no singular vectors within (1.1). Furthermore, since the spins $`j`$ in that bound belong to a finite interval, our results indicate that one cannot construct modular invariants only from the discrete series characters for the representations satisfying the unitarity bound (1.1). This means that one cannot make a consistent string theory on $`SL(2,R)=AdS_3`$ only from the spectrum from $`V_j^{\mathrm{hw}(\mathrm{lw})}`$ within (1.1).
Here some comments on the relation to Ref. may be in order. First, we discussed the cases including the representations with singular vectors and our analysis covered the cases of both $`k>2`$ and $`k<2`$. In addition, we defined the character using three variables as in (2.5). As discussed below (2.5), this was essential to obtain a modular invariant (4.10). Since we carried out a systematic analysis using the asymptotic behavior of the characters, we could derive definite conclusions without any ambiguities concerning which characters are linearly independent. Our analysis thus pinned down when the modular invariants can be constructed.
Since there exist ghosts for the discrete series outside the unitarity bound, simply adding such spectrum may not give a consistent theory. Therefore, the possibilities for a consistent theory seem (a) to use the discrete series satisfying (1.1) but include some new sectors with different characters from $`\chi _\mu `$ as in , and/or (b) to use the spectrum of other representations as in . In any case, the string theory on $`AdS_3`$ definitely deserves further investigations.
Note added
While we were proofreading the manuscript, a paper appeared which discusses the modular invariance using the discrete series $`𝒟_{\mathrm{hw}}`$ and $`𝒟_{\mathrm{lw}}`$. In , the diagonal modular invariant (4.10) is obtained from the spectrum satisfying the (more stringent) unitarity bound by (i) including additional sectors along the line of (i.e., the possibility (a) in section 5 in our paper) and by assuming ‘Wick rotation’ (see the comment in section 4). In addition, the role of $`e^{2\pi ik(uu^{})}`$ in $`Z_{\mathrm{diag}}(z,\tau ,u)`$ in our paper is played by the chiral anomaly term $`e^{\pi k\frac{(\mathrm{Im}z)^2}{\mathrm{Im}\tau }}`$ in . A relationship to the $`SL(2,C)/SU(2)`$ case is also discussed there.
The authors of argue that the Hilbert space of the consistent $`SL(2,R)`$ string theory consists of the principal continuous series, $`𝒟_{\mathrm{hw}}`$ and $`𝒟_{\mathrm{lw}}`$ including the sectors generated by the Weyl reflections. This is consistent with the discussion in section 5.
Acknowledgements
We would like to thank I. Bars, J. de Boer, A. Giveon, N. Ishibashi and S.-K. Yang for discussions and P.M. Petropoulos for useful correspondences. We would also like to thank the organizers of Summer Institute ’99 held at Fuji-Yoshida, Japan, 7-21 August, 1999, where this work was started. The work of A.K. is supported in part by the Grant-in-Aid for Scientific Research (No. 11740140) from the Ministry of Education, Science, Sports and Culture of Japan.
References |
warning/0001/hep-ex0001010.html | ar5iv | text | # Study of Two-Body 𝐵 Decays to Kaons and Pions: Observation of 𝐵→𝜋⁺𝜋⁻, 𝐵→𝐾^±𝜋⁰, and 𝐵→𝐾⁰𝜋⁰ Decays.
## Abstract
We have studied charmless hadronic decays of $`B`$ mesons into two-body final states with kaons and pions and observe three new processes with the following branching fractions: $`(B\pi ^+\pi ^{})=(4.3_{1.4}^{+1.6}\pm 0.5)\times 10^6`$, $`(BK^0\pi ^0)=(14.6_{5.13.3}^{+5.9+2.4})\times 10^6`$, and $`(BK^\pm \pi ^0)=(11.6_{2.71.3}^{+3.0+1.4})\times 10^6`$. We also update our previous measurements for the decays $`BK^\pm \pi ^{}`$ and $`B^\pm K^0\pi ^\pm `$.
preprint: CLNS 99/1650 CLEO 99-18
D. Cronin-Hennessy,<sup>1</sup> Y. Kwon,<sup>1,</sup><sup>*</sup><sup>*</sup>*Permanent address: Yonsei University, Seoul 120-749, Korea. A.L. Lyon,<sup>1</sup> E. H. Thorndike,<sup>1</sup> C. P. Jessop,<sup>2</sup> H. Marsiske,<sup>2</sup> M. L. Perl,<sup>2</sup> V. Savinov,<sup>2</sup> D. Ugolini,<sup>2</sup> X. Zhou,<sup>2</sup> T. E. Coan,<sup>3</sup> V. Fadeyev,<sup>3</sup> Y. Maravin,<sup>3</sup> I. Narsky,<sup>3</sup> R. Stroynowski,<sup>3</sup> J. Ye,<sup>3</sup> T. Wlodek,<sup>3</sup> M. Artuso,<sup>4</sup> R. Ayad,<sup>4</sup> C. Boulahouache,<sup>4</sup> K. Bukin,<sup>4</sup> E. Dambasuren,<sup>4</sup> S. Karamnov,<sup>4</sup> S. Kopp,<sup>4</sup> G. Majumder,<sup>4</sup> G. C. Moneti,<sup>4</sup> R. Mountain,<sup>4</sup> S. Schuh,<sup>4</sup> T. Skwarnicki,<sup>4</sup> S. Stone,<sup>4</sup> G. Viehhauser,<sup>4</sup> J.C. Wang,<sup>4</sup> A. Wolf,<sup>4</sup> J. Wu,<sup>4</sup> S. E. Csorna,<sup>5</sup> I. Danko,<sup>5</sup> K. W. McLean,<sup>5</sup> Sz. Márka,<sup>5</sup> Z. Xu,<sup>5</sup> R. Godang,<sup>6</sup> K. Kinoshita,<sup>6,</sup>Permanent address: University of Cincinnati, Cincinnati OH 45221 I. C. Lai,<sup>6</sup> S. Schrenk,<sup>6</sup> G. Bonvicini,<sup>7</sup> D. Cinabro,<sup>7</sup> L. P. Perera,<sup>7</sup> G. J. Zhou,<sup>7</sup> G. Eigen,<sup>8</sup> E. Lipeles,<sup>8</sup> M. Schmidtler,<sup>8</sup> A. Shapiro,<sup>8</sup> W. M. Sun,<sup>8</sup> A. J. Weinstein,<sup>8</sup> F. Würthwein,<sup>8,</sup>Permanent address: Massachusetts Institute of Technology, Cambridge, MA 02139. D. E. Jaffe,<sup>9</sup> G. Masek,<sup>9</sup> H. P. Paar,<sup>9</sup> E. M. Potter,<sup>9</sup> S. Prell,<sup>9</sup> V. Sharma,<sup>9</sup> D. M. Asner,<sup>10</sup> A. Eppich,<sup>10</sup> J. Gronberg,<sup>10</sup> T. S. Hill,<sup>10</sup> D. J. Lange,<sup>10</sup> R. J. Morrison,<sup>10</sup> H. N. Nelson,<sup>10</sup> R. A. Briere,<sup>11</sup> B. H. Behrens,<sup>12</sup> W. T. Ford,<sup>12</sup> A. Gritsan,<sup>12</sup> J. Roy,<sup>12</sup> J. G. Smith,<sup>12</sup> J. P. Alexander,<sup>13</sup> R. Baker,<sup>13</sup> C. Bebek,<sup>13</sup> B. E. Berger,<sup>13</sup> K. Berkelman,<sup>13</sup> F. Blanc,<sup>13</sup> V. Boisvert,<sup>13</sup> D. G. Cassel,<sup>13</sup> M. Dickson,<sup>13</sup> P. S. Drell,<sup>13</sup> K. M. Ecklund,<sup>13</sup> R. Ehrlich,<sup>13</sup> A. D. Foland,<sup>13</sup> P. Gaidarev,<sup>13</sup> L. Gibbons,<sup>13</sup> B. Gittelman,<sup>13</sup> S. W. Gray,<sup>13</sup> D. L. Hartill,<sup>13</sup> B. K. Heltsley,<sup>13</sup> P. I. Hopman,<sup>13</sup> C. D. Jones,<sup>13</sup> D. L. Kreinick,<sup>13</sup> M. Lohner,<sup>13</sup> A. Magerkurth,<sup>13</sup> T. O. Meyer,<sup>13</sup> N. B. Mistry,<sup>13</sup> C. R. Ng,<sup>13</sup> E. Nordberg,<sup>13</sup> J. R. Patterson,<sup>13</sup> D. Peterson,<sup>13</sup> D. Riley,<sup>13</sup> J. G. Thayer,<sup>13</sup> P. G. Thies,<sup>13</sup> B. Valant-Spaight,<sup>13</sup> A. Warburton,<sup>13</sup> P. Avery,<sup>14</sup> C. Prescott,<sup>14</sup> A. I. Rubiera,<sup>14</sup> J. Yelton,<sup>14</sup> J. Zheng,<sup>14</sup> G. Brandenburg,<sup>15</sup> A. Ershov,<sup>15</sup> Y. S. Gao,<sup>15</sup> D. Y.-J. Kim,<sup>15</sup> R. Wilson,<sup>15</sup> T. E. Browder,<sup>16</sup> Y. Li,<sup>16</sup> J. L. Rodriguez,<sup>16</sup> H. Yamamoto,<sup>16</sup> T. Bergfeld,<sup>17</sup> B. I. Eisenstein,<sup>17</sup> J. Ernst,<sup>17</sup> G. E. Gladding,<sup>17</sup> G. D. Gollin,<sup>17</sup> R. M. Hans,<sup>17</sup> E. Johnson,<sup>17</sup> I. Karliner,<sup>17</sup> M. A. Marsh,<sup>17</sup> M. Palmer,<sup>17</sup> C. Plager,<sup>17</sup> C. Sedlack,<sup>17</sup> M. Selen,<sup>17</sup> J. J. Thaler,<sup>17</sup> J. Williams,<sup>17</sup> K. W. Edwards,<sup>18</sup> R. Janicek,<sup>19</sup> P. M. Patel,<sup>19</sup> A. J. Sadoff,<sup>20</sup> R. Ammar,<sup>21</sup> A. Bean,<sup>21</sup> D. Besson,<sup>21</sup> R. Davis,<sup>21</sup> I. Kravchenko,<sup>21</sup> N. Kwak,<sup>21</sup> X. Zhao,<sup>21</sup> S. Anderson,<sup>22</sup> V. V. Frolov,<sup>22</sup> Y. Kubota,<sup>22</sup> S. J. Lee,<sup>22</sup> R. Mahapatra,<sup>22</sup> J. J. O’Neill,<sup>22</sup> R. Poling,<sup>22</sup> T. Riehle,<sup>22</sup> A. Smith,<sup>22</sup> J. Urheim,<sup>22</sup> S. Ahmed,<sup>23</sup> M. S. Alam,<sup>23</sup> S. B. Athar,<sup>23</sup> L. Jian,<sup>23</sup> L. Ling,<sup>23</sup> A. H. Mahmood,<sup>23,</sup><sup>§</sup><sup>§</sup>§Permanent address: University of Texas - Pan American, Edinburg TX 78539. M. Saleem,<sup>23</sup> S. Timm,<sup>23</sup> F. Wappler,<sup>23</sup> A. Anastassov,<sup>24</sup> J. E. Duboscq,<sup>24</sup> K. K. Gan,<sup>24</sup> C. Gwon,<sup>24</sup> T. Hart,<sup>24</sup> K. Honscheid,<sup>24</sup> D. Hufnagel,<sup>24</sup> H. Kagan,<sup>24</sup> R. Kass,<sup>24</sup> J. Lorenc,<sup>24</sup> T. K. Pedlar,<sup>24</sup> H. Schwarthoff,<sup>24</sup> E. von Toerne,<sup>24</sup> M. M. Zoeller,<sup>24</sup> S. J. Richichi,<sup>25</sup> H. Severini,<sup>25</sup> P. Skubic,<sup>25</sup> A. Undrus,<sup>25</sup> S. Chen,<sup>26</sup> J. Fast,<sup>26</sup> J. W. Hinson,<sup>26</sup> J. Lee,<sup>26</sup> N. Menon,<sup>26</sup> D. H. Miller,<sup>26</sup> E. I. Shibata,<sup>26</sup> I. P. J. Shipsey,<sup>26</sup> and V. Pavlunin<sup>26</sup>
<sup>1</sup>University of Rochester, Rochester, New York 14627
<sup>2</sup>Stanford Linear Accelerator Center, Stanford University, Stanford, California 94309
<sup>3</sup>Southern Methodist University, Dallas, Texas 75275
<sup>4</sup>Syracuse University, Syracuse, New York 13244
<sup>5</sup>Vanderbilt University, Nashville, Tennessee 37235
<sup>6</sup>Virginia Polytechnic Institute and State University, Blacksburg, Virginia 24061
<sup>7</sup>Wayne State University, Detroit, Michigan 48202
<sup>8</sup>California Institute of Technology, Pasadena, California 91125
<sup>9</sup>University of California, San Diego, La Jolla, California 92093
<sup>10</sup>University of California, Santa Barbara, California 93106
<sup>11</sup>Carnegie Mellon University, Pittsburgh, Pennsylvania 15213
<sup>12</sup>University of Colorado, Boulder, Colorado 80309-0390
<sup>13</sup>Cornell University, Ithaca, New York 14853
<sup>14</sup>University of Florida, Gainesville, Florida 32611
<sup>15</sup>Harvard University, Cambridge, Massachusetts 02138
<sup>16</sup>University of Hawaii at Manoa, Honolulu, Hawaii 96822
<sup>17</sup>University of Illinois, Urbana-Champaign, Illinois 61801
<sup>18</sup>Carleton University, Ottawa, Ontario, Canada K1S 5B6
and the Institute of Particle Physics, Canada
<sup>19</sup>McGill University, Montréal, Québec, Canada H3A 2T8
and the Institute of Particle Physics, Canada
<sup>20</sup>Ithaca College, Ithaca, New York 14850
<sup>21</sup>University of Kansas, Lawrence, Kansas 66045
<sup>22</sup>University of Minnesota, Minneapolis, Minnesota 55455
<sup>23</sup>State University of New York at Albany, Albany, New York 12222
<sup>24</sup>Ohio State University, Columbus, Ohio 43210
<sup>25</sup>University of Oklahoma, Norman, Oklahoma 73019
<sup>26</sup>Purdue University, West Lafayette, Indiana 47907
CP violation in the Standard Model (SM) arises naturally from the complex phase in the Cabibbo-Kobayashi-Maskawa (CKM) quark-mixing matrix. This picture is supported by numerous experimental constraints, as well as recent observation of direct CP violation in the kaon system, but it remains an open experimental question whether this phase is the only source of CP violation in nature. Studies of the rare charmless decays of $`B`$ mesons are likely to play an important role in constraining the CKM matrix and testing the SM picture of CP violation.
Several approaches have been suggested to extract this phase information from measurements of rare $`B`$ decays. Ratios of various $`BK\pi `$ branching fractions were shown to depend explicitly on $`\gamma Arg(V_{ub}^{})`$ with relatively modest model dependence. Within a factorization model, branching fractions of a large number of rare $`B`$ decays can be parametrized by a small number of independent physical quantities, including $`\gamma `$, which can then be extracted through a global fit to existing data. Finally, measurement of the time-dependent CP-violating asymmetry in the decay $`B^0\pi ^+\pi ^{}`$ can be used to determine the sum of $`\gamma `$ and the phase $`\beta Arg(V_{td}^{})`$. In this case additional measurements of other isospin-related $`B\pi \pi `$ processes are required to allow extraction of $`\gamma +\beta `$ .
In this Letter we present new measurements of $`BK\pi `$ and $`B\pi \pi `$ branching fractions with significantly increased statistics, superseding results from our previous publication. In particular we present first observations of the long-awaited mode $`B\pi ^+\pi ^{}`$, as well as $`BK^\pm \pi ^0`$ and $`BK^0\pi ^0`$ decays.
The data used in this analysis were collected with the CLEO II detector at the Cornell Electron Storage Ring (CESR). It consists of $`9.13\mathrm{fb}^1`$ taken at the $`\mathrm{{\rm Y}}`$(4S), corresponding to 9.66M $`B\overline{B}`$ pairs, and $`4.35\mathrm{fb}^1`$ taken below $`B\overline{B}`$ threshold, used for continuum background studies.
CLEO II is a general purpose solenoidal magnet detector, described in detail elsewhere . Cylindrical drift chambers in a 1.5T solenoidal magnetic field measure momentum and specific ionization ($`dE/dx`$) of charged particles. Photons are detected using a 7800-crystal CsI(Tl) electromagnetic calorimeter. In the CLEO II.V detector configuration, the innermost chamber was replaced by a 3-layer, double-sided silicon vertex detector, and the gas in the main drift chamber was changed from an argon-ethane to a helium-propane mixture. As a result of these modifications, the CLEO II.V portion of the data (2/3 of the total) has significantly improved particle identification and momentum resolution.
Efficient track quality requirements are imposed on charged tracks, and pions and kaons are identified by $`dE/dx`$. The separation between kaons and pions for typical signal momenta $`p2.6`$ GeV$`/c`$ is $`1.7`$ standard deviations ($`\sigma `$) for CLEO II data and $`2.0\sigma `$ for CLEO II.V data. Candidate $`K_S^0`$ are selected from pairs of tracks forming well-measured displaced vertices with a $`\pi ^+\pi ^{}`$ invariant mass within $`2\sigma `$ of the nominal $`K_S^0`$ mass. Pairs of photons with an invariant mass within 2.5$`\sigma `$ of the nominal $`\pi ^0`$ mass are kinematically fitted with the mass constrained to the nominal $`\pi ^0`$ mass. Electrons are rejected based on $`dE/dx`$ and the ratio of the track momentum to the associated shower energy in the CsI calorimeter; muons are rejected based on the penetration depth in the instrumented steel flux return.
The $`B`$ decay candidate is identified via invariant mass and total energy of its decay products. We calculate a beam-constrained $`B`$ mass $`M=\sqrt{E_\mathrm{b}^2p_B^2}`$, where $`p_B`$ is the $`B`$ candidate momentum and $`E_\mathrm{b}`$ is the beam energy. The resolution in $`M`$ is dominated by the beam energy spread and ranges from 2.5 to 3.0 $`\mathrm{MeV}`$, where the larger resolution corresponds to decay modes with a $`\pi ^0`$. We define $`\mathrm{\Delta }E=E_1+E_2E_\mathrm{b}`$, where $`E_1`$ and $`E_2`$ are the energies of the daughters of the $`B`$ meson candidate. The resolution in $`\mathrm{\Delta }E`$ is mode-dependent. For final states without $`\pi ^0`$’s, the $`\mathrm{\Delta }E`$ resolution is $`20`$ MeV ($`25`$ MeV in CLEO II). For modes with $`\pi ^0`$’s the $`\mathrm{\Delta }E`$ resolution is worse by approximately a factor of two and becomes slightly asymmetric because of energy loss out of the back of the CsI crystals. We accept events with $`M`$ within $`5.25.3`$ $`\mathrm{GeV}`$ and $`|\mathrm{\Delta }E|<200`$ MeV. This fiducial region includes the signal region and a generous sideband for background normalization. $`\pi \pi `$ and $`K\pi `$ signal events are distinguished both by $`dE/dx`$ and $`\mathrm{\Delta }E`$ observables. The $`\mathrm{\Delta }E`$ distribution for $`BK^+\pi ^{}`$, calculated under the replacement of $`m_K`$ by $`m_\pi `$, is centered at -42 MeV, giving a separation of $`2.1\sigma `$($`1.7\sigma `$ in CLEO II) between $`BK^+\pi ^{}`$ and $`B\pi ^+\pi ^{}`$.
We have studied backgrounds from $`bc`$ decays and other $`bu`$ and $`bs`$ decays and find that all are negligible for the analyses presented here. The main background arises from $`e^+e^{}q\overline{q}`$ (where $`q=u,d,s,c`$). Such events typically exhibit a two-jet structure and can produce high momentum back-to-back tracks in the fiducial region. To reduce contamination from these events, we calculate the angle $`\theta _S`$ between the sphericity axis of the candidate tracks and showers and the sphericity axis of the rest of the event. The distribution of $`\mathrm{cos}\theta _S`$ is strongly peaked at $`\pm 1`$ for $`q\overline{q}`$ events and is nearly flat for $`B\overline{B}`$ events. We require $`|\mathrm{cos}\theta _S|<0.8`$ which eliminates $`83\%`$ of the background. Using a detailed GEANT-based Monte Carlo simulation we determine overall detection efficiencies $``$ of $`1148\%`$, as listed in Table I. Efficiencies include the branching fractions for $`K^0K_S^0\pi ^+\pi ^{}`$ and $`\pi ^0\gamma \gamma `$ where applicable.
Additional discrimination between isotropic signal and rather jetty $`q\overline{q}`$ background is provided by the cosine of the angle between the candidate sphericity axis and beam axis (expected to be isotropic for signal, $`1+cos^2\theta `$ distribution for $`q\overline{q}`$ background); the ratio of Fox-Wolfram moments $`H_2/H_0`$ (expected to be smaller for signal than for background); and the distribution of the energy from the rest of the event relative to the candidate’s sphericity axis, as characterized by the energy in nine $`10^{}`$ angular bins. These 11 variables are combined by a Fisher discriminant technique as described in detail in Ref. The Fisher discriminant is a linear combination of experimental observables $`_{i=1}^N\alpha _iy_i`$, where the coefficients $`\alpha _i`$ are chosen to maximize the separation between the simulated signal and background samples.
We perform unbinned maximum-likelihood fits using $`\mathrm{\Delta }E`$, $`M`$, $``$, the angle between the $`B`$ meson momentum and beam axis, and $`dE/dx`$ (where applicable) as input information for each candidate event to determine the signal yields. Four different fits are performed, one for each topology ($`h^+h^{}`$, $`h^\pm \pi ^0`$, $`h^\pm K_S^0`$, and $`K_S^0\pi ^0`$, $`h^\pm `$ referring to a charged kaon or pion). In each of these fits, the likelihood of the event is parametrized by the sum of probabilities for all relevant signal and background hypotheses, with relative weights determined by maximizing the likelihood function $``$. The probability of a particular hypothesis is calculated as a product of the probability density functions (PDFs) for each of the input variables. Further details about the likelihood fit can be found in Ref. . The parameters for the PDFs are determined from independent data and high-statistics Monte Carlo samples. We estimate a systematic error on the fitted yield by varying the PDFs used in the fit within their uncertainties. These uncertainties are dominated by the limited statistics in the independent data samples we used to determine the PDFs. The systematic errors on the measured branching fractions are obtained by adding this fit systematic in quadrature with the systematic error on the efficiency.
Figure 1a shows the results of the likelihood fit for $`B\pi ^+\pi ^{}`$ and $`BK^\pm \pi ^{}`$. The curves represent the $`n\sigma `$ contours, which correspond to the increase in $`2\mathrm{ln}`$ by $`n^2`$. Systematic uncertainties are not included in any contour plots. The statistical significance of a given signal yield is determined by repeating the fit with the signal yield fixed to be zero and recording the change in $`2\mathrm{ln}`$. We also compute from the PDFs the event-by-event probability to be signal or continuum background, as well as the probability to be $`K\pi `$-like or $`\pi \pi `$-like. From these we form likelihood ratios, $`_{sig}=(P_{\pi \pi }^s+P_{K\pi }^s)/(P_{\pi \pi }^s+P_{K\pi }^s+P_{\pi \pi }^c+P_{K\pi }^c+P_{KK}^c)`$ and $`_\pi =P_{\pi \pi }^s/(P_{\pi \pi }^s+P_{K\pi }^s)`$. Superscripts $`s`$ and $`c`$ denote signal and continuum background respectively. Figure 1b illustrates the distribution of events in $`_{sig}`$ (vertical axis) and $`_\pi `$ (horizontal axis). The cluster of events in the upper right corner is clear evidence for $`B\pi ^+\pi ^{}`$.
Figures 1(c-f) show distributions in $`M`$ and $`\mathrm{\Delta }E`$ for events after cuts on likelihood ratios $`_{sig}`$ and $`_\pi `$ computed without $`M`$ and $`\mathrm{\Delta }E`$, respectively. The likelihood fit projections for signal and background components, suitably scaled to account for the efficiencies of the additional cuts (50-70 % for signal), are overlaid. Figure 2 shows the likelihood functions for the fits to $`BK^0\pi ^0`$ and $`Bh^\pm \pi ^0`$. Figure 3 shows $`M`$ and $`\mathrm{\Delta }E`$ distributions for $`BK^0\pi ^\pm `$, $`BK^\pm \pi ^0`$, and $`BK^0\pi ^0`$.
We summarize all branching fractions and upper limits in Table I. In addition to the first observations $`B\pi ^+\pi ^{}`$, $`BK^+\pi ^0`$, and $`BK^0\pi ^0`$, we report improved measurements for the decays $`BK^\pm \pi ^{}`$ and $`BK^0\pi ^\pm `$. The table also includes a range of theoretical predictions taken from recent literature. We see some indication for the decay $`B\pi ^\pm \pi ^0`$ with the branching fraction of $`(B\pi ^\pm \pi ^0)=(5.6_{2.3}^{+2.6}\pm 1.7)\times 10^6`$, but statistical significance of the signal yield is insufficient to claim an observation for this decay mode. We find no evidence for the decays $`BK^+K^{}`$ and $`BK^\pm K^0`$, and calculate $`90\%`$ confidence level (CL) upper limit yields by integrating the likelihood function
$$\frac{_0^{N^{UL}}_{\mathrm{max}}(N)𝑑N}{_0^{\mathrm{}}_{\mathrm{max}}(N)𝑑N}=0.90$$
(1)
where $`_{\mathrm{max}}(N)`$ is the maximum $``$ at fixed $`N`$ to conservatively account for possible correlations among the free parameters in the fit. We then increase upper limit yields by their systematic errors and reduce detection efficiencies by their systematic errors to calculate branching fraction upper limits given in Table I.
To evaluate how systematic uncertainties in the PDFs affect the statistical significance for modes where we report first observations, we repeated the fits for the $`h^+h^{},h^+\pi ^0`$ and $`K^0\pi ^0`$ modes with all PDFs changed simultaneously within their uncertainties to maximally reduce the signal yield in the modes of interest. Under these extreme conditions, the significance of the first-observation modes $`\pi ^+\pi ^{},K^+\pi ^0`$ and $`K^0\pi ^0`$ becomes $`3.2,5.3`$ and $`3.8\sigma `$ respectively. We also evaluate the branching fractions with alternative analyses using tighter and looser cuts on the continuum-suppressing variable $`|\mathrm{cos}\theta _S|`$. These variations correspond to halving and doubling the background in the fitted sample. The changes in branching fractions under these variations are insignificant compared to the statistical error of our results.
The ratio of the branching fractions $`(BK^\pm K^0)/(B\pi ^\pm K^0)`$ can be used to estimate the size of final state interactions in charmless rare $`B`$ decays. Following the method outlined above we calculate $`(BK^\pm K^0)/(B\pi ^\pm K^0)<0.3`$ at 90% CL. It has also been suggested to use the ratio of the branching fractions $`(BK^\pm \pi ^{})/(B\pi ^+\pi ^{})`$ to estimate uncertainties in the measurement of the unitarity triangle parameter $`\alpha =\pi \beta \gamma `$ via $`B^0(t)\pi ^+\pi ^{}`$. We obtain $`(BK^\pm \pi ^{})/(B\pi ^+\pi ^{})<15`$ at 90% CL which implies that an error on $`\alpha `$ obtained from time-dependent asymmetry measurements of $`B^0\pi ^+\pi ^{}`$ can be as high as $`60^{}`$ .
In summary: we have made a first observation of $`B\pi ^+\pi ^{}`$; measured branching fractions for all four exclusive $`BK\pi `$, including first observations of the decays $`BK^+\pi ^0`$ and $`BK^0\pi ^0`$; obtained improved upper limits on $`B\pi ^+\pi ^0`$ and $`BK\overline{K}`$ modes. The hierarchy of branching fractions $`KK<\pi \pi <K\pi `$ is obvious.
We thank W.-S. Hou for many useful discussions. We gratefully acknowledge the effort of the CESR staff in providing us with excellent luminosity and running conditions. This work was supported by the National Science Foundation, the U.S. Department of Energy, the Research Corporation, the Natural Sciences and Engineering Research Council of Canada, the A.P. Sloan Foundation, the Swiss National Science Foundation, and Alexander von Humboldt Stiftung. |
warning/0001/hep-th0001005.html | ar5iv | text | # 1 Vertex operator algebras and chiral blocks
## 1 Vertex operator algebras and chiral blocks
Chiral blocks – also known as conformal blocks – arise in the study of two-dimensional conformal field theory. In physics terminology, they are correlation “functions” of so-called chiral vertex operators $`\varphi _\mu `$. Thus, roughly, one deals with objects of the form
$$\varphi _{\mu _1}(p_1)\varphi _{\mu _2}(p_2)\mathrm{}\varphi _{\mu _m}(p_m)_{(\stackrel{~}{C})},$$
(1)
where $`\stackrel{~}{C}`$ is a two-dimensional manifold, some kind of ‘operator product’ between the chiral vertex operators $`\varphi _{\mu _i}`$ ‘sitting’ at $`p_i\stackrel{~}{C}`$ is understood, and $`\mathrm{}`$ stands for the operation of forming the ‘vacuum expectation value’. Yet, chiral blocks are in general neither functions, nor uniquely determined by these data. For a more detailed understanding several concepts are needed, among them in particular the following.
First, the notion of a vertex operator algebra $`𝔄=(_\mathrm{\Omega },\text{Y},v_\mathrm{\Omega },v_{\mathrm{Vir}})`$. Here $`_\mathrm{\Omega }=_n_\mathrm{\Omega }^{(n)}`$ is an infinite-dimensional $``$-graded vector space, with finite-dimensional homogeneous subspaces $`_\mathrm{\Omega }^{(n)}`$. $`_\mathrm{\Omega }`$ is endowed with infinitely many products, which are encoded in the vertex operator map $`\text{Y}:_\mathrm{\Omega }\mathrm{End}(_\mathrm{\Omega })_{}[[t,t^1]]`$, mapping $`_\mathrm{\Omega }`$ to the Laurent series in a formal variable $`t`$ with values in the endomorphisms of $`_\mathrm{\Omega }`$. The vacuum element $`v_\mathrm{\Omega }_\mathrm{\Omega }^{(0)}`$ and the Virasoro element $`v_{\mathrm{Vir}}_\mathrm{\Omega }^{(2)}`$ are distinguished vectors in $`_\mathrm{\Omega }`$, satisfying $`\text{Y}(v_\mathrm{\Omega })=\text{id}`$ and $`(\text{Y}(v;t=0))v_\mathrm{\Omega }=v`$ (Y is therefore also known as state-field correspondence).
These quantities are subject to a number of further axioms (see e.g. ), mostly not to be spelled out here. We only mention the requirement that the endomorphisms $`L_n`$ defined by the expansion $`\text{Y}(v_{\mathrm{Vir}})=_nL_nt^{n2}`$ form a basis of the Virasoro algebra $`𝒱`$ir. (More precisely, they provide a representation of $`𝒱`$ir in which the central element acts as a constant multiple, called the rank of $`𝔄`$, of the identity.)
For many purposes, it is sufficient to regard the vertex operator algebra as the Lie algebra spanned over $``$ by the FourierLaurent modes of suitable vertex operators $`\text{Y}(v)`$. Besides $`𝒱`$ir (i.e. the $`L_n`$) and its supersymmetric generalizations, examples are so-called $`𝒲`$-algebras and untwisted affine Lie algebras $`𝔤`$. The term chiral algebra is used both for the proper vertex operator algebra $`𝔄`$ and for the Lie algebra $``$($`𝔄`$).
There is a collection of irreducible $`𝔄`$-modules $`_\mu `$, with $`\mu `$ in some index set $`I`$. This includes $`\mathrm{\Omega }I`$, i.e. the vector space $`_\mathrm{\Omega }`$ underlying $`𝔄`$; this is called the vacuum sector.
The modules $`_\mu `$ are graded weight modules, with finite-dimensional weight spaces, where the weights are with respect to the zero mode $`L_0`$ of $`𝒱`$ir and a suitable collection $`\{H_0^i\}`$ of other mutually commuting modes in $``$($`𝔄`$). Hence there is the notion of characters, i.e. generating functions $`\chi _\mu (\tau ,\stackrel{}{z})=\mathrm{tr}__\mu \mathrm{e}^{2\pi \mathrm{i}\tau L_0}\mathrm{e}^{2\pi \mathrm{i}\stackrel{}{z}\stackrel{}{H}_0}`$ for weight multiplicities.
When every $`𝔄`$-module is fully reducible and $`|I|<\mathrm{}`$, one speaks of a rational vertex operator algebra, respectively rational CFT. Below we restrict to this case. In a rational theory the modules $`_\mu `$ constitute the simple objects of a modular tensor category (provided that $`𝔄`$ is ‘maximal’, which corresponds to non-degeneracy of braiding).
We can now describe chiral blocks more properly; they are certain linear forms
$$B_\stackrel{}{\mu }:_{\mu _1}_{\mu _2}\mathrm{}_{\mu _m}\stackrel{}{}_\stackrel{}{\mu }$$
(2)
on the tensor product $`\stackrel{}{}_\stackrel{}{\mu }`$ of the relevant irreducible modules $`_{\mu _i}`$. Further, to establish the connection with formula (1) one considers a complex curve $`\stackrel{~}{C}`$ with ordered marked points $`p_i\stackrel{~}{C}`$, and identifies for each $`i`$ the formal variable $`t`$ with a local holomorphic coordinate $`\zeta _i`$ at $`p_i\stackrel{~}{C}`$. (For $`\stackrel{~}{C}=^1`$ one may take $`t\widehat{=}\zeta _i=zz_i`$ with $`z`$ a quasi-global holomorphic coordinate on $`^1`$ – say the standard global coordinate on $``$ for $`^1\{\mathrm{}\}`$ – such that $`z(p_i)=z_i`$. For higher genus curves, the prescription becomes more complicated.) Then for each $`v_\mathrm{\Omega }`$ the vertex operator map provides a chiral vertex operator as appearing in (1), according to $`\text{Y}(v)\widehat{=}\varphi _\mathrm{\Omega }(v;p_i)`$. For other sectors $`\mu \mathrm{\Omega }`$, the chiral vertex operators $`\varphi _\mu (v;p_i)`$ with $`v_\mu `$ correspond in an analogous manner to intertwining operators between $`𝔄`$-modules. Then one sets
$$\varphi _{\mu _1}(v_1;p_1)\varphi _{\mu _2}(v_2;p_2)\mathrm{}\varphi _{\mu _m}(v_m;p_m)_{(\stackrel{~}{C})}=B_\stackrel{}{\mu }(v_1v_2\mathrm{}v_m).$$
(3)
Hereby the chiral algebra is interpreted as the “local implementation of the symmetries” of the system at the insertion points $`\stackrel{}{p}=(p_1,p_2,\mathrm{},p_m)`$. But we also want to study (3) in its dependence on the insertion points $`\stackrel{}{p}`$ and on the moduli $`\stackrel{}{\tau }`$ of $`\stackrel{~}{C}`$. This necessitates the construction, for each curve $`\stackrel{~}{C}`$ and number $`m`$ of insertions, of a suitable “global implementation of the symmetries”. Such an implementation, to be called a block algebra<sup>1</sup><sup>1</sup>1 Unfortunately, it is $`𝔄_{\stackrel{}{p},\stackrel{~}{C}}`$ that mathematicians sometimes call the ‘chiral algebra’, see e.g. . and denoted by $`𝔄_{\stackrel{}{p},\stackrel{~}{C}}`$, is a family (varying with the insertion points and moduli) of subalgebras of the $`m`$-fold tensor product of $``$($`𝔄`$), and can be thought of as providing a generalized co-product.
The action of $`𝔄_{\stackrel{}{p},\stackrel{~}{C}}`$ on $`\stackrel{}{}_\stackrel{}{\mu }`$ allows us to be specific about the linear forms (2). Namely, one defines the chiral blocks to be the space $`B_\stackrel{}{\mu }=\text{(}(\stackrel{}{}_\stackrel{}{\mu })_{}^{}\text{)}_{}^{𝔄_{\stackrel{}{p},\stackrel{~}{C}}}`$ of $`𝔄_{\stackrel{}{p},\stackrel{~}{C}}`$-singlets in the algebraic dual $`(\stackrel{}{}_\stackrel{}{\mu })^{}`$ – or dually, as the space
$$\overline{)B_\stackrel{}{\mu }^{}=\stackrel{}{}_\stackrel{}{\mu }_{𝔄_{\stackrel{}{p},\stackrel{~}{C}}}^{}}$$
(4)
of co-invariants of $`\stackrel{}{}_\stackrel{}{\mu }`$ with respect to $`𝔄_{\stackrel{}{p},\stackrel{~}{C}}`$. In physics terminology, this prescription says that $`B_\stackrel{}{\mu }`$ is the space of solutions to the Ward identities of the system.
## 2 Chiral blocks in WZW models
By a WZW model one means a conformal field theory for which $``$($`𝔄`$) is an untwisted affine Lie algebra $`𝔤`$ (or, more precisely, its semi-direct sum with $`𝒱`$ir) and for which the Virasoro representation is supplied by the affine Sugawara construction, which says that the Virasoro generators are quadratic expressions in the generators of $`𝔤`$, with coefficients proportional to the Killing form of the horizontal subalgebra $`𝔤_0𝔤`$.
Various notions of conformal field theory have very concrete WZW realizations:
One can regard the affine Lie algebra $`𝔤`$ as being obtained via the loop construction from a finite-dimensional simple Lie algebra $`\overline{𝔤}`$; the formal variable $`t`$ of the vertex operator formalism is closely related to the indeterminate of the loop construction. In this presentation $`𝔤`$ has a basis $`\{J_n^a\}`$ with $`n`$ and $`\{J^a\}`$ a basis of $`\overline{𝔤}`$ (together with a central element $`K`$ and a derivation $`D`$), and the simple Lie algebra $`\overline{𝔤}`$ can be identified with the horizontal subalgebra $`𝔤_0`$, which is spanned by the zero modes $`J_0^a`$.
The spaces $`_\mu `$ are irreducible highest weight modules over $`𝔤`$ with integrable highest weight $`\mu `$ of fixed level $`k_{>0}`$. There are only finitely many such weights $`\mu `$ (in fact the theory is rational), namely those whose horizontal part $`\overline{\mu }`$ is a dominant integral $`\overline{𝔤}`$-weight with inner product $`(\overline{\mu },\overline{\theta }^{})k`$ with the highest coroot $`\overline{\theta }^{}`$. For instance, for $`\overline{𝔤}=𝔰𝔩(2)`$ only the $`𝔰𝔩(2)`$-weights $`\overline{\mu }=\mathrm{\hspace{0.17em}0},1,\mathrm{},k`$ are allowed.
The vacuum sector is the basic $`𝔤`$-module, which has highest weight $`\mathrm{\Omega }k\mathrm{\Lambda }_{(0)}`$.
One has $`L_0=D`$, and $`\stackrel{}{H}_0`$ form a basis of the Cartan subalgebra of $`\overline{𝔤}`$. The characters $`\chi _\mu `$ are obtained from the corresponding formal $`𝔤`$-characters – i.e. elements of the group algebra spanned by formal exponentials in the weights – by interpreting the formal exponentials as functions on weight space. They are convergent for $`\mathrm{}(\tau )>\mathrm{\hspace{0.17em}0}`$.
The block algebra $`𝔄_{\stackrel{}{p},\stackrel{~}{C}}`$ is then (for details see e.g. ) the tensor product
$$𝔤_{\stackrel{}{p},\stackrel{~}{C}}=\overline{𝔤}_{\stackrel{}{p},\stackrel{~}{C}},$$
(5)
with $`_{\stackrel{}{p},\stackrel{~}{C}}`$ the algebra of functions holomorphic on $`\stackrel{~}{C}\{\stackrel{}{p}\}`$ and with (at most) finite order poles at the $`p_i`$. The action of $`𝔤_{\stackrel{}{p},\stackrel{~}{C}}`$ on $`\stackrel{}{}_\stackrel{}{\mu }`$ is given by $`\text{(}R_\stackrel{}{\mu }(\overline{x}f)\text{)}(v_1v_2\mathrm{}v_m):=_{i=1}^mv_1v_2\mathrm{}`$ $`R_{\mu _i}(x_i)v_i\mathrm{}v_m`$ for $`v_1v_2\mathrm{}v_m\stackrel{}{}_\stackrel{}{\mu }`$, where the $`x_i\overline{x}f_{p_i}`$ are to be regarded as elements of $`𝔤`$ ($`f_{p_i}`$ denotes the local expansion of $`f`$ at $`p_i`$).
For general CFTs, much less is known about block algebras and their action on tensor products. Roughly, one must ‘couple’ Virasoro-(quasi)primary fields to meromorphic sections of suitable powers of the canonical bundle of $`\stackrel{~}{C}`$; in the WZW case this power is zero, hence one deals with functions (5) and can be very explicit. Thus for general CFTs many facets of what is reported below are not at all rigorous – to us it is a major challenge in CFT to improve this. <sup>2</sup><sup>2</sup>2 It is as yet unclear whether the vertex operator framework is broad enough for a rigorous discussion of all issues of interest in CFT, or whether one must resort to formulations involving e.g. von Neumann algebras. (The latter would be unfortunate, as one would give up on treating non-unitary models, like ghost systems in string theory, at an equal footing as unitary ones.) In the WZW case, such a formulation follows by studying the loop group $`\mathrm{LG}`$ of the compact, connected and simply connected real Lie group $`\mathrm{G}`$ whose Lie algebra is the compact real form of $`\overline{𝔤}`$, as well as the associated local loop groups and their representations on Hilbert spaces; compare e.g. .
In this context, note that in the vertex operator setting no topology is chosen on the vector spaces $`_\lambda `$, i.e. even in the unitary case they are only pre-Hilbert spaces. In fact, for certain purposes – e.g. when trying to achieve that the generators of $``$($`𝔄`$) act continuously on the (dual) blocks – other topologies than the Hilbert space topology based on the standard norm can be more convenient. In contrast, for WZW models already enough is known so as to make precise statements and establish rigorous proofs.
## 3 Bundles of chiral blocks
Chiral block spaces have been studied in quite some detail for several reasons. (In CFT their significance emanates from the fact that they contain the physical correlation functions as special elements, see below.) Among the pertinent results are:
In all known cases (not only for rational CFTs), for fixed insertion points $`\stackrel{}{p}`$ and fixed moduli $`\stackrel{}{\tau }`$ of $`\stackrel{~}{C}`$, the space $`B_\stackrel{}{\mu }`$ is a finite-dimensional vector space $`B_\stackrel{}{\mu }(\stackrel{}{p},\stackrel{}{\tau })`$. This has a counterpart in the associated tensor category: all morphism spaces in a $`C^{}`$-tensor category with conjugates and irreducible unit are finite-dimensional .
The spaces $`B_\stackrel{}{\mu }(\stackrel{}{p},\stackrel{}{\tau })`$ fit together to the total space of a finite rank vector bundle $`_\stackrel{}{\mu }`$ over the moduli space of genus $`g`$ complex curves with $`m`$ ordered marked points.
Using $`𝒱`$ir one constructs a projectively flat ‘KnizhnikZamolodchikov’ connection on $`_\stackrel{}{\mu }`$. (Some authors reserve the term ‘chiral block’ for flat sections of $`_\stackrel{}{\mu }`$.)
The block bundles are in general not irreducible (i.e. the fibers decompose into a direct sum of vector spaces, in a manner compatible with the transition functions).
One of the major reasons for independent mathematical interest in chiral blocks is the role played in algebraic geometry by the WZW one-point blocks with vacuum insertion $`\varphi _\mathrm{\Omega }`$. Namely, the Picard group of the moduli space $`_{\mathrm{G},\stackrel{~}{C}}`$ of holomorphic principal $`\mathrm{G}_\mathrm{c}`$-bundles (with $`\mathrm{G}_\mathrm{c}`$ the complexification of G) over $`\stackrel{~}{C}`$ modulo stable equivalence, <sup>3</sup><sup>3</sup>3 $`_{\mathrm{G},\stackrel{~}{C}}`$ possesses several other interpretations as well, such as: the set of equivalence classes of flat principal G-bundles; the space of semi-stable holomorphic vector bundles $`\mathrm{E}`$ over $`\stackrel{~}{C}`$ such that the sheaf of sections of the determinant bundle is the structure sheaf of $`\stackrel{~}{C}`$, $`\mathrm{det}\mathrm{E}=𝒪_{\stackrel{~}{C}}`$; and the phase space $`𝒜_{}/𝒢`$ (flat connections modulo gauge transformations) of ChernSimons gauge theory.
In the first place, these are just bijections of sets. But each of the sets comes equipped with its own natural structures. One can translate those, so that indeed one gets a set with various different interesting structures. A crucial input for establishing these relations is BorelWeilBott theory. is generated by the determinant line bundle $`\mathrm{L}`$. $`\mathrm{L}=𝒪(\theta )`$ is a locally free rank-one sheaf of meromorphic functions on $`_{\mathrm{G},\stackrel{~}{C}}`$, where $`\theta `$ is the Theta divisor. Now for every $`k_{>0}`$, the space of holomorphic sections of $`\mathrm{L}^k`$ is canonically isomorphic to the space of one-point blocks on $`\stackrel{~}{C}`$ with insertion $`\varphi _\mathrm{\Omega }`$ at level $`k`$:
$$H^0(_{\mathrm{G},\stackrel{~}{C}},\mathrm{L}^k)B_{k\mathrm{\Lambda }_{(0)}}(\stackrel{~}{C}).$$
(6)
With traditional methods this ‘space of generalized Theta functions’ had been accessible only in a few special cases (for more information, see e.g. ).
## 4 Dimensions
When studying chiral blocks, the first quantity of interest that comes to mind is the rank of the bundle $`_\stackrel{}{\mu }`$, i.e. the dimension $`\mathrm{N}_{\stackrel{}{\mu };\stackrel{~}{C}}=\mathrm{dim}B_{\stackrel{}{\mu };\stackrel{~}{C}}`$ of the spaces $`B_{\stackrel{}{\mu };\stackrel{~}{C}}`$. In CFT, the integers $`\mathrm{N}_{\lambda ,\mu }^\nu \mathrm{N}_{\lambda ,\mu ,\nu _{}^+;^1}`$ give the fusion rules, i.e. the number $`\mathrm{\#}(\varphi _\lambda \varphi _\mu \varphi _\nu )`$ of ‘independent couplings’ between families of fields. Factorization (see section 7 below) implies that the fusion rules constitute the structure constants of a commutative semi-simple associative algebra with unit and involution, which is called the fusion rule algebra. They can be expressed in terms of a unitary symmetric matrix $`S`$ by the Verlinde formula
$$\mathrm{N}_{\lambda ,\mu }^\nu =\underset{\kappa I}{}S_{\kappa ,\lambda }S_{\kappa ,\mu }S_{\kappa ,\nu }^{}/S_{\kappa ,\mathrm{\Omega }}.$$
(7)
It is worth pointing out that the existence of a diagonalizing matrix $`S`$ obeying (7) is an immediate by-product of the representation theory of fusion rule algebras. The contents of the Verlinde conjecture is not formula (7) in itself, but rather that it is one and the same matrix $`S`$ that appears in (7) and that affords the modular transformation $`\tau 1/\tau `$ on the characters $`\chi _\mu `$. This implies in particular concrete expressions for $`S`$, e.g. the KacPeterson formula for WZW models, and similarly for coset models and WZW orbifolds.
By factorization (see below), the Verlinde formula (7) generalizes as
$$\overline{)\mathrm{N}_{\stackrel{}{\mu };\stackrel{~}{C}}=\underset{\kappa I}{}|S_{\kappa ,\mathrm{\Omega }}|^{22g}\underset{i=1}{\overset{m}{}}\frac{S_{\kappa ,\mu _i}}{S_{\kappa ,\mathrm{\Omega }}}}$$
(8)
to an arbitrary number $`m`$ of insertions and arbitrary genus $`g`$. (8) has been proven rigorously only for WZW models (in particular by algebraic geometry means, cf. e.g. and also ). But there is enormous evidence that it holds in general; in particular it was verified for very many theories that the numbers (8) are in $`_0`$.
## 5 Traces
The dimensions (8) are only the most basic characteristics of blocks. Other quantities are, of course, of interest as well. As the blocks are in general not irreducible as vector bundles, a natural generalization are the dimensions of irreducible sub-bundles. For many chiral blocks, a non-trivial sub-bundle structure follows from the presence of some group $`𝒮`$ of automorphisms $`\sigma `$ of $`𝔄_{\stackrel{}{p},\stackrel{~}{C}}`$, which in turn come from automorphisms of $`𝔄`$. There are then linear bijections $`\mathrm{\Theta }_\sigma `$ between the $`_\mu `$ satisfying the twisted intertwiner property $`\mathrm{\Theta }_\sigma ^1\text{Y}(\sigma v;z)\mathrm{\Theta }_\sigma =\text{Y}(v;z)`$ and descending to linear maps $`\mathrm{\Theta }_\stackrel{}{\sigma }`$ on the blocks. The $`\mathrm{\Theta }_\stackrel{}{\sigma }`$ realize $`𝒮`$ projectively, and the sub-bundles are obtained by the simultaneous eigenspace decomposition of the blocks with respect to these maps .
Some information on the dimensions of such sub-bundles is available, too. The dimensions are most favorably expressed in terms of traces of the twisted intertwiners $`\mathrm{\Theta }_\stackrel{}{\sigma }`$, to which they are related by Fourier transformation with respect to the subgroup of $`𝒮`$ that corresponds to the center of a twisted group algebra, where the twist is by the cocycle defining the projectivity of the action on the blocks. Concretely, there are generalizations of the Verlinde conjecture for two important types of automorphisms:
First, for automorphisms associated to simple currents $`\varphi _\mathrm{J}`$. A simple current is a unit of the fusion algebra; it can be characterized by the equality $`S_{\mathrm{J},\mathrm{\Omega }}=S_{\mathrm{\Omega },\mathrm{\Omega }}`$. Simple currents of WZW models correspond to symmetries of the Dynkin diagram of the underlying affine Lie algebra $`𝔤`$ and thereby to certain outer automorphisms $`\sigma _\mathrm{J}`$ of $`𝔤`$. The proposed formula for the traces of the corresponding twisted intertwiners $`\theta _\stackrel{}{\mathrm{J}}\theta _{\stackrel{}{\sigma }_\mathrm{J}}`$ reads
$$\overline{)\mathrm{Tr}_{_{\stackrel{}{\mu };^1}}\text{(}\mathrm{\Theta }_{\mathrm{J}_1,\mathrm{J}_2,\mathrm{},\mathrm{J}_m}\text{)}=\underset{\kappa :\mathrm{J}_{\mathrm{}}\kappa =\kappa }{}|S_{\kappa ,\mathrm{\Omega }}|_{}^2\underset{i=1}{\overset{m}{}}\frac{S_{\kappa ,\mu _i}^{\mathrm{J}_i}}{S_{\kappa ,\mathrm{\Omega }}}.}$$
(9)
Here it is assumed that $`\mathrm{J}_1\mathrm{J}_2\mathrm{}\mathrm{J}_m=\mathrm{\Omega }`$ as well as $`\mathrm{J}_i\mu _i=\mu _i`$ for all $`i=\mathrm{\hspace{0.17em}1},2,\mathrm{},m`$ (when formally extended to other cases, the expression (9) just yields zero). Further, $`S^\mathrm{J}`$ is the modular matrix for one-point blocks with insertion $`\varphi _\mathrm{J}`$ on an elliptic curve.
Second, for automorphisms $`\sigma `$ of $`𝔄`$ that preserve $`v_{\mathrm{Vir}}`$ one finds
$$\overline{)\mathrm{Tr}_{_{\stackrel{}{\mu };^1}}\text{(}\mathrm{\Theta }_{\sigma ,\sigma ,\mathrm{},\sigma }\text{)}=\underset{\dot{\kappa }}{}|S_{\dot{\kappa },\mathrm{\Omega }}^{}|_{}^2\underset{i=1}{\overset{m}{}}\frac{S_{\dot{\kappa },\mu _i}^{}}{S_{\dot{\kappa },\mathrm{\Omega }}^{}}.}$$
(10)
Here $`S^{}`$ is an ingredient of the modular S-matrix for an orbifold theory that is formed from the original CFT by quotienting out $`\sigma `$. Note that $`S^{}`$ has two distinct types of labels; they correspond to the $`\sigma `$-twisted versus the untwisted sector of the orbifold.
In the WZW case, (9) is closely related to a Verlinde formula for non-simply connected groups, see eq. (19) below, while $`S^{}`$ coincides with the ordinary S-matrix for (a pair of) twisted affine Lie algebras (being ‘genuinely twisted’ iff $`\sigma `$ is outer).
But both (9) and (10) are conjectured for arbitrary rational CFTs – they originate from structures present in every rational CFT. Even for WZW they are far from being proven rigorously – a proof would in particular imply a proof of the Verlinde formula itself. But there are definite ideas for WZW models and also for some derived theories like coset models. In addition, there is enormous numerical evidence: one obtains non-negative integers for the ranks, <sup>4</sup><sup>4</sup>4 A surprising empirical observation is that in the case of (9) the traces are actually integral themselves, even when the order of $`\sigma _\mathrm{J}`$ is larger than 2. This is reminiscent of a description of $`\mathrm{dim}B_{\stackrel{}{\mu };^1}`$ as the Euler number of a suitable BGG-like complex, and hence suggests that the traces may possess a homological interpretation as well. (On the other hand, the acyclicity result that implies non-negativity of the dimensions cannot generalize. Similar structures have appeared in .) even though they are obtained as complicated sums of arbitrary (well, they all lie in a cyclotomic extension of $``$) complex numbers.
## 6 Chiral versus full CFT
The sub-bundle structure described by the result (10) can be understood in the framework of orbifold CFTs. As it turns out, the very same chiral concepts play a role in the study of symmetry breaking boundary conditions. Therefore in the rest of this note we address this conceptually different (and at first sight totally unrelated) topic.
As a first step, let us point out that the whole discussion so far concerns what we like to call chiral conformal field theory, that is, <sup>5</sup><sup>5</sup>5 Sometimes the term ‘chiral CFT’ is used in a slightly different fashion. CFT on a compact two-dimensional manifold without boundary that has a complex structure, or in short, CFT on a complex curve $`\stackrel{~}{C}`$. The analytic properties of $`\stackrel{~}{C}`$ enter in particular in the definition of block algebras. In contrast, in most applications in physics, <sup>6</sup><sup>6</sup>6 Among them are string theory and many condensed matter phenomena. But there do exist applications where it is chiral CFT proper that is relevant. An important example is three-dimensional topological field theory – ChernSimons theory in the WZW case – and thereby the (fractional) quantum Hall effect. (While a priori in the quantum Hall effect there is thus no natural place for modular invariance on the torus, arguments assigning a physical role to it were given in the literature.) one must consider CFT on a real two-dimensional manifold $`C`$<sup>7</sup><sup>7</sup>7 You might have already wondered why the symbol $`\stackrel{~}{C}`$ was used above rather than just $`C`$. with conformal structure. $`C`$ may be non-orientable or have a boundary, and it does not come with a natural orientation even when it is orientable. We will refer to CFT on $`C`$ as full conformal field theory.
While in chiral CFT one deals with the chiral algebra $`𝔄`$, chiral vertex operators, chiral blocks, characters, and fusion rules, the key notions in full CFT are fields, correlation functions, the torus partition function, operator products, and boundary conditions.
Chiral CFT is of interest in its own right. But it also serves as a convenient intermediate step in the analysis of full CFT, since it allows to exploit the power of complex geometry. At the geometrical level, the relation between chiral and full conformal field theory is pretty simple. The surface $`C`$ possesses an oriented two-sheeted Schottky cover $`\stackrel{~}{C}`$, branched over the boundary $`C`$, from which one recovers $`C`$ by dividing out a suitable anticonformal involution $``$. Here are the simplest examples:
Orientable, no boundary: $`C=S^2`$ (sphere) $``$ $`\stackrel{~}{C}=^1^1`$, $`:(z,\dot{z})(\dot{z}^{},z^{})`$.
Orientable, with boundary: $`C=D^2`$ (disk) $``$ $`\stackrel{~}{C}=^1`$, $`:z\mathrm{\hspace{0.17em}1}/z^{}`$.
Non-orientable: $`C=^2`$ (projective plane / ‘crosscap’) $``$ $`\stackrel{~}{C}=^1`$, $`:z1/z^{}`$.
To implement the transition from $`\stackrel{~}{C}`$ to $`C`$ at the field theory level requires more work; e.g. for connected $`\stackrel{~}{C}`$ the block algebras are understood only in the simplest cases. For now suffice it to say that, in a rough sense, in many respects the transition amounts to taking two copies of chiral objects. In particular, each single (bulk) field on $`C`$ comes from two chiral vertex operators on $`\stackrel{~}{C}`$ (physically speaking, one has ‘image charges’). Thus it carries two chiral labels; we denote it by $`\varphi _{\mu ,\dot{\mu }}`$. In addition one must impose some additional constraints and identifications, to be given below.
Before proceeding, let us recall that for WZW models various structures can be made fully explicit which for other classes of CFTs are not yet worked out in detail. Fortunately, this is an issue mainly for chiral CFT. Once the chiral theory is taken for granted, considerations in full CFT turn out to be essentially model independent.
## 7 Correlation functions
Correlation functions are the ‘vacuum expectation values’ of suitable products of ‘(quantum) fields’. They constitute the quantities of most direct interest in applications. For instance, by integrating them over moduli space one obtains string scattering amplitudes. A message to be remembered is that fields and their correlation functions are objects in full CFT and hence ‘live’ on the quotient $`C`$ of $`\stackrel{~}{C}`$, while chiral blocks ‘live’ on $`\stackrel{~}{C}`$. Thus the blocks cannot be physical correlation functions; rather, a correlation function for $`C`$ is a specific element in a corresponding space of blocks on $`\stackrel{~}{C}`$. <sup>8</sup><sup>8</sup>8 Or what is the same, after choosing some (natural) basis in the block space: a specific linear combination of basis blocks. When $`\stackrel{~}{C}`$ is disconnected, this is usually written as a sesqui-linear combination of separate basis blocks for the two connected components of $`\stackrel{~}{C}`$. That element is determined by various constraints, coming in three types:
Locality: Correlators are (single-valued) functions of the insertion points $`p_i`$ – unlike generic sections of the block bundle, which typically is not a trivial vector bundle.
Locality: They are also functions of the moduli of $`C`$ (modulo the Weyl anomaly).
Factorization: They are compatible with desingularization. This amounts to a restriction on the allowed intermediate states that contribute in singular limits, and is thereby closely related to the existence of operator product expansions.
Technically: $`\stackrel{~}{C}`$ is a stable algebraic curve with at worst ordinary double points as singularities. When $`\widehat{C}`$ is a partial desingularization of $`\stackrel{~}{C}`$ that resolves a double point $`p\stackrel{~}{C}`$ in two points $`p^{},p^{\prime \prime }\widehat{C}`$, then factorization gives a canonical isomorphism
$$_{\nu I}B_{\stackrel{}{\mu },\nu ,\nu _{}^+;\widehat{C}}B_{\stackrel{}{\mu },\stackrel{~}{C}}.$$
(11)
A priori neither existence nor uniqueness of a solution to these constraints is clear.
A prominent example is provided by the correlator for $`m=\mathrm{\hspace{0.17em}0}`$ and $`g=\mathrm{\hspace{0.17em}1}`$. Then $`C`$ is a torus, and $`\stackrel{~}{C}`$ is the disconnected sum $`E_\tau E_{\dot{\tau }}`$ of two elliptic curves with opposite orientation, $`\dot{\tau }=\tau ^{}`$. The 0-point correlator on $`C`$ is the torus partition function $`Z`$, while a basis for the 0-point blocks on $`\stackrel{~}{C}`$ are tensor products of ($`𝒱`$ir-specialized) irreducible characters $`\chi _\mu `$ and $`\dot{\chi }_{\dot{\mu }}`$. So $`Z`$ is a sesqui-linear combination of characters:
$$Z(\tau )=_{\mu ,\dot{\mu }}Z_{\mu ,\dot{\mu }}\chi _\mu (\tau )\text{(}\dot{\chi }_{\dot{\mu }}(\tau )\text{)}^{}.$$
(12)
$`Z`$ is highly constrained by the property of modular invariance, i.e. locality with respect to the modulus $`\tau `$ of $`E_\tau `$. The solution of this constraint is of interest in its own right.
## 8 Boundary conditions
Solving the factorization and locality constraints is not easy at all, in general, and little is known about uniqueness. <sup>9</sup><sup>9</sup>9 Just think of the case of the torus partition function (12), where the constraint (modular invariance) looks quite innocent, but is still hard to solve. See e.g. for review and references. But for special correlation functions, which are still of great interest, a lot can be done explicitly and in much generality. An especially fortunate example is given by the 1-point functions $`\varphi _{\mu ,\dot{\mu }}`$ for bulk fields on the disk. These are important because, due to factorization, there is only a small number of basic building blocks. As long as only closed orientable surfaces $`C`$ are studied, already the 3-point functions on $`S^2`$ are sufficient. In the general situation, in addition the 1-point functions $`\varphi _{\mu ,\dot{\mu }}`$ on the disk and on $`^2`$ are needed (as well as 3-point functions for boundary fields, which correspond to open string vertex operators) .
The chiral blocks for $`\varphi _{\mu ,\dot{\mu }}`$ are two-point blocks on the Schottky cover $`^1`$ of the disk. Now one has $`\mathrm{N}_{\mu ,\dot{\mu }}=\delta _{\dot{\mu },\mu _{}^+}`$, so only a single coefficient needs to be determined:
$$\varphi _{\mu ,\dot{\mu }}(vv^{})_a=_{\mu ,\dot{\mu };\mathrm{\Omega }^{}}^a_{}B_{\mu ,\dot{\mu }}(vv^{}).$$
(13)
Again by factorization, the complex number $`_{\mu ,\dot{\mu };\mathrm{\Omega }^{}}^a_{}`$ can be interpreted as a reflection coefficient, which appears in the bulk-boundary operator product
$$\varphi _{\mu ,\dot{\mu }}(z)_{\nu I}(1|z|^2)^{2\mathrm{\Delta }_\mu +\mathrm{\Delta }_\nu }_{\mu ,\dot{\mu };\nu ^{}}^a_{}\mathrm{\Psi }_\nu ^{a,a}(\mathrm{arg}z)\text{for }|z|\mathrm{\hspace{0.17em}1}.$$
(14)
For closed orientable $`C`$ it is generally expected that the constraints possess a unique solution. In contrast, the one-point functions on the disk are in general not unique, but an additional label $`a`$ is needed. This indicates that the disk can come with several distinct boundary conditions labelled by $`a`$. A boundary condition is essentially the same as a consistent collection of one-point functions of bulk fields on the disk.
One of the most fundamental tasks in CFT is to determine, assuming the theory to be known at the chiral level, all consistent conformally invariant boundary conditions.
## 9 Classifying algebras
Thus let us address the task of determining the conformal boundary conditions for a CFT that is known at the chiral level. Un(?)fortunately the requirement of conformal invariance is rather weak, simply because $`𝔄`$ is typically much larger than just $`𝒱`$ir. As a result, there will in general (e.g. already for free boson theories) be infinitely many conformal boundary conditions. Usually they will be difficult to survey.
A pragmatic way out of this dilemma is to impose invariance under all of $`𝔄`$, or at least under a sufficiently large consistent chiral subalgebra $`\overline{𝔄}`$ of $`𝔄`$, rather than only under the Virasoro algebra $`𝒱`$ir. (Here ‘invariance’ means that the behavior at the boundary $`C`$, cf. formula (14), is identical for all bulk fields that are associated to vectors in a given $`\overline{𝔄}`$-submodule of an $`𝔄`$-module $`_\mu `$.) Then one can achieve a ‘rational’ situation, with only finitely many boundary conditions. Now by comparing two different factorization limits of the two-point function $`\varphi _{\mu ,\dot{\mu }}\varphi _{\nu ,\dot{\nu }}_a`$, one can show that
$$_{\lambda ,\dot{\lambda };\mathrm{\Omega }^{}}^a_{}_{\mu ,\dot{\mu };\mathrm{\Omega }^{}}^a_{}=_\nu \stackrel{~}{\mathrm{N}}_{\lambda ,\mu }^\nu _{\nu ,\dot{\nu };\mathrm{\Omega }^{}}^a_{}$$
(15)
with numbers $`\stackrel{~}{\mathrm{N}}_{\lambda ,\mu }^\nu `$ which are combinations of fusing matrices and operator product coefficients. At first glance, these expressions look very complicated. But there is a crucial insight: manifestly, $`\stackrel{~}{\mathrm{N}}_{\lambda ,\mu }^\nu `$ does not dependent on the boundary condition $`a`$.
This observation allows us to interpret the reflection coefficients $`_{\mu ,\dot{\mu };\mathrm{\Omega }}^a`$ as furnishing a one-dimensional irreducible representation of an algebra $`𝒞`$($`\overline{𝔄}`$) with structure constants $`\stackrel{~}{\mathrm{N}}_{\lambda ,\mu }^\nu `$, termed the classifying algebra. The results of may be summarized by the statement that the classifying algebra $`𝒞(𝔄)`$ for boundary conditions preserving the full bulk symmetry $`𝔄`$ (and with charge conjugation as torus partition function) is nothing but the fusion algebra of the CFT. Thus $`𝒞(𝔄)`$ is a semi-simple associative algebra, its structure constants are expressible through the Verlinde matrix $`S`$ as in (7), and both a basis of $`𝒞(𝔄)`$ and the boundary conditions $`a`$ are labelled by the set $`I`$ of chiral labels $`\mu `$. (Yet, an explicit verification of $`\stackrel{~}{\mathrm{N}}_{\lambda ,\mu }^\nu =\mathrm{N}_{\lambda ,\mu }^\nu `$ was achieved only in special cases where the relevant operator products and fusing matrices are known.)
When $`\overline{𝔄}𝔄`$, then the situation is more complicated, though the factorization arguments go through. For such symmetry breaking boundary conditions one finds:
One still has one-dimensional irreducible representations of some algebra $`𝒞=𝒞(\overline{𝔄})`$.
But the correlation functions are different. Namely, they are formed as different combinations of the chiral blocks for $`𝔄`$-descendant fields that are $`\overline{𝔄}`$-primaries.
The labelling $`\{\stackrel{~}{\mu }\}`$ of basis elements of $`𝒞`$ and $`\{a\}`$ of boundary conditions is more subtle. In particular the two sets of labels are distinct; both differ from the set $`I`$.
When the unbroken part of the bulk symmetries constitutes the fixed point algebra
$$\overline{𝔄}_{}=𝔄^G$$
(16)
with respect to any finite abelian group $`G`$ of automorphisms of $`𝔄`$, then the boundary conditions can be analyzed via $`G`$-orbifold and simple current techniques.
## 10 Interlude: Simple current extensions
One of the CFT concepts that was instrumental for arriving at conjectures (9) and (10) is the simple current extension of a rational CFT. It will show up again in the study of boundary conditions below. Assume that the following data are given: <sup>10</sup><sup>10</sup>10 While usually this is formulated by saying that one has some CFT with corresponding properties (and indeed there are many CFTs with those properties), here we need not refer directly to CFT.
A set $`\{\chi _\mu \}`$ ($`\mu I`$, $`|I|<\mathrm{}`$) of functions of $`\tau `$, convergent for $`\mathrm{}(\tau )>\mathrm{\hspace{0.17em}0}`$ and forming a basis of a unitary module $`V`$ over SL(2,$``$) for which $`S=S^\mathrm{t}`$ and $`T=\mathrm{diag}`$.
A vacuum label $`\mathrm{\Omega }I`$, satisfying $`S_{\mathrm{\Omega },\mu }_{>0}`$ for all $`\mu I`$, and an involution $`\mu \mu ^+`$ on $`I`$ such that $`\mathrm{\Omega }^+=\mathrm{\Omega }`$ as well as $`S_{\lambda ,\mu ^+}=S_{\lambda ,\mu }^{}`$ and $`T_{\mu ^+}=T_\mu `$ for all $`\lambda ,\mu I`$.
A subset $`𝒢I`$ such that $`S_{\mathrm{J},\mathrm{\Omega }}=S_{\mathrm{\Omega },\mathrm{\Omega }}`$ and $`T_\mathrm{J}=T_\mathrm{\Omega }`$ for all $`\mathrm{J}𝒢`$.
(In CFT terms: $`\mathrm{J}𝒢`$ has the same quantum dimension (namely unity) and the same conformal weight $`mod`$ (namely zero) as $`\mathrm{\Omega }`$, i.e. is an integer spin simple current.)
The numbers $`\mathrm{N}_{\lambda ,\mu }^\nu `$, regarded as defined by formula (7), are non-negative integers.
In this situation one defines a fusion ring with product ‘$``$’ on the vector space spanned by $`\{\phi _\mu |\mu I\}`$ by $`\phi _\lambda \phi _\mu :=_{\nu I}\mathrm{N}_{\lambda ,\mu ,\nu ^+}\phi _\nu `$, and can show rigorously :
$`𝒢`$ is a finite abelian group w.r.t. ‘$``$’ – the group of units of the fusion ring.
$`𝒢`$ organizes $`I`$ into orbits $`[\mu ]:=\{\mathrm{J}\mu |\mathrm{J}𝒢\}`$, with $`\varphi _\mathrm{J}\varphi _\mu =:\varphi _{\mathrm{J}\mu }\varphi _{\mathrm{J}\mu }`$.
Defining the stabilizer subgroup $`𝒮_\lambda :=\{\mathrm{J}𝒢|\mathrm{J}\lambda =\lambda \}`$, the combination
$$Z(\tau )=\underset{\genfrac{}{}{0pt}{}{[\mu ]:\mu I,}{T_{\mathrm{J}\mu }=T_\mu \mathrm{J}𝒢}}{}\text{(}|𝒮_\mu ||\underset{\mathrm{J}𝒢/𝒮_\mu }{}\chi _{\mathrm{J}\mu }(\tau )|^2\text{)}$$
(17)
is SL(2,$``$)-invariant. $`Z`$ is called a simple current extension modular invariant. <sup>11</sup><sup>11</sup>11 Many of these modular invariants are interesting. Examples include the $`D_{\mathrm{even}}`$ type invariants of the $`𝔰𝔩(2)`$ WZW model and the invariant $`Z=|\chi _1+\chi _{35}+\chi _{35^{}}+\chi _{35^{\prime \prime }}|^2+4|\chi _{28}|^2`$ for $`D_4`$ level 2.
To justify this name, one must be able to interpret (17) as the diagonal invariant for some extended CFT. This was achieved in , where the following was proven:
The extended labels are equivalence classes of pairs $`[\mu ,\widehat{\psi }]`$ with $`T_{\mathrm{J}\mu }=T_\mu `$ and $`\widehat{\psi }`$ a character of the untwisted stabilizer $`𝒰_\mu :=\{\mathrm{J}𝒮_\mu |F_\mu (\mathrm{J},\mathrm{J}^{})=1\mathrm{J}^{}𝒮_\mu \}𝒮_\mu `$. Here $`F_\mu `$ is an alternating bihomomorphism on $`𝒮_\mu `$, and hence the commutator cocycle $`F_\mu (\mathrm{J},\mathrm{J}^{})=_\mu (\mathrm{J},\mathrm{J}^{})/_\mu (\mathrm{J}^{},\mathrm{J})`$ for some cohomology class $`_\mu H^2(𝒮_\mu ,\mathrm{U}(1))`$ . Thus the group algebra $`𝒰_\mu `$ is isomorphic to the center of the twisted group algebra $`__\mu 𝒮_\mu `$, implying that the inclusion $`𝒰_\mu 𝒮_\mu `$ is of square index $`d_\mu ^2`$, with $`d_\mu `$ the dimension of the irreducible $`__\mu 𝒮_\mu `$-representations. <sup>12</sup><sup>12</sup>12 $`F_\mu `$ enters in calculations at various places; that for $`F_\mu \mathrm{\hspace{0.17em}1}`$ everything still nicely fits together is a strong consistency check. Non-trivial $`F_\mu `$ appear naturally via products of simple currents that individually have $`T_\mathrm{J}=T_\mathrm{\Omega }`$, e.g. when one deals with tensor products of subtheories, such as in Gepner type string compactifications.
In the $`D_4`$ example of footnote 11, one finds $`𝒮_{28}=𝒢=_2\times _2`$ but $`𝒰_{28}=\{\mathrm{\Omega }\}`$. Thus for $`\mu =\mathrm{\Lambda }_{(2)}\widehat{=}\mathrm{\hspace{0.17em}28}`$ there is only a single extended character $`\chi _{28}^{\mathrm{ext}}=\mathrm{\hspace{0.17em}2}\chi _{28}`$.
The summands in (17) are to be read as $`|𝒰_\mu ||\chi _{[\mu ]}^{\mathrm{ext}}|^2`$, i.e. for each $`[\mu ]`$ there are $`|𝒰_\mu |`$ many extended irreducible characters $`\chi _{[\mu ,\widehat{\psi }]}^{\mathrm{ext}}`$. Correspondingly the decompositions
$$_{[\mu ,\widehat{\psi }]}^{\mathrm{ext}}=_{\mathrm{J}𝒢/𝒮_\mu }^{d_\mu }_{\mathrm{J}\mu }$$
(18)
hold, and $`\chi _{[\mu ,\widehat{\psi }]}^{\mathrm{ext}}=d_\mu _{\mathrm{J}𝒢/𝒮_\mu }\chi _{\mathrm{J}\mu }`$ is the character of the extended module $`_{[\mu ,\widehat{\psi }]}^{\mathrm{ext}}`$.
Given, for every $`\mathrm{J}𝒢`$, a unitary matrix $`S^\mathrm{J}`$ satisfying the SL(2,$``$) relations as well as $`S_{\lambda ,\mu }^\mathrm{J}=S_{\mu ,\lambda }^{\mathrm{J}^1}`$ ($`\mathrm{J}𝒮_\lambda 𝒮_\mu `$) and $`S^\mathrm{\Omega }=S`$, the modular S-transformation matrix $`S^{\mathrm{ext}}`$ of the functions $`\chi _{[\mu ,\widehat{\psi }]}^{\mathrm{ext}}`$ is obtained by sandwiching the $`S^\mathrm{J}`$ between group characters:
$$\overline{)S_{[\lambda ,\widehat{\psi }_\lambda ],[\mu ,\widehat{\psi }_\mu ]}^{\mathrm{ext}}=\frac{\left|𝒢\right|^{}}{\left[\left|𝒮_\lambda \right|\left|𝒰_\lambda \right|\left|𝒮_\mu \right|\left|𝒰_\mu \right|\right]_{}^{1/2}}\underset{\mathrm{J}𝒰_\lambda 𝒰_\mu }{}\widehat{\psi }_\lambda (\mathrm{J})S_{\lambda ,\mu }^\mathrm{J}\widehat{\psi }_\mu (\mathrm{J})^{}.}$$
(19)
One has $`S_{\mathrm{J}^{}\lambda ,\mu }^\mathrm{J}=(T_\mu /T_{\mathrm{J}^{}\mu })F_\mu (\mathrm{J},\mathrm{J}^{})S_{\lambda ,\mu }^\mathrm{J}`$. $`S^{\mathrm{ext}}`$ is proven to be unitary and symmetric and to satisfy the SL(2,$``$) relations, and it was checked in a huge number of examples that it produces non-negative integers when inserted in the Verlinde formula .
There is evidence that $`S^\mathrm{J}`$ is the modular S-matrix for the one-point blocks on the torus with insertion $`\mathrm{J}`$. For WZW or coset models, $`S^\mathrm{J}`$ is the KacPeterson matrix of the orbit Lie algebra that is related <sup>13</sup><sup>13</sup>13 See , and also for background material and related work. to $`𝔤`$ by a folding of the Dynkin diagram.
Simple currents of WZW models correspond to the elements of the center of the relevant covering group G. It follows that (19) appears in the Verlinde formula for non-simply connected groups. (This result was checked in for some simple cases.)
## 11 The classifying algebra for finite abelian $`G`$
A systematic classification of boundary conditions has been achieved for all cases where $`\overline{𝔄}`$ is given as in (16), with finite abelian automorphism group $`G`$ . (The simplest case $`G=_2`$ includes e.g. Dirichlet boundary conditions for free bosons.)
All basic ingredients are already known <sup>14</sup><sup>14</sup>14 In particular one can make use of the fact that simple current extension by $`𝒢G^{}`$ provides the inverse operation to forming the orbifold with respect to the finite abelian group $`G`$. This way one can exploit both orbifold techniques and the simple current framework sketched in section 10. from chiral CFT. In particular:
The label sets $`\{\stackrel{~}{\mu }\}`$ for the basis of $`𝒞(\overline{𝔄})`$ and $`\{a\}`$ for boundary conditions arise as two different deviations from the labels appearing in (19). For $`\stackrel{~}{\mu }`$, one has a character of $`𝒮_\mu `$ rather than of $`𝒰_\mu `$, and no orbit is to be taken, but still the requirement $`T_{\mathrm{J}\mu }=T_\mu `$ is kept, while $`a`$ has the same form as extended labels, but now $`T_{\mathrm{J}\mu }T_\mu `$ is allowed:
$$\overline{)\stackrel{~}{\mu }=(\overline{\mu },\psi )}\mathrm{and}\overline{)a=[\overline{\rho }(a),\widehat{\psi }]_{}^{}}\mathrm{with}\psi 𝒮_{\overline{\mu }}^{},\widehat{\psi }𝒰_{\overline{\rho }(a)}^{}$$
(20)
(recall $`𝒰_{\overline{\lambda }}𝒮_{\overline{\lambda }}𝒢`$). (20) follows by heuristic considerations resembling ideas in .
Comparing with (19), we can make an educated guess for a diagonalizing matrix:
$$\overline{)\stackrel{~}{S}_{(\overline{\lambda },\psi _\lambda ),[\overline{\rho },\widehat{\psi }_\rho ]}=\frac{\left|𝒢\right|^{}}{\left[\left|𝒮_{\overline{\lambda }}\right|\left|𝒰_{\overline{\lambda }}\right|\left|𝒮_{\overline{\rho }}\right|\left|𝒰_{\overline{\rho }}\right|\right]_{}^{1/2}}\underset{\mathrm{J}𝒮_{\overline{\lambda }}𝒰_{\overline{\rho }}}{}\psi _{\overline{\lambda }}(\mathrm{J})S_{\overline{\lambda },\overline{\rho }}^\mathrm{J}\widehat{\psi }_{\overline{\rho }}(\mathrm{J})^{}.}$$
(21)
Then $`𝒞(\overline{𝔄})`$ is defined by prescribing the Verlinde-like formula featuring $`\stackrel{~}{S}`$:
$$\stackrel{~}{\mathrm{N}}_{\stackrel{~}{\lambda },\stackrel{~}{\mu },\stackrel{~}{\nu }}:=_a\stackrel{~}{S}_{\stackrel{~}{\lambda },a}\stackrel{~}{S}_{\stackrel{~}{\mu },a}\stackrel{~}{S}_{\stackrel{~}{\nu },a}/\stackrel{~}{S}_{\stackrel{~}{\mathrm{\Omega }},a}.$$
(22)
The structure constants are obtained from (22) by raising the third index via $`\stackrel{~}{\mathrm{N}}_{\stackrel{~}{\lambda },\stackrel{~}{\mu },\stackrel{~}{\mathrm{\Omega }}}`$.
There is as yet no rigorous derivation of formula (21). But once (21) and (22) are taken for granted, $`𝒞(\overline{𝔄})`$ can be studied with full rigor. In particular one shows :
$`\stackrel{~}{S}`$ is (weighted) unitary. (Note that it is even non-trivial that $`\stackrel{~}{S}`$ is a square matrix.)
$`𝒞(\overline{𝔄})`$ is a semi-simple commutative associative algebra with unit element $`\stackrel{~}{\mathrm{\Omega }}=\overline{\mathrm{\Omega }}`$ (the vacuum sector of the $`G`$-orbifold).
The structure constants $`\stackrel{~}{\mathrm{N}}_{\stackrel{~}{\lambda },\stackrel{~}{\mu }}^{\stackrel{~}{\nu }}`$ of $`𝒞(\overline{𝔄})`$ are diagonalized by the matrix (21). The irreducible $`𝒞(\overline{𝔄})`$-representations $`R_a`$ are one-dimensional and labelled by the boundary labels $`a`$; they yield the reflection coefficients as $`_{\stackrel{~}{\mu },\stackrel{~}{\mu }_{}^+;\stackrel{~}{\mathrm{\Omega }}}^a=R_a(\varphi _{\stackrel{~}{\mu }})=\stackrel{~}{S}_{\stackrel{~}{\mu },a}/\stackrel{~}{S}_{\stackrel{~}{\mathrm{\Omega }},a}`$.
As an algebra over $``$, $`𝒞(\overline{𝔄})`$ decomposes into ideals as
$$𝒞(\overline{𝔄})_{gG}𝒞^{(g)}(\overline{𝔄}).$$
(23)
The ideal $`𝒞^{(g)}(\overline{𝔄})`$ plays the role of a classifying algebra for boundary conditions of definite automorphism type $`g`$. The corresponding boundary states are linear combinations of $`g`$-twisted chiral blocks, which obey $`g`$-twisted Ward identities. <sup>15</sup><sup>15</sup>15 For consistent subalgebras that are not fixed point algebras, there exist boundary conditions which do not possess an automorphism type. Examples of such boundary conditions are e.g. known for the $`_2`$-orbifold of a free boson and for the $`E_6`$-type invariant of the $`𝔰𝔩(2)`$ WZW model.
The ideal $`𝒞^{(e)}(\overline{𝔄})`$ appearing in (23) is precisely the fusion rule algebra of $`𝔄`$.
It is plausible that orbifolding can be understood in terms of the folding of fusion graphs, and that the classifying algebra $`𝒞(\overline{𝔄})`$ thus coincides with the corresponding Pasquier algebra . <sup>16</sup><sup>16</sup>16 At least for cyclic $`G`$ – for non-cyclic $`G`$ one must be aware of the possibility of having non-trivial two-cocycles $`_\mu `$.
Also, in practice this is difficult to check, because the Pasquier algebra is obtained by an algorithm which does not directly produce uniform formulae for all rational CFTs. The identification seems to be established so far only for $`G=_2`$ in $`𝔰𝔩(2)`$ WZW models and Virasoro minimal models.
The statements above refer to the charge conjugation torus partition function. More general results follow via T-duality symmetries, which are similar to those of free boson theories, acting compatibly on the boundary conditions and on the torus partition function.
Knowing the classifying algebra $`𝒞(\overline{𝔄})`$ explicitly, a variety of consistency checks can be made. Most importantly, one can prove integrality of the coefficients of characters in the annulus amplitude (open string partition function). This integrality is often used as the starting point for studying boundaries; here it rather serves as an independent check.
Finally we remark that one can express the structure constants $`\stackrel{~}{\mathrm{N}}_{\stackrel{~}{\lambda },\stackrel{~}{\mu }}^{\stackrel{~}{\nu }}`$ as well as the annulus coefficients through traces of twisted intertwining operators $`\mathrm{\Theta }_\stackrel{}{\sigma }`$ on chiral block spaces. This yields the announced connection to the topic studied in section 5.
## 12 Conclusions and outlook
Let us summarize by telling what we regard as the two main messages:
First, there exists a close relation between the sub-bundle structure of chiral blocks (leading to the trace formula (10)) and symmetry breaking boundary conditions.
Second, there is a systematic classification of all boundary conditions leaving unbroken a fixed point algebra $`\overline{𝔄}=𝔄^G`$ with respect to an arbitrary finite abelian group $`G`$. Concretely, one has a general prescription, valid for all rational CFTs, for the classifying algebra, with structure constants expressed through known chiral data.
Our results illustrate that the space of boundary conditions has a rich and unexpectedly nice structure. We believe that many more issues are accessible quantitatively.
Among possible extensions of the work outlined above we mention:
One should find the diagonalizing matrix $`\stackrel{~}{S}`$ of $`𝒞(\overline{𝔄})`$ when $`\overline{𝔄}𝔄^G`$ for any group $`G`$.
2-d boundary conditions can be understood in terms of 3-d topological theory .
Non-orientable surfaces, e.g. one-point functions on $`^2`$ and the partition functions of the Klein bottle and the Möbius strip, are studied in and .
One should look for a geometric interpretation of boundary conditions for non-flat backgrounds. For WZW models this is indeed available: one obtains ‘fuzzy’ versions of (possibly twisted) conjugacy classes of the group manifold G .
More explicit information on the chiral data for further classes of models is highly welcome.
Applications to string theory include, e.g., the complete analysis of concrete compactifications and a more systematic understanding of tadpole cancellation. |
warning/0001/astro-ph0001407.html | ar5iv | text | # Broad-band BeppoSAX observation of the low-mass X-ray binary X 1822-371
## 1 Introduction
The 5.57 hr low-mass X-ray binary (LMXRB) X 1822-371 is viewed almost edge on with the central X-ray source hidden by material in the orbital plane. The X-ray lightcurve shows a partial eclipse and a smooth broad modulation with a minimum just prior to the eclipse (White et al. w:81 (1981)). The partial nature of the eclipse indicates that the X-ray emitting region is extended, and that the observed X-rays are scattered in an accretion disk corona, or ADC (White & Holt w:82 (1982)). The optical lightcurve has a similar morphology except that the eclipse is broader. White & Holt (w:82 (1982)) showed that the broad X-ray modulation can be modeled as obscuration of the ADC by the rim of the accretion disk whose thickness is greatest near phase, $`\mathrm{\Phi }`$, 0.8 (where $`\mathrm{\Phi }=0.0`$ is mid-eclipse), and least near $`\mathrm{\Phi }=0.2`$. Modeling of multi-wavelength lightcurves reveals that the effective diameter of the ADC is $`3\times 10^{10}`$ cm, about half that of the optically emitting disk, and that structure in the disk can reach a height of $`1.5\times 10^{10}`$ cm (White & Holt w:82 (1982); Mason & Córdova mc:82 (1982); Hellier & Mason h:89 (1989); Puchnarewicz et al. p:95 (1995)). No pulsations or bursts have been detected from X 1822-371 (e.g., Hellier et al. h:92 (1992)). The probable distance to X 1822-371 is 2–3 kpc and the isotropic luminosity $`10^{36}`$(d<sub>2</sub>)<sup>2</sup> erg s<sup>-1</sup> (Mason & Córdova mc:82 (1982)), where d<sub>2</sub> is distance in units of 2 kpc. The mean X 1822-371 $`\mathrm{L}_\mathrm{x}/\mathrm{L}_{\mathrm{opt}}`$ ratio of 20 compared to the average for LMXRB of 500 (van Paradijs & McClintock v:95 (1995)) implies a unobscured luminosity of $`2\times 10^{37}`$(d<sub>2</sub>)<sup>2</sup> erg s<sup>-1</sup>, similar to that observed from the brighter X-ray burst sources.
The X-ray spectrum of X 1822-371 is complex with different results being obtained from different measurements. The 2–40 keV HEAO-1 A2 spectrum can be fit by a power-law with a photon index, $`\alpha `$, of 1.3, a high-energy cutoff at 18 keV, and a broad Fe-K line with a full-width half maximum (FWHM) of 4 keV. Below 2 keV there is evidence for an excess in the Einstein Solid State Spectrometer spectrum which White et al. (w:81 (1981)) can model as a 0.25 keV thermal bremsstrahlung, a 0.16 keV blackbody or a 350 eV equivalent width (EW) emission feature at 0.8 keV. The combined 0.04–20 keV EXOSAT Channel Multiplier Array and Medium Energy Detector Array spectrum is fit by a power-law with $`\alpha =0.8`$ together with a blackbody with kT = 1.8 keV and a $``$1 keV FWHM Fe-K line (Hellier & Mason h:89 (1989)). Subsequently, Hellier et al. (h:92 (1992)) obtained a Ginga 1.5–30 keV spectrum, but were unable to find a satisfactory fit, although the model used for the EXOSAT spectrum gives the best result.
White et al. (w:97 (1997)) report on an ASCA observation of X 1822-371. The best-fit Solid-State Imaging Spectrometer (SIS) and Gas Imaging Spectrometer (GIS) 0.5–10 keV spectrum is also complex and cannot be easily fit by any of the standard models. White et al. (w:97 (1997)) characterize the spectrum using a power-law continuum with $`\alpha 0.52`$. Structured residuals remain with a pronounced dip at $``$1.5 keV and a strong decrease $`\text{ }>`$7 keV. There are no emission features present other than two weak features at 6.4 keV and 7.1 keV, consistent with Fe K$`\alpha `$ and K$`\beta `$ lines. The EW of the K$`\alpha `$ line is $``$40 eV. However, as White et al. (w:97 (1997)) note, the observed K$`\beta `$/K$`\alpha `$ ratio of $``$40% does not match the theoretical value of 13%. White et al. (w:97 (1997)) attempt to model the reduction in observed flux $`\text{ }>`$7 keV by an Fe-K absorption edge, but find that edges from a range of ionization states are necessary to explain the observed shape.
As White et al. (w:97 (1997)) point out, accurately measuring the continuum $`\text{ }>`$10 keV is vital in order to constrain the Fe-K absorption feature seen in the ASCA spectrum. Indeed, these authors used non-simultaneous EXOSAT and Ginga data to try to do this, but were unsuccessful. We report here on a BeppoSAX observation of X 1822-371 where the 0.3–40 keV spectrum was observed simultaneously with good sensitivity. We compare the results of the BeppoSAX spectral fits with those obtained by White et al. (w:97 (1997)) and those obtained from a subsequent ASCA observation.
## 2 Observations
Results from the Low-Energy Concentrator Spectrometer (LECS; 0.1–10 keV; Parmar et al. p:97 (1997)), the Medium-Energy Concentrator Spectrometer (MECS; 1.8–10 keV; Boella et al. b:97 (1997)), the High Pressure Gas Scintillation Proportional Counter (HPGSPC; 5–120 keV; Manzo et al. m:97 (1997)) and the Phoswich Detection System (PDS; 15–300 keV; Frontera et al. f:97 (1997)) on-board BeppoSAX are presented. All these instruments are coaligned and collectively referred to as the Narrow Field Instruments, or NFI. The MECS consists of two grazing incidence telescopes with imaging gas scintillation proportional counters in their focal planes. The LECS uses an identical concentrator system as the MECS, but utilizes an ultra-thin entrance window and a driftless configuration to extend the low-energy response to 0.1 keV. The non-imaging HPGSPC consists of a single unit with a collimator that was alternatively rocked on- and 180′ off-source every 96 s during the observation. The non-imaging PDS consists of four independent units arranged in pairs each having a separate collimator. Each collimator was alternatively rocked on- and 210′ off-source every 96 s during the observation.
The region of sky containing X 1822-371 was observed by BeppoSAX on 1997 September 09 12:24 UT to September 10 11:47 UT. Good data were selected in the standard way using the SAXDAS 2.0.0 data analysis package. LECS and MECS data were extracted centered on the (on-axis) position of X 1822-371 using radii of 8′ and 4′, respectively. The exposures in the LECS, MECS, HPGSPC, and PDS instruments are 14.8 ks, 37.0 ks, 19.5 ks, and 18.7 ks, respectively. Background subtraction for the imaging instruments was performed using standard files, but is not critical for such a bright source. Background subtraction for the HPGSPC was carried out using data obtained when the instrument was looking at the dark Earth and for the PDS using data obtained during intervals when the collimator was offset from the source.
The BeppoSAX data is compared with results from the Solid State Imaging Spectrometers SIS0 and SIS1 (0.6–10 keV), on-board ASCA (Tanaka et al. t:94 (1994)). The energy resolution of the SIS is a factor of a few better, except at the lowest energies, than that of the LECS and MECS. ASCA observed X 1822-371 twice. The first observation took place between 1993 October 07 03:33 and 23:55 UTC and the second between 1996 September 26 05:44 and September 27 08:05 UTC. The SIS exposures for each observation are 36.6 ks and 26.0 ks both using 1-CCD BRIGHT mode. All data were screened and processed using the standard Rev2 pipeline. The source count rates of $`\text{ }<`$6.5 s<sup>-1</sup> SIS<sup>-1</sup> means that pulse pile-up is unlikely to be significant.
## 3 Analysis and results
### 3.1 BeppoSAX lightcurve
Background subtracted lightcurves in the energy ranges 1.8–4 keV (MECS), 4–10 keV (MECS), 10–16 keV (HPGSPC) and 16–32 keV (PDS) were extracted and folded using the linear ephemeris of Hellier & Smale (h:94 (1994)). Fig. 1 shows these lightcurves together with two hardness ratios (4–10 keV counts divided by 1.8–4 keV counts and 16–32 keV counts divided by 10–16 keV counts). The lightcurves are similar to those reported previously, showing gradual reductions in flux which reach minima around $`\mathrm{\Phi }0.8`$, before the partial eclipse and rapid increases following egress. Just prior to the partial eclipse each lightcurve shows a narrow increases in flux. The depth of the partial eclipse increases with increasing energy. There are no strong variations in hardness ratio. However, the 4–8 keV/1.8–4 keV hardness ratio (HR1) shows a broad variation with a maximum around $`\mathrm{\Phi }0.4`$, as well as a narrow increase just prior to eclipse. Instead, the 16–32 keV/8–16 keV hardness ratio (HR2) does not vary appreciably with phase, except between $`\mathrm{\Phi }=0.7`$–1.1 where a harder interval, punctuated by the eclipse, is present. HR1 is similar in shape to the 6–30 keV/1–6 keV hardness ratio obtained from Ginga observations (see Fig. 1 of Hellier et al. h:92 (1992)), except that no strong softening is visible in the BeppoSAX data during the eclipse. This is in contrast to previous results, where the eclipse depths measured by EXOSAT and Ginga both showed a similar strong dependence on energy (Hellier et al. h:92 (1992)). The lack of any large change in the BeppoSAX hardness ratios justifies the use of the entire data set in the spectral analysis in Sect. 3.3.
### 3.2 Eclipse timing
Eclipse arrival times using recent BeppoSAX MECS, ASCA and RXTE Proportional Counter Array data have been determined. Since there are substantial gaps in the data and the count rates low, the BeppoSAX and ASCA lightcurves were first folded on the 5.5706 hr period of Hellier & Smale (h:94 (1994)). A model consisting of a Gaussian and a constant was then fit to the eclipse phases ($`\mathrm{\Phi }=0.95`$–1.05) of the folded lightcurves. The arrival time of the eclipse was then assigned to the eclipse which occurred closest to mid-time of the observation. For the MECS the uncertainty on the arrival time was obtained directly from the fit. Data from the 4 ASCA instruments were analyzed independently and the uncertainties in the arrival times were derived from the spreads in the obtained values, while the arrival times were defined as the averages from the 4 instruments. The RXTE observation was made between 1996 September 26 16:03 and September 27 07:19 UTC. Since only one eclipse was observed, the Gaussian and a constant model was fit directly to the eclipse interval in the lightcurve. The uncertainty in the RXTE arrival time was estimated by comparing results obtained using (1) a range of eclipse phases around those given above using the Gaussian and constant model and (2) including a linear term in the above model and fitting over the same range of phases. The spread in values obtained using these two methods was larger than the uncertainties obtained from the individual fits, and so was adopted as an estimate of the overall uncertainty. The fit results, which extend the measurements in Hellier & Smale (h:94 (1994)) by some 5 yrs, or 8000 cycles, are summarized in Table LABEL:tab:eclipse\_times.
The newly determined arrival times together with the measurements tabulated in Hellier & Smale (h:94 (1994)) were fit to obtain an updated ephemeris. Both linear ($`\chi ^2`$ = 98.5 for 17 dof) and quadratic ephemerides ($`\chi ^2`$ = 21.4 for 16 dof) were used. The difference in reduced $`\chi ^2`$ indicates that the quadratic term in the ephemeris is established with 99.9999% confidence. Fig. 2 shows the residuals with respect to both ephemerides. The updated quadratic ephemeris is given by:
$`\mathrm{T}_{\mathrm{ecl}}`$ $`=`$ $`2445615.30964(15)+0.232108785(50)\mathrm{N}`$
$`+2.06(23)\times 10^{11}\mathrm{N}^2.`$
### 3.3 BeppoSAX spectrum
The overall spectrum of X 1822-371 was first investigated by simultaneously fitting data from all the BeppoSAX NFI. All spectra were rebinned using standard procedures. Data were selected in the energy ranges 0.3–4.0 keV (LECS), 2.0–10 keV (MECS), 7.0–30 keV (HPGSPC), and 15–40 keV (PDS) where the instrument responses are well determined and sufficient counts obtained. This gives background-subtracted count rates of 1.0, 4.2, 8.8 and 4.0 s<sup>-1</sup> for the LECS, MECS, HPGSPC, and PDS, respectively. The photo-electric absorption cross sections of Morisson & McCammon (m:83 (1983)) and the abundances of Anders & Grevesse (a:89 (1989)) are used throughout. Factors were included in the spectral fitting to allow for normalization uncertainties between the instruments. These factors were constrained to be within their usual ranges during the fitting. All spectral uncertainties and upper-limits are given at 90% confidence.
Initially, a standard LMXRB spectral model consisting of a a cutoff power-law ($`\mathrm{E}^\alpha \mathrm{exp}(\mathrm{E}_\mathrm{c}/\mathrm{kT})`$) appropriate to low-luminosity sources was fit to the BeppoSAX spectrum. Since the central X-ray source in X 1822-371 is hidden from direct view the overall luminosity of X 1822-371 is highly uncertain, and an additional blackbody component, frequently observed in more luminous LMXRB (e.g., White et al. w:88 (1988)) was also included. The cutoff power-law model gives an unacceptable fit with a $`\chi ^2`$ of 670 for 133 dof. Adding a 1.5 keV blackbody reduces the $`\chi ^2`$ to 596 for 131 dof. Inspection of the residuals indicated the presence of an Fe-K line. Adding a broad line at $``$6.5 keV further reduces the $`\chi ^2`$ to 408 for 128 dof. The remaining residuals are indicative of an intense low-energy excess. We attempted to model this feature using a Gaussian emission line, thermal bremsstrahlung, blackbody and power-law components. Only the fit including a Gaussian feature was acceptable with a $`\chi ^2`$ of 144.5 for 125 dof. Thus, this model consists of a cutoff power-law with $`\alpha =0.78\pm 0.02`$ and $`\mathrm{E}_\mathrm{c}=5.69\pm 0.06`$ keV and a $`1.27\pm 0.03`$ keV blackbody continuum together with two Gaussian emission features. The nature of the low-energy feature is highly uncertain. Its energy can only be constrained to be $`<`$0.50 keV, it is extremely broad (FWHM = 1.5 keV), and its EW is implausibly high ($``$7 keV). We therefore reject this model as being physically unreasonable.
The HEAO-1 X 1822-371 spectrum of White et al. (w:81 (1981)) is well fit above 2 keV using a power-law model with a high-energy break ($`\mathrm{exp}((\mathrm{E}_{\mathrm{cut}}\mathrm{E})/\mathrm{E}_{\mathrm{fold}})`$ for $`\mathrm{E}>\mathrm{E}_{\mathrm{cut}}`$) at 17 keV. Fitting this model to the BeppoSAX NFI spectrum does not produce an acceptable fit with a $`\chi ^2`$ of 1960 for 132 dof. If a blackbody is included, then the situation improves with a $`\chi ^2`$ of 342 for 130 dof. Including a broad Fe line decreases the $`\chi ^2`$ to 198.9 for 127 dof. The major contribution to $`\chi ^2`$ comes from the region near 1 keV where the model overestimates the observed spectrum. We attempted to model feature this using partial covering (the pcfabs model in xspec, an ionized absorber the absori model in xspec), and an absorption edge. Only the edge was successful with an energy of $`1.30\pm 0.07`$ keV and an optical depth, $`\tau `$, of $`0.35\pm 0.09`$ for a $`\chi ^2`$ of 139.4 for 125 dof. The best-fit value of $`\mathrm{E}_{\mathrm{cut}}`$ of $`17.1\pm 0.7`$ keV is in good agreement with that obtained from HEAO-1 of $`17.4\pm 0.5`$ keV (White et al. w:81 (1981)). There is no evidence for the presence of low-energy thermal components, Fe-L emission near 1 keV as proposed by White et al. (w:81 (1981)), or a reflection component (the pexrav model in xspec).
Finally, we investigated whether Comptonization models such as the xspec model comptt (Titarchuk ti:94 (1994); Hua & Titarchuk h:95 (1995); Titarchuk & Lyubarskij t:95 (1995)), which self-consistently calculate the spectrum produced by the Comptonization of soft photons in a hot plasma, could be applicable to X 1822-371. This model contains as free parameters the temperature of the Comptonizing electrons kT<sub>e</sub>, the plasma optical depth with respect to electron scattering $`\tau _\mathrm{p}`$ and the input temperature of the soft photon (Wien) distribution kT<sub>W</sub>. A spherical geometry was assumed for the comptonizing region. This replaced the power-law and high-energy cutoff in the previous model and the fit repeated. A good fit is obtained with a $`\chi ^2`$ of 122.1 for 125 dof. The 1–10 keV flux is $`5.0\times 10^{10}`$ erg cm<sup>-2</sup> s<sup>-1</sup> and the blackbody contributes 43% of the total flux in this energy range. The best-fit parameters are given in Table LABEL:tab:spec\_paras. If a disk geometry is assumed for the Comptonizing region, the best-fit parameters are essentially unchanged except for $`\tau _p`$ which reduces to $`13.3\pm 0.3`$. Fig. 3 shows the count spectra and the data/model ratios with, and without, the edge. The depth of the feature ($``$20%) is significantly greater than the uncertainty in LECS calibration at the corresponding energies ($`\text{ }<`$5%). Fig. 4 shows the incident photon spectrum and illustrates how fitting a power-law in the energy range 1.0–10 keV would produce a “dip” at $``$1.5 keV as observed by White et al. (w:97 (1997)). The 90% confidence upper limits to the optical depths of any O-K and Fe-K edges at 0.54 keV and 7.1 keV are 0.18 and 0.05, respectively.
### 3.4 ASCA spectrum
We next examined whether the BeppoSAX best-fit model presented above is consistent with results from ASCA. For the 1993 ASCA observation these are presented in White et al. (w:97 (1997)). A simple power-law was first fit to the ASCA SIS spectra confirming that the overall shape of the spectrum is as reported in White et al. (w:97 (1997)) when using the latest processing. Next, the best-fit BeppoSAX model was fit to the SIS spectra. For the fit to the 1993 spectra an additional line feature at $``$7 keV was included. The results are presented in Table LABEL:tab:asca and show that the best-fit BeppoSAX model also provides a reasonable fit to the ASCA SIS spectra. During the 1993 and 1996 observations, the 1–10 keV fluxes were $`5.8\times 10^{10}`$ erg cm<sup>-2</sup> s<sup>-1</sup> and $`5.3\times 10^{10}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, and the blackbody contributed 48% and 47% of the flux in this energy range, respectively.
Comparing the best-fit parameters with those obtained with BeppoSAX reveals the following: (1) The overall spectral shape is similar in all three observations and may be modeled using comptt and blackbody continuum components. (2) There is a substantial (EW $``$100–300 eV) Fe-K feature present which may be modeled as a broad feature at $``$6.5 keV, or as the sum of two narrower features at $``$6.4 keV and $``$7.0 keV. (3) Using the relation between optical and X-ray extinctions of Predehl & Schmitt (ps:95 (1995)), all three best-fit absorptions are inconsistent with the value derived from the color excess of 0.13 in Mason et al. (m:82 (1982)) of $`8\times 10^{20}`$ atom cm<sup>-2</sup>. (4) All three spectra reveal clear evidence for a “dip” at $``$1.3 keV, which may be modeled by an absorption edge. The energies of this feature ($`1.33\pm _{0.11}^{0.05}`$ keV, $`1.37\pm 0.03`$ keV, and $`1.28\pm 0.03`$ keV) are in reasonable agreement, while there are clear variations in its optical depth ($`\tau =0.28\pm 0.06`$, $`0.093\pm 0.012`$, and $`0.15\pm 0.02`$).
## 4 Discussion
We present results of a 1997 September BeppoSAX observation of X 1822-371 and compare these with earlier ASCA observations when the source had a similar 1.0–10 keV intensity. The spectrum is unusually complex and cannot be fit by any of the usual models applied to LMXRB such as a cutoff power-law and blackbody unless an unusually strong low-energy emission feature is included. A good fit is obtained to the 0.3–40 keV BeppoSAX spectrum with the combination of a Comptonization component and a blackbody together with an Fe-K emission line and an absorption edge. The same model provides reasonable fits to ASCA SIS spectra with a similar absorption feature being required.
There are at least three highly unusual features of the X 1822-371 spectrum. The first, also pointed out by White et al. (w:97 (1997)), is that the continuum is much harder than is typical of similar luminosity LMXRBs, where $`\alpha `$ is usually 1.5–2.5. The spectrum scattered in an ADC is expected to resemble the original, unless the optical depth is large. The second is the extremely large contribution of the blackbody component ($`\text{ }>`$40% of the total), whereas fits to luminous LMXRB indicate the presence of blackbodies with typical luminosities of 16–34% of the non-thermal component (White et al. w:88 (1988)). However, we note that another LMXRB, the X-ray dip source XB 1746-371 located in the globular cluster NGC 6441, also recently been found to have a strong blackbody-like component which contributes 88% of the 1–10 keV flux (Parmar et al. p:99 (1999); see also Guainazzi et al. g:99 (2000)). The blackbody component most probably originates in an optically thick boundary layer between the accretion disk and the neutron star surface, or from the neutron star itself. It is therefore unlikely to be observed directly in X 1822-371 and so it is surprising that it appears so bright in comparison with the non-thermal component. The third unusual aspect of the spectrum is the presence of the strong low-energy feature, which can be modeled as an absorption edge. The fact that a similar component is also seen in the two ASCA observations implies that this is a stable feature of the spectrum, but not necessarily that it is correctly modeled. However, its nature is highly uncertain. The energy corresponds to K-edges of highly ionized Ne x and neutral Mg, or to an L-edge of moderately ionized Fe. Surprisingly, no strong ($`\tau >0.05`$) Fe-K or ($`\tau >0.18`$) O-K edges are evident in the spectrum.
Models for the spectrum of X 1822-371 which involve a significant Comptonized component appear plausible. A plasma with kT$`{}_{\mathrm{e}}{}^{}`$5–10 keV and a $`\tau _\mathrm{p}`$ of $``$22–26 is required. These values imply a Comptonization parameter $`\mathrm{y}=4\mathrm{k}\mathrm{T}_\mathrm{e}\tau _\mathrm{p}^2/\mathrm{m}_\mathrm{e}\mathrm{c}^2`$ of $``$25–40. Values of y$`>`$12 imply that the emerging spectrum will be saturated and have a Wien-like shape (e.g., Titarchuk t:94 (1994)). Guainazzi et al. (g:99 (2000)) show that when the BeppoSAX spectra of a number of LMXRB located in globular clusters are fit with the same continuum model as applied here (a blackbody and comptt), then the derived values of $`\tau _\mathrm{p}`$ and $`\mathrm{kT}_\mathrm{e}`$ are respectively correlated and anti-correlated with the source luminosity. Guainazzi et al. (g:99 (2000)) suggest that these correlations may be qualitatively explained if the X-ray emission at the boundary layer between the accretion disk and the neutron star surface is proportional to the accretion rate. If this results in an increase in the Comptonizing plasma optical depth this would allow Compton cooling to become more efficient, yielding a lower Comptonizing electron temperature. If the intrinsic luminosity of X 1822-371 is assumed to be $`2\times 10^{37}`$ erg s<sup>-1</sup> (see Sect. 1), then the values derived here using BeppoSAX for $`\tau _\mathrm{p}`$ and $`\mathrm{kT}_\mathrm{e}`$ are in good agreement with the relations derived from the study of the globular cluster LMXRB X-ray sources.
The overall X-ray spectrum of X 1822-371 remains poorly understood. In particular, calculations by e.g., Ko & Kallman (k:94 (1994)) and Kallman et al. (k:96 (1996)) show that photo-ionized ADC should be a rich source of line emission, which does not appear to be the case here as only a moderate EW Fe-K line is seen. Vrtilek et al. (v:93 (1993)) discuss the situation where an X-ray source is viewed at high inclination through a moderate optical depth ADC. They predict both a deep Fe-K edge and a prominant K-$`\alpha `$ emission line - again neither of which are seen. Thus the X 1822-371 spectrum cannot be easily understood in terms of current models of X-ray production and reprocessing in ADCs and we await future high quality measurements to shed more light on this complex spectrum.
###### Acknowledgements.
The BeppoSAX satellite is a joint Italian and Dutch programme. We thank the referee, Martin Still, for helpful comments and Ken Ebisawa, Koji Mukai, and the staff of the BeppoSAX Science Data Center for assistance. This research has made use of data obtained through the High Energy Astrophysics Science Archive Research Center Online Service, provided by the NASA/Goddard Space Flight Center. |
warning/0001/math-ph0001022.html | ar5iv | text | # Abstract
## Abstract
Within the framework of Lanczos-inspired perturbation theory wave functions and energies in the short-range potential $`V(x)=a\rho ^2(x)+\lambda \rho ^4(x)`$ with $`\rho (x)=\mathrm{sech}\alpha x`$ and a small coupling $`\lambda `$ are shown obtainable in closed form.
PACS 03.65.Ge 02.30.Gp 02.30.Hq 03.65.Db
Pöschl-Teller potential $`V^{(0)}(x)=\mu (\mu +1)\mathrm{sech}^2\alpha x`$ resembles the celebrated harmonic oscillator. Within the so called shape invariant family these two spatially symmetric exactly solvable potentials form a unique subset. In perturbation theory there seems to emerge one of the most significant differences between them. In contrast to an enormous interest in the various anharmonic forces there exists virtually no analysis of a perturbed Pöschl-Teller model in the current literature. Partially, we intend to fill the gap. This letter is concerned with the most elementary quartic example
$$\left[\frac{d^2}{dx^2}\frac{\mu (\mu +1)}{\mathrm{cosh}^2\alpha x}+\frac{4\lambda }{\mathrm{cosh}^4\alpha x}\right]\psi _{(PPT)}(x)=\kappa _{(PPT)}^2\psi _{(PPT)}(x).$$
(1)
The parity is conserved, $`\psi _{(PPT)}(x)=(1)^p\psi _{(PPT)}(x)`$, $`p=0,1`$, and all the unperturbed $`\lambda =0`$ bound states are available in closed form. Unfortunately, after we re-scale $`\alpha 1`$ for simplicity we immediately notice that the normalizable states form a mere finite set with $`\kappa _{(PPT)}^{(0)}=\mu 2Np=\kappa (N,p)`$ and wave functions
$$x|\psi _{2N+p}^{(0)}=\frac{\mathrm{tanh}^px}{\mathrm{cosh}^{\kappa (N,p)}x}_2F_1(\mu N+\frac{1}{2},N,1+\kappa (N,p),\frac{1}{\mathrm{cosh}^2x})$$
(2)
where $`02N+p<\mu `$ . This is the reason why these mutually orthogonal elementary functions cannot be used as an unperturbed basis.
A key to our present $`\lambda 0`$ construction will lie in the use of a non-orthogonal basis. We shall employ the ansatz
$$\psi _{(PPT)}(x)=\mathrm{tanh}^px\underset{n=0}{\overset{\mathrm{}}{}}\frac{c_n(\lambda )}{\mathrm{cosh}^{2n+\kappa _{(PPT)}}x}$$
(3)
with $`\kappa _{(PPT)}=\kappa (N,p)+2\epsilon (\lambda )`$. Its use is inspired by the general method of Lanczos and its perturbative implementations . The basis itself is taken from the particular terminating expansions (2) which define the unperturbed coefficients $`c_n^{(0)}`$ such that $`c_N^{(0)}0`$ and $`c_0^{(0)}=1`$ while $`c_{N+1}^{(0)}=c_{N+2}^{(0)}=\mathrm{}=0`$.
As always in similar constructions one inserts eq. (3) in our differential eq. (1). This gives the current recurrence relations which may be written in the three-term form
$$\lambda c_{n1}+\beta _nc_n+\alpha _{n+1}c_{n+1}=0,n=0,1,\mathrm{}$$
(4)
or as an infinite-dimensional linear algebraic problem
$$\left(\begin{array}{cccc}\beta _0& \alpha _1& & \\ \lambda & \beta _1& \alpha _2& \\ & \lambda & \beta _2& \mathrm{}\\ & & \mathrm{}& \mathrm{}\end{array}\right)\left(\begin{array}{c}c_0\\ c_1\\ c_2\\ \mathrm{}\end{array}\right)=0.$$
(5)
In general, unfortunately, equation (5) cannot be treated as an infinite-dimensional limit of its truncated matrix subsystems due to their non-variational, power-series origin . Benefits brought by the tridiagonality concern only the wave functions. Their coefficients are defined by the closed determinantal formulae
$$c_{n+1}det\left(\begin{array}{cccc}\beta _0& \alpha _1& & \\ \lambda & \beta _1& \mathrm{}& \\ & \mathrm{}& \mathrm{}& \lambda \alpha _n\\ & & \lambda & \beta _n\end{array}\right),n=0,1,\mathrm{}.$$
(6)
Once we postulate, in the spirit of the current perturbation theory,
$$c_n(\lambda )=c_n^{(0)}+\lambda c_n^{(1)}+\lambda ^2c_n^{(2)}+\mathrm{},n=1,2,\mathrm{},$$
(7)
we may generate the separate coefficients from eq. (6) by elementary algebra. This type of dependence of the wave functions upon the energy which is not known resembles the Brillouin-Wigner perturbation method able to treat the energy as an external parameter. In the similar vein the use of the variational energies has been recommended in the semi-analytic power-series context, e.g., by Hautot and Tater . The subtle problem of the determination of energies in the present perturbative example is to be settled in what follows.
For the sake of definitness, perturbations with $`c_n(\lambda )=c_n^{(0)}+\lambda h_n(\lambda )`$ will be normalized to $`h_0=0`$. In order to reduce the complexity of formulae we shall only pay attention to the ground state problem with quantum numbers $`p=0`$ and $`N=0`$. This implies that for the sufficiently small perturbations our bound state exists at any non-negative $`\mu >0`$ . The details of transition to $`p=1`$ and/or to $`N>0`$ are not too interesting in the present context. For our choice of $`p=0`$ and $`N=0`$ the accepted normalizations simplify the very first recurrence relation. It reads $`\beta _0+\lambda \alpha _1h_1=0`$, defines the coefficient $`h_1=𝒪(1)`$ “to all orders” and suggests a change of the notation, $`\beta _0=\lambda \gamma _0`$. As long as $`\gamma _0=𝒪(1)`$ we may recall our explicit $`\beta _0\epsilon (\epsilon +\mu +1/2)`$ and infer that $`\epsilon =\epsilon (\lambda )=𝒪(\lambda )`$. Using a particularly convenient “strength-reparametrization” $`a=(2\mu +1)/4`$ we shall postulate that $`\epsilon (\lambda )=\tau (\lambda )a`$ with some not yet known analytic function. In terms of its expansions, say,
$$\tau (\lambda )=a+b\lambda +c\lambda ^2+d\lambda ^3+f\lambda ^4+g\lambda ^5+O\left(\lambda ^6\right)$$
(8)
our new “energy-parameters” $`a,b,\mathrm{}`$ enter the matrix elements $`\alpha _n=n/2n^22n\tau `$, $`\beta _0=\tau ^2a^2`$ and $`\beta _n=n^2+2n\tau +\beta _0`$. For illustration we may insert our ansatz in $`\beta _0=2ab\lambda +\left(2ac+b^2\right)\lambda ^2+O\left(\lambda ^3\right)`$ giving $`\gamma _0=2ab+\left(2ac+b^2\right)\lambda +O\left(\lambda ^2\right)`$ etc.
Wave functions have already been specified by eq. (6). Apparently, there are no conditions imposed upon the energies. This is a paradox which we are going to explain now. In the first step we insert the coefficient $`h_1=\gamma _0/\alpha _1`$ in the first, second and third row of recurrences (4) or (5). The former relation becomes an identity but the next one preserves a genuine three-component character,
$$2(2\epsilon +\mu +2)h_2=1+(\epsilon +1)(\epsilon +\mu +3/2)h_1.$$
(9)
Up to a tiny perturbation the infinite rest with $`n=1,2,\mathrm{}`$ has just a two-term form
$$(n+2)(2\epsilon +n+\mu +2)h_{n+2}=(\epsilon +n+1)(\epsilon +n+\mu +3/2)h_{n+1}+\lambda h_n$$
(10)
giving immediately the asymptotic estimate valid for all $`|\lambda |1`$,
$$h_n=h_n(\lambda )\frac{\mathrm{\Gamma }(\epsilon +n)\mathrm{\Gamma }(\epsilon +n+\mu +1/2)}{n!\mathrm{\Gamma }(2\epsilon +n+\mu +1)}n^{3/2},n1.$$
(11)
This implies the convergence of $`\psi _{(PPT)}(x)`$, up to the central $`x=0`$, rigorously.
The boundary of the circle of convergence coincides with the centre of the spatial symmetry. Hence, for our particular even-parity choice of $`p=0`$ the first derivative of the wave function must have a nodal zero there,
$$_x\psi _{(PPT)}(x)\mathrm{sinh}x\underset{n=0}{\overset{\mathrm{}}{}}\frac{(2n+\kappa )h_n}{\mathrm{cosh}^{2n+\kappa +1}x}0,x0.$$
(12)
As long as $`nh_nn^{1/2}`$ the infinite sum itself is divergent in the origin. The whole expression (12) exhibits an $`0\times \mathrm{}`$ indeterminacy there. The conservation of parity must be imposed “by brute force” . Vice versa, the strict validity of the boundary condition $`_x\psi _{(PPT)}(0)=0`$ represents precisely the “seemingly lost” quantization condition.
In the nearest vicinity of $`x=0`$ the finite but very large value of the sum in eq. (12) is equal to a positive real number multiplied by the coefficient $`h_2`$. This observation follows from eq. (10) which gives, step-by-step, $`h_3(0)=2(\mu +5/2)/(3\mu +9)h_2(0)`$, $`h_4(0)=3(\mu +7/2)/(4\mu +16)h_3(0)`$ etc. In the interval $`\epsilon >(\mu +3)/2`$ all the multiplication factors remain positive. This guarantees that the parity is conserved, in a semi-infinite interval of $`\lambda `$, if and only if the coefficient $`h_2(\lambda )`$ vanishes in the leading order approximation.
This is one of our most important observations. The required leading-order change of sign of $`h_2(\lambda )`$ (i.e., relation $`h_2^{(1)}=0`$) may be interpreted as the consequence of the Sturm-Liouville oscillation theorems . In their light, for an increasing or decreasing leading-order energy $`E=\kappa ^2=[2a1/2+𝒪(\lambda )]^2`$ the nodes of $`_x\psi _{(PPT)}(x)`$ would smoothly move along the real axis of $`x`$. In this sense the condition $`h_2^{(1)}=0`$ of coincidence of one of the nodes with the origin reflects the zero-order physical interpretation of the first coefficient $`a`$ in our expansion (8) and forces us to postulate, consequently, the disappearance of all the subsequent leading-order coefficients,
$$h_2^{(1)}=h_3^{(1)}=\mathrm{}=0.$$
(13)
Before writing it down as an explicit algebra let us first move up to the next order approximation. Due to the danger of a re-introduction of the asymmetry (or, in the other words, of a spike-shaped discontinuity in $`\psi _{(PPT)}(x)`$ at $`x=0`$ ) in the higher orders of $`\lambda `$ we just have to repeat the previous argumentation. In the second order we must first notice that the vanishing condition (13) changed also our recurrences (5). The role of the “last genuine three-term” relation (without a guarantee of a sign-preservation) moves from $`n=1`$ to $`n=2`$ in eq. (4). As long as our estimate (11) remains valid within its own error bound $`1+𝒪(1/n)`$ the role of an overall sign-determining factor $`\lambda h_2(\lambda )=𝒪(\lambda ^2)`$ is taken over by the new norm $`\lambda h_3(\lambda ))=𝒪(\lambda ^2)`$. Mutatis mutandis we get the second-order condition $`h_3^{(2)}=h_4^{(2)}=\mathrm{}=0`$ and, in general,
$$c_n(\lambda )=\lambda h_n(\lambda )=\lambda ^nh_n^{(n)}+\lambda ^{n+1}h_n^{(n+1)}+\mathrm{},n=1,2,\mathrm{}.$$
(14)
This is a perturbative generalization of the standard termination rules. It modifies our above naive expectations and incorporates correctly the physical definite-parity requirement in our original ansatz (3). In a more compact notation let us put $`c_0=f_0(\lambda )=1`$ and $`\lambda h_n(\lambda )\lambda ^nf_n(\lambda )`$ with $`f_n(\lambda )=𝒪(\lambda ^0)`$. This modifies our system of equations (5),
$$\left(\begin{array}{ccccc}\beta _0& \lambda \alpha _1& & & \\ \lambda & \lambda \beta _1& \lambda ^2\alpha _2& & \\ & \lambda ^2& \lambda ^2\beta _2& \lambda ^3\alpha _3& \\ & & \lambda ^3& \lambda ^3\beta _3& \mathrm{}\\ & & & \mathrm{}& \mathrm{}\end{array}\right)\left(\begin{array}{c}f_0\\ f_1\\ f_2\\ f_3\\ \mathrm{}\end{array}\right)=0.$$
(15)
All these equations may be solved by the closed determinantal formula again,
$$f_{n+1}(\lambda )=\frac{\mathrm{\Gamma }(\mu +1+2\lambda \eta )}{\mathrm{\Gamma }(\mu +n+2+2\lambda \eta )}\left[\lambda ^{(n+1)(n+2)/2}_n(\lambda )\right],$$
(16)
$$_n(\lambda )=det\left(\begin{array}{ccccc}\beta _0(z)& \lambda \alpha _1(z)& & & \\ \lambda & \lambda \beta _1(z)& \lambda ^2\alpha _2(z)& & \\ & \lambda ^2& \mathrm{}& \mathrm{}& \\ & & \mathrm{}& \lambda ^{n1}\beta _{n1}& \lambda ^n\alpha _n\\ & & & \lambda ^n& \lambda ^n\beta _n\end{array}\right).$$
It pre-factorizes the wave function coefficients in a certain less usual way. We may re-shuffle it slightly once more, in the more traditional Padé-approximation spirit, with $`f_0=\gamma _1=1`$ and
$$f_j=\frac{\gamma _{j1}}{(\alpha _1)(\alpha _2)\mathrm{}(\alpha _j)},j=1,2,\mathrm{}.$$
(17)
This representation of the wave function coefficients is related to our final, “optimal” recurrences
$$\left(\begin{array}{ccccc}\beta _0& \lambda & & & \\ \alpha _1& \beta _1& \lambda & & \\ & \alpha _2& \beta _2& \lambda & \\ & & \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right)\left(\begin{array}{c}1\\ \gamma _0\\ \gamma _1\\ \mathrm{}\end{array}\right)=0.$$
(18)
The proof of their equivalence to our previous equations is trivial. One has only to verify that the new notation is not inconsistent since the first row of eq. (18) just reproduces our old leading-order “change of the notation” $`\beta _0=\lambda \gamma _0`$.
We are now prepared to compute the unknown auxiliary variables $`b,c,\mathrm{}`$ and, in effect, the perturbed energies. Firstly we get rid of the “already known” $`\gamma _0`$ in eq. (18). This only means that we omit the first line and abbreviate $`\widehat{A}_2=\alpha _2\gamma _0`$ and $`\widehat{B}_1=\beta _1\gamma _0\alpha _1`$ in the remaining equations
$$\left(\begin{array}{ccccc}\widehat{B}_1& \lambda & & & \\ \widehat{A}_2& \beta _2& \lambda & & \\ & \alpha _3& \beta _3& \lambda & \\ & & \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right)\left(\begin{array}{c}1\\ \gamma _1\\ \gamma _2\\ \mathrm{}\end{array}\right)=0.$$
(19)
Now we may parallel the above study where, retrospectively, $`\widehat{B}_0=\beta _0`$. In the other words we only have to require that $`\widehat{B}_1\lambda \gamma _1`$ vanishes in the unperturbed limit (cf. the first line in eq. (19)). In contrast to its trivial $`j=0`$ predecessor the new $`j=1`$ situation imposes a constraint upon $`\widehat{B}_1=2a+1/2+2ab(1+2a)+𝒪(\lambda )`$. Its zero-order component must be zero. This determines the physical value of the first moment $`b=b(a)=\tau ^{(1)}`$,
$$\tau ^{(1)}=\frac{1+4a}{4a\left(1+2a\right)}.$$
(20)
As expected on variational grounds this result means a positive first order change in our energy $`E=\kappa ^2`$ where $`\kappa =\kappa (\lambda )=[4\tau (\lambda )+2\mu 1]/4`$.
The whole $`jj+1`$ procedure can be, obviously, iterated. We arrive at the equations
$$\left(\begin{array}{ccccc}\widehat{B}_j& \lambda & & & \\ \widehat{A}_{j+1}& \beta _{j+1}& \lambda & & \\ & \alpha _{j+2}& \beta _{j+2}& \lambda & \\ & & \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right)\left(\begin{array}{c}1\\ \gamma _j\\ \gamma _{j+1}\\ \mathrm{}\end{array}\right)=0$$
(21)
with the elements $`\widehat{A}_{j+1}=\alpha _{j+1}\gamma _{j1}`$ and $`\widehat{B}_j=\beta _j\gamma _{j1}\alpha _j\gamma _{j2}`$. For any $`j`$, as we have already shown, the first-line rule $`\widehat{B}_j=\lambda \gamma _j=𝒪(\lambda )`$ determines the value of the moment $`\tau ^{(j)}`$. Beyond the formal choice of $`j=0`$ (confirming the consistency of $`a=\tau ^{(0)}=\mu /2+1/4`$) and of $`j=1`$ (giving $`b=\tau ^{(1)}`$) we must perform symbolic manipulations which become a task for the computer. Without its assistance one easily obtains just the next and quite compact $`j=2`$ formula
$$c=\tau ^{(2)}=1/32\frac{\left(1+4a\right)\left(8a^2+8a^3+3a+1\right)}{\left(1+2a\right)^3a^3\left(1+a\right)}.$$
With the assistance of the MAPLE language the evaluation of the next corrections remains entirely straightforward giving $`d`$ or
$$\tau ^{(3)}=\frac{\left(1+4a\right)𝒟_7}{128a^5\left(1+2a\right)^5\left(1+a\right)^2\left(3+2a\right)}$$
with
$$𝒟_7=3+26a+97a^2+168a^3+196a^4+288a^5+320a^6+128a^7$$
or $`f`$ or rather
$$\tau ^{(4)}=\frac{\left(1+4a\right)𝒟_{12}}{2048a^7\left(1+2a\right)^7\left(1+a\right)^3\left(3+2a\right)^2\left(2+a\right)}$$
with
$$𝒟_{12}=90+1335a+8815a^2+32715a^3+69135a^4+54250a^5$$
$$106568a^6340152a^7378096a^8165184a^9+22272a^{10}+43008a^{11}+10240a^{12}$$
etc. These results inspire the general ansatz
$$\tau ^{(K)}=\frac{2^{M(K)}\left(1+4a\right)𝒟_L}{a^{2K1}\left(1+2a\right)^{2K1}\left(a+1\right)^{K1}\left(a+3/2\right)^{K2}\left(a+2\right)^{K3}\mathrm{}\left(a+K/2\right)}$$
re-confirmed by the explicit evaluation of the next-order correction $`g`$,
$$\tau ^{(5)}=\frac{\left(1+4a\right)𝒟_{18}}{8192a^9\left(1+2a\right)^9\left(1+a\right)^4\left(3+2a\right)^3\left(2+a\right)^2\left(5+2a\right)}$$
with the rather complicated
$$𝒟_{18}=3780+80892a+794817a^2+\mathrm{}+458752a^{18}.$$
This did not disprove the obtrusive hypothesis that $`𝒟_L`$ are polynomials in $`a`$ with integer coefficients and with a growing degree $`L=L(K)=(K+1)(K+2)/23`$.
We may summarize our approach to eq. (1) as a scheme which can be easily generalized. Its use of a suitable non-orthogonal basis may find applications in different settings. In a way paralleling the Lanczos method (with a strong numerical flavor) the idea itself has already been tested on several “less solvable” semi-numerical examples . In contrast to them our present construction seems to exhibit a much closer similarity to the current perturbation constructions of anharmonic oscillators. Even in this comparison, our new example seems to hide several pleasant surprises (e.g., a possible non-zero radius $`\lambda _{max}`$ of convergence) which could make it worth of a further study.
### Acknowledgement
Our thorough clarification of the “subtle problem of the determination of energies” has been encouraged by an anonymous referee. The work was supported by the GA AS CR, grant Nr. A1048004. |
warning/0001/hep-th0001199.html | ar5iv | text | # 𝖣-branes in "AdS"₃×𝑆³×𝑆³×𝑆¹
## 1. Introduction
In we have initiated a study of the possible $`𝖣`$-brane configurations in exact superstring backgrounds based on AdS spaces and characterised by a NS-NS antisymmetric field. Here we will continue this study by considering the case of the $`\text{AdS}_3\times S^3\times S^3\times S^1`$ background. String propagation on this background was studied in . This geometry appears as the throat limit of two differently oriented coincident sets of fivebranes intersecting in one direction, together with a set of infinitely stretched strings .
The paper is organised as follows. In the next section we start with a short summary describing the bosonic background in order to set the notation and exhibit the conformal structure. In Section 3 we consider the boundary state formalism adapted to this particular model; we write down the gluing conditions which preserve conformal invariance and the affine symmetry of the underlying current algebra and solve for them, thus determining two classes of bosonic configurations. In Section 4 we identify the $`𝖣`$-brane configurations that each of the two classes of solutions give rise to. The first class of solutions describe $`𝖣`$-brane configurations which can be thought of as a straightforward generalisation of the D-type configurations in $`\text{AdS}_3\times S^3\times T^4`$ . By contrast, the second class of solutions has no such analogue in the $`\text{AdS}_3\times S^3\times T^4`$ background. Therefore we devote Section 5 to determining the twisted conjugacy classes of $`\mathrm{SU}(2)\times \mathrm{SU}(2)`$ and analysing their geometry and topology. In Section 6 we extend our analysis to the $`N=1`$ supersymmetric case and we find that all the bosonic configurations determined previously can be made into $`N=1`$ configurations. We then analyse, in Section 7, the fraction of supersymmetry preserved by these configurations, and we find that they preserve half of the spacetime supersymmetry. We end, in Section 8, with a summary of results.
## 2. The bosonic $`\text{AdS}_3\times S^3\times S^3\times S^1`$ background
Bosonic string propagation on the background $`\text{AdS}_3\times S^3\times S^3\times S^1`$ can be described using a free compact boson for the $`S^1`$ factor, and a WZW model with semisimple group $`\mathrm{SL}(2,\mathrm{})\times \mathrm{SU}(2)\times \mathrm{SU}(2)`$ for the $`\text{AdS}_2\times S^3\times S^3`$ part.
The theory of the compact free boson for the $`S^1`$ factor is described by the action
$$I_{S^1}[\phi ]=_\mathrm{\Sigma }\phi \overline{}\phi ,$$
where $`\mathrm{\Sigma }`$ is an orientable Riemann surface. The operator product algebra is standard:
$$\phi (z)\phi (w)=\frac{1}{(zw)^2}+\text{reg},$$
(1)
with similar operator product expansions for the antiholomorphic sector.
The WZW model describing the $`\text{AdS}^3\times S^3\times S^3`$ part has as target the semisimple group $`𝐆=𝐆_\mathrm{𝟏}\times 𝐆_\mathrm{𝟐}\times 𝐆_\mathrm{𝟑}`$, with $`𝐆_\mathrm{𝟏}=\mathrm{SL}(2,\mathrm{})`$, and $`𝐆_\mathrm{𝟐}`$ and $`𝐆_\mathrm{𝟑}`$ two copies of $`\mathrm{SU}(2)`$. The corresponding action will therefore be a sum of three terms
$$I=I_{\mathrm{SL}(2,\mathrm{})}[g_1]+I_{\mathrm{SU}(2)}[g_2]+I_{\mathrm{SU}(2)}[g_3],$$
(2)
where each term is of the standard form
$$\frac{2}{k_i}I[g_i]=_\mathrm{\Sigma }g_i^1g_i,g_i^1\overline{}g_i+\frac{1}{6}_Bg_i^1dg_i,[g_i^1dg_i,g_i^1dg_i],$$
with $`B=\mathrm{\Sigma }`$. Each of the fields $`g_i`$ is a map from $`\mathrm{\Sigma }`$ to the Lie group $`𝐆_𝐢`$, $`i=1,2,3`$. We denote by $`𝔤=𝔤_1𝔤_2𝔤_3`$ the corresponding Lie algebra, where $`𝔤_1=𝔰𝔩(2,\mathrm{})`$ and $`𝔤_2,𝔤_3=𝔰𝔲(2)`$. For these algebras we choose the following bases of generators: $`\{T_a\}`$ for $`𝔤_1`$, $`\{X_a\}`$ for $`𝔤_2`$ and $`\{Y_a\}`$ for $`𝔤_3`$ satisfying
$$[T_1,T_2]=T_3,[T_2,T_3]=T_1,[T_3,T_1]=T_2,$$
and
$$[X_1,X_2]=X_3,[X_2,X_3]=X_1,[X_3,X_1]=X_2,$$
and similar Lie brackets for the $`Y_a`$’s. We also need to specify an invariant metric on $`𝔤`$, which has a diagonal form
$$\eta =\left(\begin{array}{ccc}\eta _1& 0& 0\\ 0& \eta _2& 0\\ 0& 0& \eta _3\end{array}\right),$$
with three components given by $`(\eta _1)_{ab}=\mathrm{diag}(+,+,)`$ and $`(\eta _2)_{ab}=(\eta _3)_{ab}=\mathrm{diag}(+,+,+)`$.
A group element $`g=(g_1,g_2,g_3)`$ in $`𝐆`$ can be parametrised as follows:
$$g_1=e^{\theta _2T_2}e^{\theta _1T_1}e^{\theta _3T_3},g_2=e^{\varphi _2X_2}e^{\varphi _1X_1}e^{\varphi _3X_3},$$
(3)
$$g_3=e^{\sigma _2Y_2}e^{\sigma _1Y_1}e^{\sigma _3Y_3},$$
(4)
where $`\theta _\mu `$, $`\varphi _\mu `$ and $`\sigma _\mu `$, $`\mu =1,2,3`$, play the rôle of the spacetime fields. In terms of them (2) becomes a sigma-model action, whose spacetime metric and antisymmetric field can be straightforwardly obtained (the expressions for $`\mathrm{SL}(2,\mathrm{})`$ and $`\mathrm{SU}(2)`$ were explicitly written down in ).
This model possesses, as is well-known, an infinite-dimensional symmetry group $`𝐆(z)\times 𝐆(\overline{z})`$ characterised by the conserved currents $`\mathrm{𝕁}(z)=gg^1`$ and $`\overline{\mathrm{𝕁}}(\overline{z})=g^1\overline{}g`$, which underlies the exact conformal invariance of this background. The conserved currents generate, at the quantum level, an affine Lie algebra, $`\widehat{𝔤}_1\widehat{𝔤}_2\widehat{𝔤}_3`$, described by
$$\mathrm{𝕁}_a(z)\mathrm{𝕁}_b(w)=\frac{h_{ab}}{(zw)^2}+\frac{f_{ab}^{}{}_{}{}^{c}\mathrm{𝕁}_c(w)}{zw}+\text{reg},$$
(5)
where the indices run over the whole Lie algebra, $`a,b=1,\mathrm{},9`$, and the coefficients $`h_{ab}`$ define the symmetric bilinear form
$$h=\left(\begin{array}{ccc}k_1\eta _1& 0& 0\\ 0& k_2\eta _2& 0\\ 0& 0& k_3\eta _3\end{array}\right).$$
(6)
The parameters $`k_i`$ are related to the level $`x_i`$ of the corresponding affine algebra by $`k_i=x_i+g_i^{}`$, where $`g_i^{}`$ is the dual Coxeter number.
The CFT corresponding to this string background is then described by the energy-momentum tensor
$$𝖳=\mathrm{\Omega }^{ab}(\mathrm{𝕁}_a\mathrm{𝕁}_b)+(\phi \phi ),$$
where $`\mathrm{\Omega }^{ab}`$ are the components of the inverse $`\mathrm{\Omega }^1`$ of the invariant metric
$$\mathrm{\Omega }=\left(\begin{array}{ccc}\mathrm{\Omega }_1& 0& 0\\ 0& \mathrm{\Omega }_2& 0\\ 0& 0& \mathrm{\Omega }_3\end{array}\right),$$
with the components given by $`\mathrm{\Omega }_1=2(k_1+1)\eta _1`$, $`\mathrm{\Omega }_2=2(k_21)\eta _2`$, and $`\mathrm{\Omega }_3=2(k_31)\eta _3`$. The central charge of this CFT is given by
$$c=\frac{3k_1}{k_1+1}+\frac{3k_2}{k_21}+\frac{3k_3}{k_31}+1.$$
## 3. Boundary states
The strategy that we will follow in order to determine the $`𝖣`$-brane configurations which can be consistently defined in type IIB string theory on $`\text{AdS}_3\times S^3\times S^3\times S^1`$ is similar to the one used in and further developed in . We consider a class of gluing conditions, which is defined in terms of a Lie algebra automorphism, $`R:𝔤𝔤`$, which preserves the metric $`\eta `$:
$$[R(Z_a),R(Z_b)]=R([Z_a,Z_b]),$$
(7)
$$R^T\eta R=\eta ,$$
(8)
where $`\{Z_a\}`$ is a given basis in $`𝔤`$, in terms of which $`R`$ is given by $`R(Z_a)=Z_bR_{}^{b}{}_{a}{}^{}`$. The gluing conditions read
$$\mathrm{𝕁}_a(z)R_{}^{b}{}_{a}{}^{}\overline{\mathrm{𝕁}}_b(\overline{z})=0,$$
(9)
and can be easily seen to preserve the current algebra of the bulk theory.
These gluing conditions have to satisfy the basic consistency requirement, which is conformal invariance. In the bosonic case, this comes down to imposing
$$𝖳(z)=\overline{𝖳}(\overline{z}),$$
at the boundary. In this case, the requirement of conformal invariance translates into the condition
$$R^T\mathrm{\Omega }R=\mathrm{\Omega }.$$
(10)
As explained in , $`𝖣`$-branes in a WZW model with group $`𝐆`$ are classified by the group $`\mathrm{Out}_o(𝐆)`$ of metric-preserving outer automorphisms of $`𝐆`$, which is defined as the quotient $`\mathrm{Aut}_o(𝐆)/\mathrm{Inn}_o(𝐆)`$ of the group of metric-preserving automorphisms by the invariant subgroup of inner automorphisms. For the case at hand, and ignoring the $`S^1`$ factor for which no automorphism is inner, $`\mathrm{Out}_o(𝐆)\mathrm{}_2`$, whence there are two distinct types of $`𝖣`$-branes on $`\mathrm{SL}(2,\mathrm{})\times \mathrm{SU}(2)\times \mathrm{SU}(2)`$. Let us see this.
Group automorphisms for simply connected groups come by exponentiating automorphisms of the Lie algebra. In our case, the Lie algebra $`𝔤`$ of $`𝐆`$ is a direct sum of three terms $`𝔤_1`$, $`𝔤_2`$ and $`𝔤_3`$, each of them being a three-dimensional simple Lie algebra. Since $`𝔰𝔩(2,\mathrm{})`$ and $`𝔰𝔲(2)`$ are non-isomorphic simple Lie algebras, there is no nontrivial homomorphism between them. We therefore deduce that the matrix of boundary conditions defined by the automorphism $`R:𝔤_1𝔤_2𝔤_3𝔤_1𝔤_2𝔤_3`$, must take one of the following two forms
$$R_I=\left(\begin{array}{ccc}R_{11}& 0& 0\\ 0& R_{22}& 0\\ 0& 0& R_{33}\end{array}\right),R_{II}=\left(\begin{array}{ccc}R_{11}& 0& 0\\ 0& 0& R_{23}\\ 0& R_{32}& 0\end{array}\right),$$
(11)
where $`R_{ij}:𝔤_j𝔤_i`$, for any $`i,j=1,2,3`$. We thus have two classes of solutions: the first, described by $`R_I`$, exists for any values of the parameters $`k_i`$, whereas the second, given by $`R_{II}`$, exists only for particular values of the parameters, such that $`k_2=k_3`$. Moreover, from (10) and (8), it follows that $`R_{11}`$ belongs to $`\mathrm{O}(2,1)`$, and $`R_{ij}`$, with $`i,j=2,3`$, belong to $`\mathrm{O}(3)`$. On the other hand, from (7) we deduce that $`R_{ii}`$ are Lie algebra automorphisms, corresponding to $`𝔰𝔩(2,\mathrm{})`$ and $`𝔰𝔲(2)`$, whereas $`R_{23}`$ and $`R_{32}`$ are Lie algebra isomorphisms. Explicitly, each of these conditions translates into a condition on the corresponding matrix, that reads
$$det(R_{ij})=1,$$
which makes $`R_{11}`$ belong to $`\mathrm{SO}(2,1)`$, and $`R_{ij}`$, for $`i,j=2,3`$, to $`\mathrm{SO}(3)`$.
Clearly, the main difference between these two classes of solutions is that $`R_I`$ describe inner automorphisms, whereas $`R_{II}`$ does not. More precisely, we have
$$R_{II}=TR_I,$$
where the matrix $`T`$ is given by
$$T=\left(\begin{array}{ccc}\mathbb{𝟙}& 0& 0\\ 0& 0& \mathbb{𝟙}\\ 0& \mathbb{𝟙}& 0\end{array}\right),$$
with obvious notation.
These results can be summarised as follows. We consider a set of gluing conditions on the group manifold $`\mathrm{SL}(2,\mathrm{})\times \mathrm{SU}(2)\times \mathrm{SU}(2)`$ which preserve conformal invariance and the infinite-dimensional symmetry of the current algebra of the bulk theory. These gluing conditions are described in terms of metric-preserving automorphisms of the Lie algebra $`𝔰𝔩(2,\mathrm{})𝔰𝔲(2)𝔰𝔲(2)`$. They admit two classes of solutions, characterised by different matrices of gluing conditions: for generic values of the levels $`k_i`$, the solutions are parametrised by elements of $`\mathrm{SO}(2,1)\times \mathrm{SO}(3)\times \mathrm{SO}(3)`$; additionally, for the particular values of $`k_i`$ such that $`k_2=k_3`$ we have an extra set of solutions, parametrised again by the elements of the group $`\mathrm{SO}(2,1)\times \mathrm{SO}(3)\times \mathrm{SO}(3)`$.
## 4. $`𝖣`$-brane solutions
Let us begin with the $`𝖣`$-brane configurations produced by the inner automorphisms $`R_I`$. This case represents a straightforward generalisation of the corresponding analysis performed for the $`\text{AdS}_3\times S^3\times T^4`$ background. Indeed, since the matrix of gluing conditions $`R_I`$ is block diagonal, the resulting $`𝖣`$-brane configurations take a product form, $`𝒟_{\mathrm{SL}(2,\mathrm{})}\times 𝒟_{\mathrm{SU}(2)}\times 𝒟_{\mathrm{SU}(2)}`$, where $`𝒟_𝐆`$ represent the $`𝖣`$-brane configurations in the group manifold $`𝐆`$. The possible $`𝖣`$-brane configurations in $`\mathrm{SL}(2,\mathrm{})`$ and $`\mathrm{SU}(2)`$ have been studied in detail in . We therefore obtain in our case that the $`𝖣`$-brane solution passing through a point $`g`$ in $`𝐆`$ and being described by a set of gluing conditions defined by the inner automorphism $`R_I`$ is characterised by a worldvolume which lies along a product of conjugacy classes shifted by group elements determined by $`R_I`$:
$$𝒞_{\mathrm{SL}(2,\mathrm{})}(g_1r_1^1)r_1\times 𝒞_{\mathrm{SU}(2)}(g_2r_2^1)r_2\times 𝒞_{\mathrm{SU}(2)}(g_3r_3^1)r_3,$$
where $`R_{11}=\mathrm{Ad}_{r_1}`$, $`R_{22}=\mathrm{Ad}_{r_2}`$, $`R_{33}=\mathrm{Ad}_{r_3}`$, for some $`(r_1,r_2,r_3)`$ in $`𝐆`$. The conjugacy classes of $`\mathrm{SU}(2)`$ are parametrised by $`S^1/\mathrm{}_2`$, which we can understand as the interval $`\theta [0,\pi ]`$. The conjugacy classes corresponding to $`\theta =0,\pi `$ are points, corresponding to the elements $`\pm e`$ in the centre of $`\mathrm{SU}(2)`$, whereas the classes corresponding to $`\theta (0,\pi )`$ are 2-spheres. If we picture $`\mathrm{SU}(2)`$, which is homeomorphic to the 3-sphere, as the one-point compactification of $`\mathrm{}^3`$ where the sphere at infinity is collapsed to a point, the foliation of $`\mathrm{SU}(2)`$ by its conjugacy classes coincides with the standard foliation of $`\mathrm{}^3`$ by 2-spheres with two degenerate spheres at the origin and at infinity. For $`\mathrm{SL}(2,\mathrm{})`$, on the other hand, we have three types of metrically nondegenerate<sup>1</sup><sup>1</sup>1This condition is necessary for their interpretation as D-branes. conjugacy classes (for details, see ): two point-like ones, corresponding to the two elements in the centre of $`\mathrm{SL}(2,\mathrm{})`$, a family of two-dimensional classes with planar topology and a family of two-dimensional classes with cylindrical topology.
We now turn to the $`𝖣`$-branes produced by outer automorphisms $`R_{II}`$. Notice first of all that, since the matrix of gluing conditions $`R_{II}`$ has a block diagonal form, the resulting $`𝖣`$-brane configurations take, also in this case, a product form, $`𝒟_{\mathrm{SL}(2,\mathrm{})}\times 𝒟_{\mathrm{SU}(2)\times \mathrm{SU}(2)}`$, where $`𝒟_{\mathrm{SL}(2,\mathrm{})}`$ represent the same $`𝖣`$-brane configurations in the group manifold $`\mathrm{SL}(2,\mathrm{})`$ which were discussed above. It remains to analyse the $`𝒟`$-brane configurations on the product group $`\mathrm{SU}(2)\times \mathrm{SU}(2)`$, corresponding to the gluing conditions
$$\stackrel{~}{g}\stackrel{~}{g}^1=\stackrel{~}{R}_{II}(\stackrel{~}{g}^1\overline{}\stackrel{~}{g}),$$
(12)
where $`\stackrel{~}{g}=(g_2,g_3)`$, and the corresponding matrix of gluing conditions
$$\stackrel{~}{R}_{II}=\left(\begin{array}{cc}0& R_{23}\\ R_{32}& 0\end{array}\right)$$
can be easily shown to be of the form
$$\stackrel{~}{R}_{II}=\stackrel{~}{T}\mathrm{Ad}_{\stackrel{~}{r}},\stackrel{~}{T}=\left(\begin{array}{cc}0& \mathbb{𝟙}\\ \mathbb{𝟙}& 0\end{array}\right),$$
where $`\stackrel{~}{r}=(r_2,r_3)`$, while $`R_{32}=\mathrm{Ad}_{r_2}`$ and $`R_{23}=\mathrm{Ad}_{r_3}`$. This implies that it is sufficient to analyse the $`𝖣`$-brane configurations produced by $`\stackrel{~}{R}_{II}=\stackrel{~}{T}`$ in $`\mathrm{SU}(2)\times \mathrm{SU}(2)`$; any other $`\stackrel{~}{R}_{II}`$ will lead to configurations which differ from these only by translations in the group manifold. Indeed, a set of gluing conditions described by a generic $`\stackrel{~}{R}_{II}`$
$$\stackrel{~}{g}\stackrel{~}{g}^1=\stackrel{~}{T}\mathrm{Ad}_{\stackrel{~}{r}}(\stackrel{~}{g}^1\overline{}\stackrel{~}{g}),$$
for some $`\stackrel{~}{r}`$ in $`\mathrm{SU}(2)\times \mathrm{SU}(2)`$, can be written as
$$hh^1=\stackrel{~}{T}(h^1\overline{}h),$$
with $`h=\stackrel{~}{g}\stackrel{~}{r}^1`$. According to the general theory developed in (for a somewhat different approach, see ) the $`𝖣`$-brane configurations produced by the gluing conditions (12) with $`\stackrel{~}{R}_{II}=\stackrel{~}{T}`$ are nothing but the twisted conjugacy classes determined by the outer automorphism $`\stackrel{~}{T}`$, which we denote by $`𝒞_{\mathrm{SU}(2)\times \mathrm{SU}(2)}^{\stackrel{~}{T}}(\stackrel{~}{g})`$. In summary, the $`𝖣`$-brane configurations passing through a point $`g`$ in $`\mathrm{SL}(2,\mathrm{})\times \mathrm{SU}(2)\times \mathrm{SU}(2)`$ and described by the matrix of gluing conditions $`R_{II}`$ have worldvolumes which lie along a product of twisted conjugacy classes
$$𝒞_{\mathrm{SL}(2,\mathrm{})}(g_1r_1^1)r_1\times 𝒞_{\mathrm{SU}(2)\times \mathrm{SU}(2)}^{\stackrel{~}{T}}(\stackrel{~}{g}\stackrel{~}{r}^1)\stackrel{~}{r}.$$
Thus, in order to complete our analysis, we have to determine the twisted conjugacy classes of $`\mathrm{SU}(2)\times \mathrm{SU}(2)`$. Next section will be devoted to solving this mathematical problem. There we will show that the twisted conjugacy classes of $`\mathrm{SU}(2)\times \mathrm{SU}(2)`$ corresponding to the outer automorphism $`\stackrel{~}{T}`$ are basically determined by the (ordinary) conjugacy classes of $`\mathrm{SU}(2)`$. Indeed, if we denote by $`m`$ the group multiplication, $`m:\mathrm{SU}(2)\times \mathrm{SU}(2)\mathrm{SU}(2)`$, which assigns to every element $`(g_2,g_3)`$ in $`\mathrm{SU}(2)\times \mathrm{SU}(2)`$ an element in $`\mathrm{SU}(2)`$ given by $`m(g_2,g_3)=g_2g_3`$, then we have that
$$𝒞_{\mathrm{SU}(2)\times \mathrm{SU}(2)}^{\stackrel{~}{T}}(g_2,g_3)=m^1𝒞_{\mathrm{SU}(2)}(g_2g_3).$$
Using this result we will see that the group $`\mathrm{SU}(2)\times \mathrm{SU}(2)`$ has two types of conjugacy classes: two three-dimensional ones, diffeomorphic to $`S^3`$, and a family of five-dimensional classes, diffeomorphic to $`S^2\times S^3`$.
## 5. Some twisted conjugacy classes in $`\mathrm{SU}(2)\times \mathrm{SU}(2)`$
The primary goal of this section is to determine the twisted conjugacy classes of $`\mathrm{SU}(2)\times \mathrm{SU}(2)`$ and understand their geometry and topology. However we think it might be useful to consider in the beginning a slightly more general problem, which basically consists in replacing $`\mathrm{SU}(2)`$ with an arbitrary group.
Let $`G`$ be a Lie group and let $`D=G\times G`$ be the product group. Let $`\tau :DD`$ be the twist $`\tau (x,y)=(y,x)`$. It is clearly an outer automorphism of $`D`$.
Let $`\mathrm{Ad}_\tau `$ denote the twisted adjoint action of $`D`$ on $`D`$,
$$\mathrm{Ad}_\tau (x,y)(x_0,y_0)=(xx_0y^1,yy_0x^1),$$
for all $`(x,y),(x_0,y_0)D`$. By the twisted conjugacy class of a point $`(x,y)D`$, we mean the orbit of $`(x,y)`$ under the twisted adjoint action of $`D`$ on $`D`$. We denote it $`𝒪_{(x,y)}`$. In other words,
$$𝒪_{(x,y)}=\{(wxz^1,zyw^1)w,zG\}.$$
In the notation of the last section, $`𝒪_{(x,y)}=𝒞_{G\times G}^{\stackrel{~}{T}}(x,y)`$.
We would like to determine the twisted conjugacy classes of $`D`$ and explore their topology. If, as in the case of interest, $`G`$ (and hence $`D`$) possesses a bi-invariant metric, we also would like to say something about the geometry of the twisted conjugacy classes as submanifolds of $`D`$.
We start with an example. The twisted conjugacy class of the identity $`(e,e)`$ is the “anti-diagonal”, a submanifold of $`D`$ diffeomorphic to $`G`$, but which is not a subgroup:
$$\mathrm{Ad}_\tau (x,y)(e,e)=(xy^1,yx^1)=(z,z^1)$$
for $`z:=xy^1`$. Hence the orbit $`𝒪_{(e,e)}`$ of $`(e,e)`$ is
$$𝒪_{(e,e)}=\{(z,z^1)|zG\}G,$$
where the isomorphism is one of differentiable manifolds.
This example points the way to determining the rest of the twisted conjugacy classes. We start with a series of observations.
Group multiplication gives a natural surjection $`m:D=G\times GG`$. We will write it simply as $`m(x,y)=xy`$. It is easy to see that the inverse image of a point is diffeomorphic to $`G`$. Indeed, suppose $`xy=x^{}y^{}`$. Then $`x^1x^{}=y(y^{})^1=z`$, say. Therefore, $`x^{}=xz`$ and $`y^{}=z^1y`$, for some $`zG`$. In other words,
$$m^1(xy)=\{(xz,z^1y)zG\}G.$$
This means that $`G\times G`$ is a bundle over $`G`$ with fibre $`G`$, but the fibration is different than the standard fibrations $`\mathrm{pr}_1:G\times GG`$ sending $`(x,y)x`$ and $`\mathrm{pr}_2:G\times GG`$ sending $`(x,y)y`$. Nevertheless $`m:G\times GG`$ is a principal $`G`$-bundle. To prove this we must simply exhibit a free action of $`G`$ which preserves the fibres. The typical fibre is given by
$$m^1(xy)=\{(xz,z^1y)zG\}.$$
Let $`gG`$ and consider the action $`(x,y)(xg^1,gy)`$. This is clearly free and moreover $`xg^1gy=xy`$, whence it preserves the fibre.
We can now state the main result of this section.
###### Theorem 1.
Let $`𝒪_{(x,y)}`$ be the twisted conjugacy class of $`(x,y)D`$ and let $`𝒞_z`$ be the (standard) conjugacy class of $`zG`$. Then,
$$𝒪_{(x,y)}=m^1𝒞_{xy}.$$
In other words, twisted conjugacy classes in $`D`$ are the inverse images under the group multiplication of the (standard) conjugacy classes in $`G`$.
###### Proof.
A typical element in $`𝒪_{(x,y)}`$ is $`(uxv^1,vyu^1)`$. The product of these two elements is $`uxyu^1`$, which is conjugate to $`xyG`$. In other words, $`m(𝒪_{(x,y)})=𝒞_{xy}`$, whence $`𝒪_{(x,y)}m^1𝒞_{xy}`$. To prove the reverse inclusion, we will prove that if $`ab`$ and $`cd`$ are conjugate in $`G`$, then $`(a,b)`$ and $`(c,d)`$ belong to the same twisted conjugacy class in $`D`$. If $`ab`$ and $`cd`$ are conjugate, then $`cd=zabz^1`$ for some $`zG`$. This is equivalent to
$`z^1cd=abz^1`$ $`a^1z^1c=bz^1d^1=w^1wG`$
$`c=zaw^1\text{and}d=wbz^1`$
$`(c,d)=(zaw^1,wbz^1),`$
whence $`(c,d)`$ and $`(a,b)`$ are in the same twisted conjugacy class. ∎
As a corollary we have that twisted conjugacy classes in $`D=G\times G`$ are principal $`G`$-bundles over conjugacy classes of $`G`$.
We now specialise to $`G=\mathrm{SU}(2)`$. The Lie group $`\mathrm{SU}(2)`$ has two kinds of conjugacy classes: points (at $`\pm e`$) and $`2`$-spheres everywhere else. Now, the inverse image under $`m`$ of the points are topologically 3-spheres, whereas the inverse image of a $`2`$-sphere is a principal $`\mathrm{SU}(2)`$-bundle over $`S^2`$. Since principal $`G`$-bundles over $`S^2`$ are classified up to homotopy by $`\pi _1(G)`$, the fact that $`\mathrm{SU}(2)`$ is simply-connected implies that the bundle is trivial. In other words, the Lie group $`\mathrm{SU}(2)\times \mathrm{SU}(2)`$ has two types of twisted conjugacy classes, diffeomorphic to $`S^3`$ or to $`S^2\times S^3`$.
Notice furthermore that the $`S^3`$ orbits are homologically nontrivial. This is because the maps induced in homology by the canonical projections $`\mathrm{pr}_1`$ and $`\mathrm{pr}_2`$ send the homology classes of the orbits to the fundamental class of $`\mathrm{SU}(2)`$, up to orientation.
For example, if $`𝒪=𝒪_{(e,e)}`$, notice that
$$\mathrm{pr}_1(z,z^1)=z\text{and}\mathrm{pr}_2(z,z^1)=z^1.$$
Therefore $`(\mathrm{pr}_1)_{}[𝒪]=1\mathrm{}H_3(S^3)`$ and $`(\mathrm{pr}_2)_{}[𝒪]=1\mathrm{}H_3(S^3)`$.
On the other hand, the five-dimensional classes are homologically trivial, since $`H_5(S^3\times S^3)=0`$. Nevertheless a similar argument shows that they are not homotopically trivial.
It remains to see how the twisted conjugacy classes are embedded geometrically relative to the bi-invariant metric on $`G\times G`$. We will first consider the small orbits (of dimension $`dimG`$) obtained from point-like conjugacy classes in $`G`$. As we will show, they are totally geodesic submanifold of $`D`$; and, in particular, they are minimal.
Consider $`𝒪=𝒪_{(e,e)}`$. It is not a subgroup of $`D`$, but we will see that it is a subgroup of $`D^{}`$, a Lie group which shares the same underlying manifold as $`D`$, but whose group multiplication is different. Moreover, the bi-invariant metric on $`D`$ is also bi-invariant in $`D^{}`$. It will follow that $`𝒪`$ is totally geodesic as a consequence of the following well-known result (see, for example, Exercise 6.6 in ):
###### Theorem 2.
Let $`G`$ be a Lie group with a bi-invariant metric. Then any subgroup $`H`$ is a totally geodesic submanifold.
The group $`D^{}`$ is defined as $`G\times G^{\mathrm{opp}}`$, where $`G^{\mathrm{opp}}`$, the opposite group, is the group sharing the same underlying manifold with $`G`$ but with the opposite multiplication law:
$`m^{\mathrm{opp}}:G^{\mathrm{opp}}\times G^{\mathrm{opp}}`$ $`G^{\mathrm{opp}}`$
$`(x,y)`$ $`yx.`$
In other words, $`m^{\mathrm{opp}}=m\tau `$. Clearly $`G^{\mathrm{opp}}`$ and $`G`$ are isomorphic as Lie groups: the isomorphism $`GG^{\mathrm{opp}}`$ being defined by $`xx^1`$.
Under the group multiplication in $`D^{}=G\times G^{\mathrm{opp}}`$,
$`m^{}:D^{}\times D^{}`$ $`D^{}`$
$`((x,y),(u,v))`$ $`(xu,vy),`$
it is clear that $`𝒪`$ is a subgroup:
$$(x,x^1)(y,y^1)=(xy,y^1x^1)=(xy,(xy)^1).$$
Since $`D^{}`$ has the same underlying manifold as $`D`$, we have a $`D`$-bi-invariant metric on it. To be able to apply the theorem, it remains to show that this metric is also $`D^{}`$-bi-invariant. The bi-invariant metric on $`D=G\times G`$ is the riemannian product of the *same* bi-invariant metric on each of the factors: this guarantees that $`\tau `$ is an isometry. Therefore this metric on $`G\times G`$ will be bi-invariant under $`D^{}`$ if and only if the metric on $`G`$ is $`G^{\mathrm{opp}}`$-bi-invariant. But this follows trivially from the following observation. For $`xG`$ let $`L(x)`$ and $`R(x)`$ denote the left- and right-multiplication by $`x`$ in $`G`$, respectively. Similarly, let $`L^{\mathrm{opp}}(x)`$ and $`R^{\mathrm{opp}}(x)`$ denote the similar operations in $`G^{\mathrm{opp}}`$. Then one has
$$L(x)=R^{\mathrm{opp}}(x)\text{and}R(x)=L^{\mathrm{opp}}(x).$$
This means that left-invariance under $`G`$ is equivalent to right-invariance under $`G^{\mathrm{opp}}`$ and viceversa. In particular, bi-invariance under $`G`$ is equivalent to bi-invariance under $`G^{\mathrm{opp}}`$. We conclude that since the metric on $`G`$ is $`G`$-bi-invariant, it is also $`G^{\mathrm{opp}}`$-bi-invariant. In summary, we have just proven the following:
###### Theorem 3.
The twisted conjugacy class $`𝒪_{(e,e)}`$ is totally-geodesic relative to the bi-invariant metric on $`D`$. In particular, it is minimal.
How about the other small twisted conjugacy classes? These are the inverse images by the multiplication $`m`$ of point-like conjugacy classes in $`G`$, hence of elements in the centre of $`G`$. Let $`z`$ be an element in the centre of $`G`$. Then the twisted conjugacy class $`𝒪_{(e,z)}:=m^1(z)`$ is given by
$$𝒪_{(e,z)}=\{(x,x^1z)xG\}.$$
This is the translate (both left and right) of $`𝒪_{(e,e)}`$ by the element $`(e,z)`$. Since the metric on $`D`$ is bi-invariant, $`𝒪_{(e,z)}`$ is isometric to $`𝒪_{(e,e)}`$ as submanifolds of $`D`$. In particular, since $`𝒪_{(e,e)}`$ is totally geodesic, so is $`𝒪_{(e,z)}`$. This proves the following:
###### Theorem 4.
Let $`zG`$ be any element in the centre. The twisted conjugacy class $`𝒪_{(e,z)}`$ is totally-geodesic relative to the bi-invariant metric on $`D`$. In particular, it is minimal.
Specialising to $`G=\mathrm{SU}(2)`$ we have that the twisted conjugacy classes $`𝒪_{(e,e)}`$ and $`𝒪_{(e,e)}`$ in $`\mathrm{SU}(2)\times \mathrm{SU}(2)`$ are totally-geodesic three-spheres.
How about the larger orbits? In this case, it is possible to argue that it is metrically a fibre product with totally-geodesic fibres diffeomorphic to $`G`$. But the total space of the bundle is certainly not totally geodesic. In fact, it is in general not even minimal.
## 6. The $`N=1`$ supersymmetric extension
Let us now consider the $`N=1`$ supersymmetric extension of the affine Lie algebra $`\widehat{𝔤}`$, which we will denote by $`\widehat{𝔤}_{N=1}=\widehat{𝔰𝔩}(2,\mathrm{})_{N=1}\widehat{𝔰𝔲}(2)_{N=1}`$, with generators $`(\mathrm{𝕁}_a,\mathrm{\Psi }_a)`$ satisfying
$`\mathrm{𝕁}_a(z)\mathrm{𝕁}_b(w)`$ $`={\displaystyle \frac{h_{ab}}{(zw)^2}}+{\displaystyle \frac{f_{ab}^{}{}_{}{}^{c}\mathrm{𝕁}_c(w)}{zw}}+\text{reg},`$ (13)
$`\mathrm{𝕁}_a(z)\mathrm{\Psi }_b(w)`$ $`={\displaystyle \frac{f_{ab}^{}{}_{}{}^{c}\mathrm{\Psi }_c(w)}{zw}}+\text{reg},`$ (14)
$`\mathrm{\Psi }_a(z)\mathrm{\Psi }_b(w)`$ $`={\displaystyle \frac{h_{ab}}{zw}}+\text{reg},`$ (15)
with $`h_{ab}`$ defined as in (6). The free fields $`(\phi ,\lambda )`$ on $`S^1`$ satisfy the standard OPEs
$`\phi (z)\phi (w)`$ $`={\displaystyle \frac{1}{(zw)^2}}+\text{reg},`$ (16)
$`\lambda (z)\lambda (w)`$ $`={\displaystyle \frac{1}{zw}}+\text{reg}.`$ (17)
Then the generators of the $`N=1`$ SCA will be given by
$`𝖳(z)`$ $`=\frac{1}{2}h^{ab}(\stackrel{~}{\mathrm{𝕁}}_a\stackrel{~}{\mathrm{𝕁}}_b)+\frac{1}{2}h^{ab}(\mathrm{\Psi }_a\mathrm{\Psi }_b)+\frac{1}{2}(\phi \phi )+\frac{1}{2}(\lambda \lambda )`$
$`𝖦(z)`$ $`=h^{ab}(\stackrel{~}{\mathrm{𝕁}}_a\mathrm{\Psi }_b)+(\phi \lambda )\frac{1}{6k^2}f^{abc}(\mathrm{\Psi }_a\mathrm{\Psi }_b\mathrm{\Psi }_c),`$
where we have introduced the so-called decoupled currents, $`\stackrel{~}{\mathrm{𝕁}}_a\mathrm{𝕁}_a\frac{1}{2}h^{bd}f_{ab}^{}{}_{}{}^{c}(\mathrm{\Psi }_c\mathrm{\Psi }_d)`$, in terms of which the superconformal generators take a relatively simple form. The coefficients $`h^{ab}`$ are the components of $`h^1`$. The central charge of this SCFT is given by
$$c=\frac{3}{2}+\frac{3(k_11)}{k_1}+\frac{3}{2}+\frac{3(k_2+1)}{k_2}+\frac{3}{2}+\frac{3(k_3+1)}{k_3}+\frac{3}{2}.$$
In order to have a critical superstring theory the levels must satisfy the following relation
$$\frac{1}{k_1}\frac{1}{k_2}\frac{1}{k_3}=0,$$
in which case $`c=15`$. And, since we have a similar structure for the antiholomorphic sector as well, we actually have a $`(1,1)`$ SCFT.
The gluing conditions are given by
$$\mathrm{𝕁}_a(z)R_{}^{b}{}_{a}{}^{}\overline{\mathrm{𝕁}}_b(\overline{z})=0,\mathrm{\Psi }_a(z)S_{}^{b}{}_{a}{}^{}\overline{\mathrm{\Psi }}_b(\overline{z})=0,$$
(18)
where the coefficients $`R_{}^{b}{}_{a}{}^{}`$ and $`S_{}^{b}{}_{a}{}^{}`$ are defined by $`R,S:𝔤𝔤`$, with $`R(Z_a)=Z_bR_{}^{b}{}_{a}{}^{}`$ and $`S(Z_a)=Z_bS_{}^{b}{}_{a}{}^{}`$, for any $`Z_a`$ in $`𝔤`$. These conditions are to be understood as supersymmetric generalisations of the gluing conditions written down in Section 3; henceforth $`R`$ is taken to be an automorphism of $`𝔤`$ which preserves the metric. At this point we do not need to impose any specific condition on $`S`$, since this will be fixed, as we will see in a moment, by supersymmetry considerations.
The gluing conditions (18) have to satisfy a similar consistency requirement as in the bosonic case. In this context, consistency means that the holomorphic SCFT is set equal to the antiholomorphic SCFT up to an automorphism of the $`N=1`$ SCA; in other words, at the boundary we must have
$$𝖳(z)=\overline{𝖳}(\overline{z})\text{and}𝖦(z)=\pm \overline{𝖦}(\overline{z}).$$
These conditions have been written down previously in , in the context of Kazama–Suzuki models.
The first requirement translates into a number of conditions on the matrices $`R`$ and $`S`$. Thus, from the quadratic terms in the currents we obtain that
$$R^T\eta R=\eta ,S^T\eta R=\pm \eta ,$$
which immediately implies that
$$S=\pm R,$$
(19)
as one would expect from supersymmetry. Further, from the cubic terms in the currents we have that
$$[S(Z_a),S(Z_b)]=\pm S([Z_a,Z_b]),[R(Z_a),S(Z_b)]=S([Z_a,Z_b]),$$
which, together with (19), implies that
$$[R(Z_a),R(Z_b)]=R([Z_a,Z_b]).$$
(20)
In other words, the conditions that $`R`$ must satisfy in order for the corresponding $`𝖣`$-brane configurations to preserve superconformal invariance match exactly the assumptions already made on $`R`$. Furthermore, it follows that these gluing conditions preserve the infinite-dimensional symmetry of the $`N=1`$ current algebra (13)-(15).
Finally, since we know from the bosonic case that $`R`$ must take one of two particular forms (11), we obtain that $`S`$ must have a similar form
$$S_I=\pm \left(\begin{array}{ccc}R_{11}& 0& 0\\ 0& R_{22}& 0\\ 0& 0& R_{33}\end{array}\right),S_{II}=\pm \left(\begin{array}{ccc}R_{11}& 0& 0\\ 0& 0& R_{23}\\ 0& R_{32}& 0\end{array}\right),$$
(21)
We therefore conclude that every bosonic configuration that we determined can be made into an $`N=1`$ supersymmetric configuration without having to impose additional conditions.
## 7. Spacetime supersymmetry
In this section we analyse the fraction of spacetime supersymmetry preserved by the $`𝖣`$-brane configurations we determined before. In the context of superconformal field theories spacetime supersymmetry appears as a by-product of $`N=2`$ superconformal invariance, being related, via bosonisation, to the $`U(1)`$ current. Instead of following this standard approach, here we will analyse the spacetime symmetry preserved by the $`𝖣`$-branes we found using a different route, which was described in .
We will therefore consider the spacetime supercharges to be constructed directly from the $`N=1`$ SCFT. To this end we introduce the fermionic fields $`\psi _i`$ for $`\mathrm{SL}(2,\mathrm{})`$, $`\chi _i`$ and $`\omega _i`$ for the two copies of $`\mathrm{SU}(2)`$, with $`i=1,2,3`$. Further, we choose five fermion bilinears and bosonise them into five scalar fields $`H_I`$, with $`I=1,\mathrm{},5`$ as follows
$$H_1=\psi _1\psi _2,H_2=\chi _1\chi _2,H_3=\omega _1\omega _2,$$
$$H_4=i\left(\sqrt{\frac{k_1}{k_2}}\chi _3+\sqrt{\frac{k_1}{k_3}}\omega _3\right)\psi _3,H_5=\left(\sqrt{\frac{k_1}{k_2}}\chi _3\sqrt{\frac{k_1}{k_3}}\omega _3\right)\lambda .$$
The corresponding spacetime supercharges will be required to be BRST invariant and to pass the GSO projection. This will yield
$$Q=𝑑ze^{\frac{\varphi }{2}}S(z),$$
(22)
where $`\varphi `$ is the scalar field which appears in the bosonised superghost system of the fermionic string, and the corresponding spin fields are given by (for a detailed discussion see )
$`S_1(z)`$ $`=e^{\frac{i}{2}(H_1+H_2+H_3+H_4+H_5)},`$
$`S_2(z)`$ $`=e^{\frac{i}{2}(H_1H_2H_3H_4H_5)},`$
$`S_3(z)`$ $`=e^{\frac{i}{2}(H_1+H_2+H_3H_4+H_5)},`$
$`S_4(z)`$ $`=e^{\frac{i}{2}(H_1H_2H_3+H_4H_5)},`$
$`S_5(z)`$ $`=\sqrt{{\displaystyle \frac{k_1}{k_2}}}e^{\frac{i}{2}(H_1H_2+H_3H_4+H_5)}+\sqrt{{\displaystyle \frac{k_1}{k_3}}}e^{\frac{i}{2}(H_1H_2+H_3+H_4H_5)},`$
$`S_6(z)`$ $`=\sqrt{{\displaystyle \frac{k_1}{k_2}}}e^{\frac{i}{2}(H_1+H_2H_3+H_4H_5)}+\sqrt{{\displaystyle \frac{k_1}{k_3}}}e^{\frac{i}{2}(H_1+H_2H_3H_4+H_5)},`$
$`S_7(z)`$ $`=\sqrt{{\displaystyle \frac{k_1}{k_2}}}e^{\frac{i}{2}(H_1H_2+H_3+H_4+H_5)}+\sqrt{{\displaystyle \frac{k_1}{k_3}}}e^{\frac{i}{2}(H_1H_2+H_3H_4H_5)},`$
$`S_8(z)`$ $`=\sqrt{{\displaystyle \frac{k_1}{k_2}}}e^{\frac{i}{2}(H_1+H_2H_3H_4H_5)}+\sqrt{{\displaystyle \frac{k_1}{k_3}}}e^{\frac{i}{2}(H_1+H_2H_3+H_4+H_5)}.`$
For any given $`𝖣`$-brane configuration, the boundary conditions satisfied by the fermionic fields will lead to a certain set of boundary conditions satisfied by the supercharges. Let us consider the those configurations described by $`R_I`$, with $`R_{11}`$ of the form of a spatial rotation. The corresponding boundary conditions satisfied by the fermions read (we use the notation introduced in )
$`\psi _1\mathrm{cos}\alpha \overline{\psi }_1\mathrm{sin}\alpha \overline{\psi }_2=0,`$ $`\chi _1\mathrm{cos}\beta \overline{\chi }_1\mathrm{sin}\beta \overline{\chi }_2=0,`$
$`\psi _2+\mathrm{sin}\alpha \overline{\psi }_1\mathrm{cos}\alpha \overline{\psi }_2=0,`$ $`\chi _2+\mathrm{sin}\beta \overline{\chi }_1\mathrm{cos}\beta \overline{\chi }_2=0,`$
$`\psi _3\overline{\psi }_3=0,`$ $`\chi _3\overline{\chi }_3=0,`$
$`\omega _1\mathrm{cos}\gamma \overline{\omega }_1\mathrm{sin}\gamma \overline{\omega }_2=0,`$
$`\omega _2+\mathrm{sin}\gamma \overline{\omega }_1\mathrm{cos}\gamma \overline{\omega }_2=0,`$ $`\lambda \pm \overline{\lambda }=0,`$
$`\omega _3\overline{\omega }_3=0,`$
where we have systematically ignored a $`\pm `$ sign coming from (21), which does not affect the fermion bilinears in the expression of $`S(z)`$. The $`\pm `$ sign in the boundary condition on $`\lambda `$ corresponds to having Neumann or Dirichlet boundary conditions on $`S^1`$, respectively. We will see however that only one choice will give rise to states preserving some spacetime supersymmetry. From the above relations it follows that the scalar fields $`H_I`$ satisfy the following boundary conditions:
$$H_I=\overline{H}_I,I=1,2,3,4,H_5=\pm \overline{H}_5,$$
where the $`\pm `$ sign in the boundary condition for $`H_5`$ reflects the one in the boundary condition along $`S^1`$. Therefore, in order to preserve some spacetime supersymmetry, we have to impose a Dirichlet boundary condition along the flat direction of the target space. In this way we obtain that these configurations describe odd-dimensional $`𝖣`$-branes (as we expect in the case of a type IIB theory). It immediately follows that, for these configurations, the spacetime supercharges satisfy
$$Q_\alpha =\overline{Q}_\alpha ,\alpha =1,\mathrm{},8.$$
(23)
Hence we conclude that the $`𝖣`$-brane configurations defined by $`R_I`$, with $`R_{11}`$ a spatial rotation in $`\mathrm{SO}(2,1)`$, and characterised by a Dirichlet condition along $`S^1`$ preserve half of the spacetime supersymmetry, and therefore are BPS states.
In order to analyse the configurations described by $`R_I`$ with $`R_{11}`$ given by a boost in $`\mathrm{SO}(2,1)`$ we need a slight change in the way we define the fermion bilinears and the corresponding supercharges. Essentially, we need to switch the places of $`\psi _2`$ and $`\psi _3`$ in the definition of $`H_1`$ and $`H_4`$. Then by using the following boundary conditions for the fermions on $`\text{AdS}_3`$
$`\psi _1\mathrm{cosh}\alpha \overline{\psi }_1\mathrm{sinh}\alpha \overline{\psi }_3=0,`$
$`\psi _2\overline{\psi }_2=0,`$
$`\psi _3\mathrm{sinh}\alpha \overline{\psi }_1\mathrm{cosh}\alpha \overline{\psi }_3=0,`$
and unchanged conditions for all the other fermions, we obtain that $`H_I=\overline{H}_I`$, for all $`I`$, provided we set again a Dirichlet boundary condition along $`S^1`$. From this it follows that the spacetime supercharges satisfy (23) as before. Hence all the corresponding odd-dimensional $`𝖣`$-branes describe BPS states that preserve half of the spacetime supersymmetry.
Finally, in the case where $`R_{11}`$ is given by a null rotation in $`\mathrm{SO}(2,1)`$, due to the particular form of the boundary conditions and to the nonlocal nature of the dependence of the spacetime supercharges on the fermionic fields, it is rather difficult to determine the fraction of spacetime supersymmetry preserved by this particular type of boundary states.
Let us now turn to the $`𝖣`$-brane configurations described by $`R_{II}`$. If we start with configurations characterised by a component $`R_{11}`$ of the form of a spatial rotation in $`\mathrm{SO}(2,1)`$ then the corresponding boundary conditions on the fermions read
$`\psi _1\mathrm{cos}\alpha \overline{\psi }_1\mathrm{sin}\alpha \overline{\psi }_2=0,`$ $`\chi _1\mathrm{cos}\beta \overline{\omega }_1\mathrm{sin}\beta \overline{\omega }_2=0,`$
$`\psi _2+\mathrm{sin}\alpha \overline{\psi }_1\mathrm{cos}\alpha \overline{\psi }_2=0,`$ $`\chi _2+\mathrm{sin}\beta \overline{\omega }_1\mathrm{cos}\beta \overline{\omega }_2=0,`$
$`\psi _3\overline{\psi }_3=0,`$ $`\chi _3\overline{\omega }_3=0,`$
$`\omega _1\mathrm{cos}\gamma \overline{\chi }_1\mathrm{sin}\gamma \overline{\chi }_2=0,`$
$`\omega _2+\mathrm{sin}\gamma \overline{\chi }_1\mathrm{cos}\gamma \overline{\chi }_2=0,`$ $`\lambda \pm \overline{\lambda }=0,`$
$`\omega _3\overline{\chi }_3=0,`$
This time, in order to be able to preserve some fraction of the spacetime supersymmetry, we must impose a Neumann boundary condition along $`S^1`$. Then we obtain
$$H_1=\overline{H}_1,H_2=\overline{H}_3,H_3=\overline{H}_2,H_4=\overline{H}_4,H_5=\overline{H}_5.$$
This in turn implies that the corresponding supercharges (where, we recall, $`k_2=k_3`$) will satisfy the following conditions
$$Q_\alpha =A_{}^{\beta }{}_{\alpha }{}^{}\overline{Q}_\beta ,\alpha ,\beta =1,\mathrm{},8$$
(24)
where the coefficients $`A_{}^{\beta }{}_{\alpha }{}^{}`$ are the elements of the matrix $`A`$
$$A=\left(\begin{array}{cccc}\mathbb{𝟙}& & & \\ & \mathbb{𝟙}& & \\ & & \sigma & \\ & & & \sigma \end{array}\right),$$
where $`\sigma =\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ and $`\mathbb{𝟙}`$ is the two-dimensional identity matrix. This means that these configurations too preserve half the spacetime supersymmetry and thus constitute BPS states.
Similar results hold in the case where the component $`R_{11}`$ of $`R_{II}`$ is given by a boost, provided we use the appropriate choice for the scalar fields $`H_I`$. And also similarly, we can not say much about the boundary states having the $`R_{11}`$ of the form of a null rotation.
The results are summarised in Table 1.
## 8. Conclusions
In this paper we have studied, using the SCFT framework and the boundary state formalism, the possible $`𝖣`$-brane configurations which one can consistently define in an $`\text{AdS}_3\times S^3\times S^3\times S^1`$ background characterised by a purely NS-NS B field.
We have analysed a certain type of gluing conditions (type-D, according to the nomenclature used in ), which are characterised by the fact that they preserve not only the superconformal structure of the background but also the underlying symmetry of the $`N=1`$ current algebra. We have seen that the solutions fall in two different classes. The first class, produced by gluing conditions defined in terms of inner automorphisms of the corresponding Lie algebra, describes D-brane configurations whose worldvolumes are products of shifted conjugacy classes
$$𝒞_{\mathrm{SL}(2,\mathrm{})}r_1\times 𝒞_{\mathrm{SU}(2)}r_2\times 𝒞_{\mathrm{SU}(2)}r_3,$$
giving thus rise to odd-dimensional D-branes embedded in $`\mathrm{SL}(2,\mathrm{})\times \mathrm{SU}(2)\times \mathrm{SU}(2)`$ and even-dimensional D-branes wrapped on the flat $`S^1`$. It is however the odd-dimensional D-branes that turn out to also preserve half of the spacetime supersymmetry of the background.
The second class of solutions, produced by gluing conditions defined in terms of outer automorphisms, describes D-brane configurations whose worldvolumes are products of shifted conjugacy classes in $`\mathrm{SL}(2,\mathrm{})`$ with twisted conjugacy classes in $`\times \mathrm{SU}(2)\times \mathrm{SU}(2)`$
$$𝒞_{\mathrm{SL}(2,\mathrm{})}r_1\times 𝒞_{\mathrm{SU}(2)\times \mathrm{SU}(2)}^{\stackrel{~}{T}}\stackrel{~}{r}.$$
We have studied in some detail the twisted conjugacy classes of $`\mathrm{SU}(2)\times \mathrm{SU}(2)`$, showing that they can be characterised as the inverse images, under the group multiplication, of the standard conjugacy classes of $`\mathrm{SU}(2)`$. In particular, they consist of two totally-geodesic three-spheres and a family of homologically trivial but homotopically nontrivial five-dimensional submanifolds diffeomorphic to $`S^2\times S^3`$. These solutions give rise to even-dimensional branes embedded in $`\mathrm{SL}(2,\mathrm{})\times \mathrm{SU}(2)\times \mathrm{SU}(2)`$ and odd-dimensional D-branes wrapped on the flat $`S^1`$, and it is again the odd-dimensional branes that preserve half of the spacetime supersymmetry.
## Acknowledgements
It is a pleasure to thank AA Tseytlin for many useful discussions in the early stages of this project. This work was partially supported by a PPARC Postdoctoral Fellowship. |
warning/0001/hep-ph0001020.html | ar5iv | text | # The 𝜉(2230) Meson and The Pomeron Trajectory
## Abstract
We examine the possibility that the $`\xi (2230)`$ meson is a member of the Pomeron trajectory. A method of connecting the $`\xi p\overline{p}`$ decay width and the $`pp`$ cross sections through the Pomeron residue function is presented. We have used a relativistic, singularity-free form factor to make the analytic continuation of the residue function between crossed channels. We predict that if the $`\xi (2230)`$ meson is a Pomeron, then it should have a $`\xi p\overline{p}`$ decay width of about 2 MeV.
Keywords: Regge poles, Pomeron, mesons.
PACS: 12.39Mk, 12.40.Nn, 12.40.Yx, 12.90.+b
Prior to the advent of the quantum chromodynamic (QCD) theory, the Regge theory was extensively used in describing the strong interaction. In spite of its success, the nature of the Pomeron, which plays a crucial role in explaining the asymptotic behavior of the hadron-hadron interaction, remains unclear. Over the years, precise phenomenological fits to high-energy data have determined that the spin and the mass of the Pomerons satisfy the relation $`\alpha (t)=\alpha (0)+\alpha ^{}t`$ with $`\alpha (0)1.08`$ and $`Re[\alpha ^{}]0.2`$(GeV)<sup>-2</sup>. If the Pomeron($`𝒫`$) trajectory is more than a theoretical device, then it should have member constituencies identifiable with physical states. To date, no such states have been found.
In this report we show that the $`\xi (2230)`$ meson could be a good candidate for Pomeron. This meson was first observed in the radiative decay of the $`J/\mathrm{\Psi }`$ to $`K\overline{K}`$ by the MARK III collaboration. Interest in $`\xi (2230)`$ resurged when the BES collaboration measured the reaction $`J/\mathrm{\Psi }\gamma \xi `$ and determined that this state has a mass of 2230 MeV with a total width $`\mathrm{\Gamma }_{tot}=20\pm 17`$ MeV. Its spin is still uncertain, being either $`2^{++}`$ or $`4^{++}`$, but will be determined in future measurements.
If $`J=2`$ is confirmed, then with $`J=Re[\alpha (t)]=2`$ at $`t=M_\xi ^2=(2.23)^2=5`$(GeV)<sup>2</sup> the $`\xi `$ satisfies the above-mentioned spin-mass relation of the Pomeron. In the following we shall go beyond this relation and show what would be the constraint on the $`\xi p\overline{p}`$ decay width, $`\mathrm{\Gamma }_{\xi p\overline{p}}`$, when $`\xi `$ is on the $`𝒫`$-trajectory. In particular, we show how we can relate it to the residue function of the Pomeron. We then use this relation and the measured $`pp`$ cross sections to predict the width. In our analysis, we have used for the first time a relativistic and singularity-free form factor. We emphasize that our analysis does not rely on any assumption about the subhadronic content of the $`\xi `$ meson. Hence, it is model-independent of the latter. The signifcance of our result will be elucidated.
Let us examine the $`pp`$ elastic scattering (denoted $`1234`$) mediated by the exchange of the $`\xi (2230)`$ meson. The corresponding $`t`$-channel process is, therefore, the $`p\overline{p}`$ scattering (denoted $`1\overline{3}\overline{2}4)`$ via the formation and decay of the $`\xi `$. Let $`m(0.94`$ GeV), $`M_\xi (2.23`$ GeV), and $`J_\xi `$ be, respectively, the mass of the $`p(\overline{p})`$, the mass of the $`\xi `$, and the spin of the $`\xi `$. The $`t`$-channel helicity Feynman amplitude can be written as
$$M_{\lambda _{\overline{2}}\lambda _4;\lambda _1\lambda _{\overline{3}}}(\overline{s},\overline{t})=4m^2C_I[4\pi (2J_\xi +1)]<\lambda _{\overline{2}}\lambda _4^{J_\xi }(\overline{s})\lambda _1\lambda _{\overline{3}}>d_{\lambda \lambda ^{}}^{J_\xi }(z_t)$$
(1)
where $`\lambda \lambda _1\lambda _{\overline{3}},\lambda ^{}\lambda _{\overline{2}}\lambda _4`$, and $`C_I=1/2`$ is the isospin factor. For clarity, we denote the square of the total c.m. energy and the momentum transfer in the $`s`$-channel by $`s`$ and $`t`$, respectively, and those in the $`t`$-channel by $`\overline{s}=(p_1+p_{\overline{3}})^2`$ and $`\overline{t}=(p_1p_{\overline{2}})^2`$. In the c.m. of the $`p\overline{p}`$, $`\overline{s}=4(m^2+k_t^2),\overline{t}=2k_t^2(1cos\theta _t)2k_t^2(1z_t)`$. Because $`p_{\overline{3}}=p_3,p_{\overline{2}}=p_2`$, we have $`\overline{s}=t`$ and $`\overline{t}=s`$. In Eq.(1)
$$<\lambda _{\overline{2}}\lambda _4^{J_\xi }(\overline{s})\lambda _1\lambda _{\overline{3}}>=\frac{G_\lambda ^{}H_{\lambda _{\overline{2}}\lambda _4;J_\xi }(\overline{s})G_\lambda H_{J_\xi ;\lambda _1\lambda _{\overline{3}}}(\overline{s})}{\overline{s}M_\xi ^2+iM_\xi \mathrm{\Gamma }_{tot}},$$
(2)
The $`H`$ denotes the form factor for the $`\xi p\overline{p}`$ vertex in the helicity basis and $`G`$ the coupling constant.
The Regge-pole amplitude due to the $`𝒫`$-trajectory can be written as
$$A_{\lambda _{\overline{2}}\lambda _4;\lambda _1\lambda _{\overline{3}}}(\overline{s},\overline{t})=C_I\frac{4\pi ^2(2\alpha +1)\beta _{\lambda \lambda ^{}}(t)(1)^{\alpha +\lambda }\frac{1}{2}[1+(1)^\alpha ]}{sin\pi (\alpha +\lambda ^{})}d_{\lambda \lambda ^{}}^\alpha (z_t)𝒜_{\lambda \lambda ^{}}^t,$$
(3)
where $`\beta _{\lambda \lambda ^{}}(t)`$ stands for $`\beta _{\lambda _{\overline{2}}\lambda _4;\lambda _1\lambda _{\overline{3}}}(\overline{s}t)`$. The contribution of the $`\xi `$ meson to the $`𝒫`$-trajectory is obtained by calculating their overlap as follows:
$$\frac{1}{2}_1^{+1}M_{\lambda _{\overline{2}}\lambda _4;\lambda _1\lambda _{\overline{3}}}(\overline{s},\overline{t})d_{\lambda \lambda ^{}}^{J_\xi }(z_t)𝑑z_t=lim_{\alpha J_\xi }\frac{1}{2}_1^{+1}A_{\lambda _{\overline{2}}\lambda _4;\lambda _1\lambda _{\overline{3}}}(\overline{s},\overline{t})d_{\lambda \lambda ^{}}^{J_\xi }(z_t)𝑑z_t.$$
(4)
By using $`_1^{+1}d_{\lambda \lambda ^{}}^{J_\xi }(z_t)d_{\lambda \lambda ^{}}^{J_\xi }(z_t)𝑑z_t=2/(2J_\xi +1)`$ and the following relation, which we have derived,
$$\frac{1}{2}_1^{+1}d_{\lambda \lambda ^{}}^\alpha (z_t)d_{\lambda \lambda ^{}}^{J_\xi }(z_t)𝑑z_t=\frac{sin\pi \alpha }{\pi (\alpha J_\xi )(\alpha +J_\xi +1)}$$
(5)
with
$$(1)^{J_\xi }\left(\frac{(\alpha \lambda )!(J_\xi +\lambda )!}{(\alpha +\lambda )!(J_\xi \lambda )!}\right)^{\frac{1}{2}},(\lambda =0,\pm 1),$$
(6)
we obtain
$`{\displaystyle \frac{16m^2\pi G^2H_{\lambda _{\overline{2}}\lambda _4;J_\xi }(t)H_{J_\xi ;\lambda _1\lambda _{\overline{3}}}(t)}{tM_\xi ^2+iM_\xi \mathrm{\Gamma }_{tot}}}`$ (7)
$`=`$ $`\left[{\displaystyle \frac{4\pi (2\alpha +1)\beta _{\lambda ,\lambda ^{}}(t)}{(\alpha J_\xi )(\alpha +J_\xi +1)}}{\displaystyle \frac{(1)^{\alpha +\lambda }\frac{1}{2}[1+(1)^\alpha ]sin\pi \alpha }{sin\pi (\alpha +\lambda ^{})}}\right]_{\alpha J_\xi }`$ (8)
$`=`$ $`{\displaystyle \frac{4\pi \beta _{\lambda \lambda ^{}}(t_r)/\alpha _R^{}}{tt_r+i\alpha _I(t_r)/\alpha _R^{}}}.`$ (9)
Hence, for even $`J_\xi `$ we have the identifications $`t_r=M_\xi ^2,\alpha _I(t_r)/\alpha _R^{}=M_\xi \mathrm{\Gamma }_{tot}`$, and
$$\beta _{\lambda \lambda ^{}}(t_r)=\alpha _R^{}4m^2G_\lambda ^{}H_{\lambda _{\overline{2}}\lambda _4;J_\xi }(t_r)G_\lambda H_{J_\xi ;\lambda _1\lambda _{\overline{3}}}(t_r).$$
(10)
In obtaining Eq.(9) we have defined $`\alpha _RRe(\alpha ),\alpha _IIm(\alpha )`$, and have used the Taylor expansion $`\alpha (t)=\alpha _R(t_r)+(\alpha _R^{}+i\alpha _I^{}(t_r))(tt_r)+\alpha _I(t_r)+\mathrm{}`$, with $`\alpha _R(t_r)J_\xi `$, $`\alpha _I\alpha _R,`$ and $`\alpha _I^{}\alpha _R^{}.`$
The diagonal matrix element of Eq.(10) is directly related to the $`\xi p\overline{p}`$ decay width. Using the method of , we obtain
$`{\displaystyle \frac{1}{2}}\mathrm{\Gamma }_{\xi p\overline{p}}`$ $`=`$ $`C_I{\displaystyle \frac{m^2}{M_\xi ^2}}{\displaystyle \frac{𝐪_𝐫}{16\pi ^2}}{\displaystyle \underset{\lambda _p\lambda _{\overline{p}}}{}}G_\lambda H_{\lambda _p\lambda _{\overline{p}}}(𝐪_{𝐫}^{}{}_{}{}^{2})^2`$ (11)
$`=`$ $`C_I{\displaystyle \frac{𝐪_𝐫}{64\pi ^2\alpha _R^{}M_\xi ^2}}{\displaystyle \underset{\lambda }{}}\beta _{\lambda \lambda }(t_r),`$ (12)
where $`𝐪_{𝐫}^{}{}_{}{}^{2}=t_r/4m^2=0.36`$ (GeV)<sup>2</sup> and $`\lambda =\lambda _p\lambda _{\overline{p}}`$. The $`𝐪_𝐫`$ is the value of the c.m. momentum $`𝐪`$ of the $`p\overline{p}`$ system at $`\overline{s}=M_\xi ^2t_r`$. On the other hand, the $`\beta _{\lambda \lambda ^{}}(t)`$ at $`t0`$ is related to the $`pp`$ cross sections. The spin-averaged total $`pp`$ cross section is given by
$$\sigma _{tot}^{pp}=\left(\frac{1}{4}\right)\left(\frac{4\pi }{k_s}\right)\frac{1}{8\pi k_s}\underset{\lambda _3\lambda _4;\lambda _1\lambda _2}{}Im[𝒜_{\lambda _3\lambda _4;\lambda _1\lambda _2}^s(s,t)]_{_{t=0}},$$
(13)
where $`k_s`$ and $`𝒜^s`$ are the $`s`$-channel c.m. momentum and amplitudes, and $`k_s^2s/4`$. Furthermore, $`_{\lambda _3\lambda _4;\lambda _1\lambda _2}Im[𝒜_{\lambda _3\lambda _4;\lambda _1\lambda _2}^s]=_{\lambda _{\overline{2}}\lambda _4;\lambda _1\lambda _{\overline{3}}}Im[𝒜_{\lambda _{\overline{2}}\lambda _4;\lambda _1\lambda _{\overline{3}}}^t]`$ due to crossing symmetry. There are sixteen helicity amplitudes for $`p\overline{p}`$ scattering. However, only five amplitudes are independent as a result of the parity, total spin conservation, and time-reversal invariance. We denote them as $`𝒜_i^t`$. The indices $`i=1,\mathrm{},5`$ correspond successively to $`(\lambda _{\overline{2}}\lambda _4;\lambda _1\lambda _{\overline{3}})=`$ $`(++;++),`$ $`(++;),`$ $`(+;+),`$ $`(+;+),`$ $`(++;+)`$ with $`\pm `$ denoting, respectively, the helicities $`\pm \frac{1}{2}`$. In the limit $`s\mathrm{}`$ and, hence, $`z_t\mathrm{}`$
$$d_{\lambda \lambda ^{}}^\alpha \frac{(1)^{(\lambda \lambda ^{})/2}(2\alpha )!(z_t/2)^\rho }{\sqrt{(\alpha +M)!(\alpha M+\lambda \lambda ^{})!(\alpha M)!(\alpha +M\lambda \lambda ^{})!}}$$
(14)
and
$$\left(\frac{z_t}{2}\right)^\rho (1)^\rho \left(\frac{s}{4m^2t}\right)^\rho $$
(15)
where $`\rho \alpha M+\lambda \lambda ^{}/2+\lambda +\lambda ^{}/2`$ and $`Mmax(\lambda ,\lambda ^{})`$. Upon introducing Eqs.(14) and (15) into Eq.(3), one sees that $`𝒜_1^t=𝒜_2^t`$ and $`𝒜_3^t=𝒜_4^t`$. Consequently,
$$\sigma _{tot}^{pp}=\frac{1}{2s}Im[4𝒜_1^t+8𝒜_5^t]_{t=0}.$$
(16)
The $`pp`$ elastic differential cross section is given by
$$\frac{d\sigma ^{pp}}{dt}=\frac{1}{16\pi s^2}\left(4𝒜_1^t^2+4𝒜_3^t^2+8𝒜_5^t^2\right).$$
(17)
In the above equations, $`𝒜_1^t,𝒜_3^t,𝒜_5^t`$ depend, respectively, on $`\beta _{00}(t),\beta _{11}(t),\beta _{10}(t)`$.
The $`\beta (t)`$ can be calculated from $`\beta (t_r)`$ with the use of a model that we specify below. First, we note that in deriving Eq.(10) one only requires $`lim_{tt_r}\alpha _R(t)=J_\xi `$. From the mathematical point of view, the position of $`t`$ at which the limit is taken can be anywhere along the Regge/Pomeron trajectory. Hence, we have the functional equality
$$\beta _{\lambda \lambda ^{}}(t)=\alpha _R^{}4m^2G_\lambda ^{}H_{\lambda _{\overline{2}}\lambda _4;J_\xi }(t)G_\lambda H_{J_\xi ;\lambda _1\lambda _{\overline{3}}}(t)$$
(18)
along this trajectory. In fact, the theory of analytic functions implies that the equality exists in a region where both $`\beta `$ and $`H`$ are analytic functions. In general the functional form of $`H(t)`$ can also depend on $`t`$. However, because there are no bound states and other resonances at $`t<M_\xi ^2`$, we can assume that the functional form of $`H(t)`$ is the same in the entire region $`tM_\xi ^2`$. From Eqs.(10) and (18) we have
$$\beta _{\lambda \lambda ^{}}(t)=\beta _{\lambda \lambda ^{}}(t_r)\frac{H_{\lambda _4,\lambda _4;J_\xi }(t)H_{J_\xi ;\lambda _1,\lambda _1}(t)}{H_{\lambda _4,\lambda _4;J_\xi }(t_r)H_{J_\xi ;\lambda _1,\lambda _1}(t_r)}.$$
(19)
Once the form factor $`H`$ is known, the $`\beta _{\lambda \lambda ^{}}`$ can be calculated from their values at $`t_r`$ which, by Eq.(12), are related to the decay width. It is advantageous to model the form factor in the $`LS`$basis because in this latter basis the form factor has a well-known $`q^L`$ threshold behavior that can be explicitly incorporated into the model. Because the helicity basis is related to the $`LS`$basis by a unitary transformation, one has $`_\lambda G_\lambda H_{\lambda _p,\lambda _{\overline{p}};J_\xi }^2=_{LS}g_{LS}F_{LS}^2`$ in Eq.(12). Here $`F_{LS}`$ and $`g_{LS}`$ denote, respectively, the form factor and the dimensionless coupling constant in the $`LS`$basis. For $`p\overline{p}`$ system with $`J^P=2^+`$, the parity conservation leads to $`L=1,3`$ and $`S=1`$ only. We will henceforth omit the subscript $`S`$.
In the literature, the most commonly used hadronic form factor is either of an exponential form or of a multipole form. However, these form factors are unsuitable to analyses involving channel-crossing. For example, the exponential form factor $`exp(t/\mathrm{\Lambda }^2)`$ is analytic in the $`s`$-channel where $`t0`$ but diverges in the $`t`$-channel where $`t>0`$. The multipole form factor $`[\mathrm{\Lambda }^2/(\mathrm{\Lambda }^2t)]^n(n=1,2,\mathrm{})`$ has no pole on the real axis of $`t`$ in the $`s`$-channel (where $`t0`$) but will have it in the $`t`$-channel (where $`t>0`$). Conversely, if we use $`exp(t/\mathrm{\Lambda }^2)`$ or $`[\mathrm{\Lambda }^2/(\mathrm{\Lambda }^2+t)]^n`$, then the situation will be reversed. We propose the following form factor which is singularity free:
$$F_L(t)^2=\left(\frac{t/4m^2}{𝐪_{𝐫}^{}{}_{}{}^{2}}\right)^L\left(\frac{e^{t_r/\lambda _t^2}}{R(x_t)+e^{t/\lambda _t^2}}\right)^2\left(\frac{1+e^{t_r/\lambda _s^2}}{R(x_s)+e^{t/\lambda _s^2}}\right)^2$$
(20)
where $`(t/4m^2)^Lf_L=𝐪^{2L}`$ reflects the $`q^L`$-dependence of the $`F_L`$. Because $`F_L(t_r)1`$ we have from Eq.(12) and from the relation between $`H`$ and $`F`$ that $`\mathrm{\Gamma }_{\xi p\overline{p}}g_1^2+g_3^2`$.
In Eq.(20), $`x_s(t2m^2)/\lambda _s^2`$ and $`x_t(t2m^2)/\lambda _t^2`$. The function $`R`$ is analytic and is defined by
$$R(x)=\frac{1}{2}(1+tanh(ax))=\frac{e^{ax}}{e^{ax}+e^{ax}}.$$
(21)
It rapidly changes from 0 to 1 when $`x`$ changes from $`<0`$ to $`>0`$, with $`a`$ controling the transition speed at $`x=0`$. With $`a>10`$, $`R`$ will be very close to a step function but does not have the discontinuity of the latter. Consequently, the form factor is a continuous function of $`t`$. In the $`s`$-channel ($`t0`$), $`F(t)^2f_L[1+exp(t/\lambda _t^2)]^2[exp(t/\lambda _s^2)]^2t^Lexp(2t/\lambda _s^2)`$, which goes to 0 as $`t\mathrm{}`$. Hence, $`\lambda _s`$ controls the form factor. Since $`F(t)`$ does not have the physical-channel energy $`s`$ as an explicit variable, it does not diverge with $`s`$. When $`t`$ reaches the $`t`$-channel physical domian ($`t>4m^2`$), $`F(t)^2f_L[exp(t/\lambda _t^2)]^2[1+exp(t/\lambda _s^2)]^2t^Lexp(2t/\lambda _t^2)0`$ when $`t\mathrm{}`$, exhibiting the correct energy behavior in the $`t`$-channel. Here in the $`t`$-channel the $`\lambda _t`$ controls the form factor. The above well-behaved $`t`$-dependence of $`F(t)`$ makes the latter a good tool for continuing $`\beta `$ between the direct and crossed channels.
Eqs.(12) and (16) to (19) allow us to predict the decay width $`\mathrm{\Gamma }_{\xi p\overline{p}}`$ from the measured $`pp`$ total and elastic cross sections, and vice versa. We have determined the form factor parameters from the experimental $`pp`$ total and elastic cross sections data at $`\sqrt{s}=53`$ and 62 GeV. As anticipated, $`\lambda _s`$ is mainly determined by the diffraction peak of the $`d\sigma /dt`$.
The result of our fit is given in Table I where the wider parameter ranges at 62 GeV are due to the larger experimental error bars. Table I indicates that if $`\xi `$ lies on the $`𝒫`$-trajectory, then it has a $`\mathrm{\Gamma }_{\xi p\overline{p}}`$ between 1.5 and 2 MeV. This decay width, in addition to spin, can be used to check whether $`\xi `$ is a Pomeron. We note that the BES collaboration cited nearly equal branching ratios for the four observed decays modes. Using these equal branching ratios and the above $`\mathrm{\Gamma }_{\xi p\overline{p}}`$ we conclude that the $`\mathrm{\Gamma }_{tot}`$ for the $`\xi `$ is at least 8 MeV. While this lower bound is compatible with the published data, it is, however, necessary to ascertain in future experiments that no other important decay channels than those observed in ref. are left out.
We recall that the theoretical modeling of the Pomeron started in the perturbative QCD sector. Various gluon-exchange models were proposed. However, all these models gave an $`\alpha (0)`$ much greater than the phenomenologcial value of 1.08. Recently, there is a growing interest in a possible connection between the Pomeron and the glueball. We emphasize that at this time there is no convincing proof that Pomeron is a glueball. As to the $`\xi (2230)`$ meson itself, it can either be a tensor glueball or a $`q\overline{q}`$glue mixture. Frank Close has shown that if $`\xi `$ is a tensor glueball then its total width $`\mathrm{\Gamma }_{tot}`$ would be of the order of 25 MeV; but no discussion was made on the $`\mathrm{\Gamma }_{\xi p\overline{p}}`$. Although the $`\mathrm{\Gamma }_{tot}`$ of the BES result falls within this limit, that measurement needs to be improved. We believe that high-statistics data on many decay modes of all the tensor states in this mass region are needed to pin down the gluonic content of the $`\xi `$. In this work, we do not study the microscopic composition of the $`\xi `$. Instead, we investigate the conditions that $`\xi `$ could be a Pomeron. If the Pomeron status of the $`\xi `$ is established, then the subhadronic structure of the $`\xi `$ will be directly relevant to that of the Pomeron. In this respect, our study will help clarify the Pomeron-glue connection.
In summary, we have derived the relation between the residue function of the Pomeron trajectory and decay widths of its member states. Meaurements of the spin and the $`p\overline{p}`$-decay width of the $`\xi `$ are important for determining if the $`\xi (2230)`$ meson could be the long-sought Pomeron.
We thank Drs. J.C. Peng, L.S. Kisslinger, P.R. Page, and L. Burakovsky for helpful discussions. |
warning/0001/astro-ph0001533.html | ar5iv | text | # A time-delay determination from VLA light curves of the CLASS gravitational lens B1600+434
## 1 Introduction
Gravitational lens systems can be used to determine the Hubble parameter, H<sub>0</sub> (Refsdal 1964). However, both a good mass model of the lens galaxy as well as a time delay between an image pair are necessary ingredients to accomplish this. The mass model can be constrained from the lens-image properties, whereas a time delay can be obtained through correlations between two image light curves. Intrinsic source variability should occur in all light curves, lagging by time delays which depend on the source and lens redshifts, the lens-mass distribution and the cosmological parameters, most prominently H<sub>0</sub> (e.g. Schneider et al. 1992).
Recently, time delays and values of H<sub>0</sub> were determined from several different gravitational lens (GL) systems: Q0957+561 (e.g. Schild & Thomson 1997; Kundić et al. 1997a; Haarsma et al. 1997, 1999; Press, Rybicki & Hewitt 1992; Pelt et al. 1994; Pelt et al. 1996; Grogin & Narayan 1996; Bernstein et al. 1997; Fischer et al. 1997; Falco et al. 1997; Bernstein & Fischer 1999; Barkana et al. 1999), PG1115+080 (e.g. Schechter et al. 1997; Courbin et al. 1997; Kundić et al. 1997b; Keeton & Kochanek 1997; Barkana 1997; Saha & Williams 1997; Pelt et al. 1998; Impey et al. 1998; Romanowsky & Kochanek 1999), B0218+357 (e.g. Biggs et al. 1999), B1608+656 (e.g. Fassnacht et al. 1999; Koopmans & Fassnacht 1999) and PKS 1830-211 (e.g. van Ommen et al. 1995; Lovell et al. 1998; Kochanek & Narayan 1992; Nair, Narasimha & Rao 1993).
Assuming that the lens galaxies in these GL systems can be described by an isothermal mass distribution, one finds that the values of H<sub>0</sub> derived from these GL systems – except for PG1115+080 – are consistent within their 1-$`\sigma `$ errors and agree with the local, SNe Ia and S-Z determinations of H<sub>0</sub> (e.g. Koopmans & Fassnacht 1999). However, one has to keep in mind that a similar mass-sheet or change in the radial mass profile of the lens galaxies introduce a similar change in the determination of H<sub>0</sub> from each of these GL systems (e.g. Falco, Gorenstein & Shapiro 1985; Gorenstein, Shapiro & Falco 1988; Grogin & Narayan 1992; Keeton & Kochanek 1997; Wucknitz & Refsdal 1999). Hence, also for mass models other than isothermal one might expect the values of H<sub>0</sub> to agree to first order. Strong deviations of the lens-galaxy mass distribution from isothermal, however, would lead to systematic differences between the values of H<sub>0</sub> determined from lensing and those determined from other methods.
In this paper, we present a determination of the time delay between the two lens images of the CLASS gravitational lens B1600+434, using the light curves obtained during an eight-month VLA 8.5-GHz monitoring campaign. In section 2, we describe the VLA observations of B1600+434 and the data reduction. In section 3, we apply the minimum dispersion method from Pelt et al. (1996) and the PRH-method (Press et al. 1992) to determine a time delay between the lens images and use that to estimate a tentative value for H<sub>0</sub>, keeping the above-mentioned problems with the mass-model degeneracies in mind. In section 4, our conclusions are summarized.
## 2 Data & Reduction
### 2.1 Observations
We observed the CLASS gravitational lens B1600+434 (Jackson et al. 1995; Jaunsen & Hjorth 1997; Koopmans, de Bruyn & Jackson 1998) with the VLA in A- and B-arrays at 8.5 GHz (X-band), during the period 1998 February 13 to October 14. The typical angular resolution of the radio images ranges from about 0.2 arcsec in A-array to 0.7 arcsec in B-array, sufficient to resolve the 1.4-arcsec double. In total we obtained 75 epochs of about 30 min each, including phase- and flux-calibrator observations and slewing time. The average time interval between epochs was 3.3 days. A typical observing run consisted of the sequence listed in Table 1. This sequence was repeated, typically twice, until about 30 min observing time was filled. Several epochs consisted of only 10 or 20 min, making it necessary to reduce the number of sequences, or the time spent on each of the sources. At the end of a sequence, we again observed 2 min on the phase calibrator J1549+506 (Fig.1; Patnaik et al. 1992).
The flux calibrator, B1634+627 (3C343), is a slightly extended steep-spectrum source (Fig.1; e.g. van Breugel et al. 1992), which should not be variable. Its flux density at 8.5 GHz is 0.84$`\pm `$0.01 Jy, which we obtained with the Westerbork-Synthesis-Radio-Telescope (WSRT) in December 1998. This value agrees with the average flux density of 0.83 Jy, obtained with the 26-m University of Michigan Radio Astronomy Observatory (UMRAO) telescope, which has been monitoring this source for the past 15 years. We use the WSRT flux-density determination to bring all flux densities to the correct absolute scale (Sect.2.4).
### 2.2 Calibration
The initial flux and phase calibration is done in the NRAO data-reduction package AIPS (version 15OCT98). We fixed the flux density of the phase calibrator at 1.17 Jy<sup>1</sup><sup>1</sup>1Taken from the VLA Calibrator Manual. at all epochs and subsequently solve for the telescope phase and gain solutions. The phase calibrator (J1549+506; Patnaik et al. 1992) is an unresolved flat-spectrum point source ($`<`$10 mas)<sup>2</sup><sup>2</sup>2 See the VLBA Calibrators List. at 8.5 GHz when observed with the VLA in A-array (Fig.1). A single delta function is therefore sufficient to describe the source structure. Because the phase calibrator is a flat spectrum and compact source, it could vary significantly over a monitoring period of eight months. Thus, the flux density of 1.17 Jy could be wrong by a large factor. We therefore use the flux calibrator to determine the proper flux-density scale and correct for ‘variability’ introduced in both images of B1600+434, due to intrinsic flux-density variability of the phase-calibrator. As we will see in Sect.2.4, this strategy works well.
We smooth (1-min intervals) and interpolate the phase and gain solutions between the phase calibrator scans for the entire observation period of about 30 min. These are then used to determine initial phase and gain solutions for the flux calibrator B1634+627 and the gravitational lens system B1600+434. No flagging was done in AIPS, because it is hard to assess if phase and/or gain errors can be corrected through self-calibration. Flagging is only done in the DIFMAP package (Shepherd 1997), where the visibility amplitudes and phases can be compared with a model of the source structure and possibly be corrected through self-calibration. Only if the latter fails, do we manually flag those visibilities which most strongly deviate from the model.
### 2.3 Imaging and model fitting
First we make maps of B1600+434 for all epochs in DIFMAP (version 2.2c), using a process of iterative model fitting and self-calibration. Because the two lens images are compact ($``$1 mas) at 8.5-GHz, determined from VLBA observations (Fig.3), and also show no extended emission from either the lens or quasar on mas to arcsec scales in MERLIN 5-GHZ observations (Koopmans et al. 1998), we can safely model the lens image structure by two delta functions, for which we determine the positions and flux densities. We fit these delta functions to the visibilities, until the model-$`\chi ^2`$ converges. We perform a phase-only self-calibration, using this model. Typically, this decreases the $`\chi ^2`$ of the model and the rms noise in the map significantly. We repeat this process several times, until no further decrease in either $`\chi ^2`$ or rms noise level is obtained. Finally, a global gain-self-calibration is performed, solving for small gain errors of each telescope. This gain-self-calibration does not significantly change the flux densities of the images (i.e. changes that are much less than the statistical error on the flux densities) and individual telescope gain corrections are typically less than a few percent. We repeat the process of model fitting and self-calibration, until it converges again. The visibilities are then compared with the best model and obviously errant points are manually flagged. Once more, we iterate between model fitting and self-calibration until convergence is reached. Finally, we average the visibilities over 300 sec and repeat the convergence process, after having flagged those points which still deviate significantly ($``$3-sigma) from the model.
The flux-density ratio between images A and B changes by at most a few tenths of a per cent before and after these calibration cycles, as long as there are no visibilities which deviate by orders of magnitude from the model. The reduced-$`\chi ^2`$-value of the final model fit typically lies between 1.0 and 1.1. The residual maps show no spurious features due to bad visibilities. Fig.3 shows a radio map of B1600+434, created from the combined A-array data-sets of seven epochs.
### 2.4 Flux calibration
We subsequently make maps of the flux calibrator, B1634+ 627, using the same procedure as for B1600+434. B1634+627 can be well represented by a single Gaussian component with a 1-$`\sigma `$ major axis of 90 mas, an axial ratio of 0.8 and a position angle of 70 deg, in broad agreement with the source structure as seen in 50-cm VLBI images (Nan et al. 1991). We use this component to model fit the image structure and self-calibrate the visibilities. We subsequently remove this Gaussian component and clean the map to find a better description of the extended structure of the source. The ratio between the flux density in the Gaussian component and the total flux density in all extended emission (seen with the VLA in A and B-arrays on a scale $``$10 arcsec) is 0.98. The ratio is independent of epoch (within the errors). The determination of the flux density using a single Gaussian component is better defined than the flux density derived from an iterative cleaning procedure. The latter procedure seems to introduce $``$1% errors, depending on how the maps were cleaned, what array was used and inside which box the total flux density was determined. Using only a single Gaussian to represent the source, does not involve cleaning or choosing a box size.
We determine the normalized light curve of the flux calibrator B1634+627 by dividing its flux-density light curve by the flux density of 0.84$`\pm `$0.01 Jy (Sect.2.1). The resulting curve shows a linear increase of approximately 5% over the eight-month observing period, which we attribute to a similar change in the flux density of the flat-spectrum phase calibrator, J1549+506. We also see that we have overestimated the flux density of B1634+627 by some 20–30%. In other words, our initial estimate of the flux density of the phase calibrator, based on the value given in the VLA calibrator manual, was 20–30% too high, which comes as no surprise for a variable flat-spectrum radio source.
From May 12 1998 onwards, we added a second flux calibrator<sup>3</sup><sup>3</sup>3Taken from the VLA Calibrator Manual., B1358+624 (Fig.1), to the observations, to estimate the reliability of the flux-density determination of B1634+627. We followed the same calibration and mapping procedure as for B1634+627. B1358+624 also shows the same 5% linear increase in flux density, supporting the idea that this is the result of variability of the phase calibrator. We scale the flux density curve of B1358+624 to fit the calibrated light curve of B1634+627 and find its flux density to be 1.14$`\pm `$0.02 Jy. The final normalized flux density curve of B1358+624 is also shown in Fig.2.
### 2.5 Light curves of B1600+434
To correct the light curves of B1600+434 for flux-density calibration errors, we divide them by the average of the normalized light curves of B1634+627 and B1358+624.
We assume that both flux calibrators do not vary intrinsically over the eight-month observing period. We also assume that the phase and gain solutions found from J1549+ 506 do not change significantly over a time span of several minutes, such that interpolation can be used to make a first-order correction for phase and gain errors in data of B1600+434, B1634+627 and B1358+624. This flux density correction removes the largest errors in the light curves of B1600+434, after which only statistical errors, second-order systematic errors and intrinsic variations are left. Any deviations between the two normalized light curves are most likely due to short-term calibration errors, either instrumental or atmospheric. The final flux-calibrated light curves of the lens images A and B are shown in Fig.4. The calibrated flux-densities are listed in Table 2.
Both light curves show a gradual decrease of about 0.02 mJy day<sup>-1</sup> over a period of 243 days (Sect.3.1). Superposed on this gradual long-term variability, image A also exhibits strong (up to 11% peak-to-peak) modulations on a time scale of a few days to several weeks. Image B only shows a few (up to 6% peak-to-peak) features, but separated in time by about one month or more. The modulation indices (i.e. fractional rms variability) around the gradual long-term decrease in flux density (indicated by the dashed lines in Fig.4) are 2.8% and 1.9% for images A and B, respectively.
### 2.6 External variability
In Koopmans & de Bruyn (1999) it is shown that most of the observed short–term variability is of external origin (at the 14.6-$`\sigma `$ confidence level). A number of possible causes of this short-term variability are examined: (i) scintillation caused by the ionized component of the Galactic ISM and (ii) radio microlensing of a core-jet structure by massive compact objects in the lens galaxy. Based on Galactic scintillation models (e.g. Narayan 1992; Taylor & Cordes 1993; Rickett et al. 1995), one predicts a strong increase in the modulation index towards longer wavelengths in the case of scintillation, whereas a strong decrease is observed in the long-term monitoring data obtained with the Westerbork-Synthesis-Radio-Telescope (WSRT) at 1.4 and 5 GHz (Koopmans & de Bruyn 1999). The quantitative decrease in the modulation index from 5 to 1.4 GHz, seen in this WSRT monitoring data, agrees remarkably well with that predicted from microlensing simulations, but differs by a factor of $``$8 from that predicted from the scintillation models. A strong case for the occurrence of radio microlensing in B1600+434 can therefore be made, which can complicate a straightforward determination of the time delay from these light curves. However, based on two independent methods of analysis (Sect.3.1), we are convinced that the obtained time delay is little affected by this external variability.
### 2.7 Error analysis
The errors on the light curves are a combination of thermal noise errors and systematic errors (e.g. modeling, self-calibration, instrumental, atmospheric, etc.). The noise errors on the $``$1 Jy flux density calibrators are of the order of 0.01% after a few minutes of integration. The noise error on each of the lens images is about 0.3%, determined from the residual maps (i.e. the radio image after subtracting the model of the source structure). This noise level agrees well with the theoretically expected value for the typical integration time of $``$10 min.
To estimate the systematic errors (i.e. systematic in the sense that they affect the flux-densities of both image A and B), we compare the two normalized light curves of B1634+627 and B1358+624. In Fig.2, we see that both curves follow each other extremely well. Their ratio has an rms scatter of 0.7%, determined from fitting a Gaussian to its distribution function. Assuming the errors on both normalized light curves are similar, the errors on the individual points are therefore 0.7$`/\sqrt{2}`$$``$0.5%. This error is probably a mixture of modeling, self-calibration and short-term atmospheric and instrumental effects, which are hard to remove. We conservatively assume that the data of B1600+434 contains a similar 0.5% error.
During the B-array observations, six points lie clearly outside the 3-$`\sigma `$ region (Fig.2, upper panel), whereas during the A- and BnA-array observations, the ratio seems much more more stable. This stability during A- and BnA-array observations is reflected in the extremely small scatter in the measured distance between the two lens images (Fig.4, upper panel), which is similar to the theoretical expectation value of $`\mathrm{\Delta }r_{\mathrm{AB}}`$=0.5 mas, where we used $`\mathrm{\Delta }r_{\mathrm{AB}}=\sqrt{2}\times \mathrm{\Delta }\theta /(2\mathrm{SNR})`$, with $`\mathrm{\Delta }\theta `$ being the beam size of 0.2 (0.7) arcsec in A-array (B-array) and SNR the signal-to-noise ratio of about 1/0.003$``$330.
The six ‘outliers’ are given an error equal to their difference in normalized flux density divided by $`\sqrt{2}`$, which is the expectation value if their errors are equal and drawn from a Gaussian distribution. The errors are 1.8–4.8%. Although this approach appears rather ad-hoc, we will later on in the determination of a time-delay use the light-curves both with and without these points, to investigate the effect they have on our analysis. As it will turn out, the effect is negligible (Sect.3.1.3).
To explain why we find these outliers, we investigated the system temperature (T<sub>sys</sub>) as function of time. Fast or systematic changes in T<sub>sys</sub> could indicate instrumental, atmospheric (e.g. precipitation) problems or electromagnetic interference. During day 223 (i.e. JD$``$2450858), T<sub>sys</sub> shows rapid changes of up to 20% on time scales of a few minutes, which could explain the large difference in the normalized flux density from the running mean and the large difference in the ratio between the normalized flux densities of B1634+627 and B1358+624 from unity. During the observations, cumuloform type clouds were forming with $``$50% sky coverage over the array<sup>4</sup><sup>4</sup>4Information obtained from the VLA observing logs, as kept by the VLA operators., possibly indicating strong interference caused by nearby thunderstorms (i.e. lightning). Also the other five outliers, during B-array, show some erratic behavior of T<sub>sys</sub>, although less serious than on day 223 and typically for only several of the telescopes. Only day 165 shows a $``$50% decrease in T<sub>sys</sub> over a 1h time interval for all telescopes. Because this decrease is relatively smooth, gain-self-calibration can solve for most of the errors. No such behavior is found between days 45 to 52 (in A-array) for example, which shows a much more gradual change in T<sub>sys</sub> and maximum differences less than 10%. On day 45, T<sub>sys</sub> behaves similar to epochs with no severe data problems. Its system temperature, however, is on average higher, which explains why the normalized flux density is lower. The higher system temperature for each of the telescopes is explained by the fact that it was snowing during the observations over the entire array (100% sky coverage). However, because T<sub>sys</sub> changes only gradually, the error on the corrected flux densities of images A and B will be similar to those of the other well behaved epochs.
Finally, we average the overlapping parts of the two normalized light curves, such that the errors decrease by a factor $`\sqrt{2}`$. Adding all errors quadratically, we find a total error of 0.8% on the light curves of B1600+434 A and B in the region where the normalized flux calibrator light curves do not overlap, and 0.7% where they do overlap. In Sect.3.1.4, we show that an error of 0.7–0.8% is statistically highly plausible, lending credibility to it.
Moreover, at a given epoch the systematic errors in the flux densities of images A and B are the same (i.e. because of their small angular separation of 1.4 arcsec, instrumental and atmospheric errors should be the same, as well as initial phase calibrator errors, which have been transferred to both images). Their flux density ratio is therefore much better determined and has an error of only $`\sqrt{2}\times `$0.3% $``$0.4%. The errors at different epochs, however, are independent, which is important if the light curves are shifted to determine a time delay (Sect.3.1).
## 3 Analysis
In this section we use the VLA 8.5-GHz light curves to put constraints on the time delay between the two lens images.
### 3.1 The time delay
A simple estimate of the time delay could be obtained, using the long-term gradients in both light-curves combined with the intrinsic flux-density ratio. Fitting a straight line to both curves (Fig.4) gives a flux-density decrease of $``$1.98$`10^2`$ mJy day<sup>-1</sup> for curve A and $``$1.87$`10^2`$ mJy day<sup>-1</sup> for curve B. Correcting the latter value by multiplying it by the flux-density ratio of 1.212 (see below) gives $``$2.27$`10^2`$ mJy day<sup>-1</sup>. This value is different from that of curve A by some 15%, indicating that the rate of decrease in the flux-density changes over the time-scale of the observations. One can therefore not simply divide the difference in flux-density between curve A and curve B (multiplied by the intrinsic flux-density ratio) by the rate of decrease in flux-density to obtain a time delay.
Moreover, the rapid strong modulations seen in the light curve of image A (Fig.4) makes interpolation questionable. This excludes the use of either the $`\chi ^2`$-minimization or cross-correlation methods in determining the time delay, because the light curves have to be resampled on a similar grid through some form of interpolation.
We have therefore chosen to use the non-parametric minimum-dispersion method developed by Pelt et al. (1996). In Sect.3.1.4 we will also derive the time delay using the PRH-method from Press, Rybicki & Hewitt (1992). As an additional constraint we use a flux density ratio of 1.212$`\pm `$0.005, determined from 28 epochs of VLA 8.5-GHz observations during a period of $``$4 months in 1996-1997 in which there was relatively little variability (C.B. Moore 1999, private communication). Because this period is significantly longer than the time delay between images A and B (Sect.3.1.2), the low rms variability implies that the above value should represent the intrinsic flux density ratio quite closely. However, we emphasize the preliminary nature of this value, which might still change slightly in a final analysis.
Because the fainter image (B) lags the brighter image (A) (Koopmans et al. 1998) and the flux densities of both images decrease almost linearly over a period of eight months, the flux-density ratio will on average be smaller than the flux-density ratio of 1.212. Shifting the light curve of image B back in time will increase the flux-density ratio between the overlapping parts of the light curves. When the shift in time equals the time delay between the lens images, the average flux-density ratio between the light curves should be 1.212. Hence, if we multiply the light curve of image B with the flux-density ratio of 1.212, we expect the minimum dispersion between the two light curves to occur near the intrinsic time delay.
#### 3.1.1 Minimum dispersion method
From the simple consideration that the observed flux-density ratio ($``$1.16) is smaller than the ratio of 1.212, we immediately see that image B lags image A in time, in agreement with lens models (Koopmans et al. 1998). Thus, the time delay $`\mathrm{\Delta }t_{\mathrm{B}\mathrm{A}}`$ is positive. A composite light curve is created by multiplying light curve B with the flux-density ratio of 1.212 and shifting it backward in time by $`\mathrm{\Delta }t_{\mathrm{B}\mathrm{A}}`$. The dispersion of the composite light curve is calculated as in Pelt et al. (1996), using the dispersion measure
$$D_4^2(\mathrm{\Delta }t)=\frac{_{n=1}^{N1}_{m=n+1}^NS_{\mathrm{m},\mathrm{n}}W_{m,n}G_{m,n}(C_\mathrm{n}C_\mathrm{m})^2}{_{m=n+1}^NS_{\mathrm{m},\mathrm{n}}W_{m,n}G_{m,n}}$$
where $`C_\mathrm{n}`$ is the $`n`$-th point on the composite light curve and $`G_{m,n}`$=1 (0), if $`C_\mathrm{n}`$ and $`C_\mathrm{m}`$ are from different (the same) light curves. We calculate the dispersion for all time delays $`\mathrm{\Delta }t_{\mathrm{B}\mathrm{A}}`$=0–100 days, in steps of 1 day. This process is performed for four different decorrelation time scale, $`\delta `$=1,2,4,8, in units of the average time span between epochs, i.e. 3.3 days. We use a decorrelation weight function $`S_{\mathrm{m},\mathrm{n}}=\mathrm{exp}[(t_\mathrm{m}t_\mathrm{n})^2/(2\mathrm{\Delta }^2)]`$, where $`\mathrm{\Delta }=3.3\times \delta `$ is the decorrelation time scale in days. The statistical weights are $`W_{m,n}=(W_nW_m)/(W_n+W_m)`$, where $`W_i=1/\sigma _i^2`$ and $`\sigma _i`$ the 1-$`\sigma `$ error on the flux density at the $`i`$-th epoch.
In Fig.5 the dispersion $`D_4^2`$ is plotted versus the time delay, for $`\delta `$=1,2,4 and 8. The dispersion minimizes near a time delay of 46 to 51 days.
#### 3.1.2 Median time delay and statistical error range
To determine a statistical confidence region for the time delay, we performed Monte-Carlo simulations for $`\delta `$=1,2,4 and 8. First, we re-sampled the light curves, using the sampling-interval distribution determined from the VLA observations. Because the light curves exhibit variability due to external causes (Sect.2.6), we do not create a composite light curve, by combining the image light curves, as has been done for B0218+357 (Biggs et al. 1999) and B1608+656 (Fassnacht et al. 1999). Such a light curve only resembles the true underlying light curve, if all variability were intrinsic to the source, which is not the case for B1600+434. We therefore linearly interpolate the observed light curves and errors to obtain the flux densities and errors at the re-sampled intervals. Subsequently, Gaussian distributed errors are added to (i) each point on the re-sampled light curves (1-$`\sigma `$ equal to the interpolated flux-density error) and (ii) the assumed intrinsic flux-density ratio ($`\sigma _r`$=0.005). The process described above was repeated 5000 times for each decorrelation time scales. The delays were stored, where $`D_4^2`$ minimizes. The resulting time-delay probability distribution functions (PDF) for the four decorrelation time scales are shown in Fig.5. From the PDFs of $`\mathrm{\Delta }t_{\mathrm{B}\mathrm{A}}`$, we determine the median values for the time delay and the statistical confidence regions. The final result of this procedure, for the different decorrelation scales, are listed in Table 3.
For $`\delta `$=1, multiple strong peaks are found (Fig.5). For $`\delta `$$`>`$1, only one peak is found. The small decorrelation time scale for $`\delta =1`$ makes the dispersion measure especially sensitive to the modulations in both light curves. These peaks, however, are much stronger in the light curve of image A and the result of external causes (Sect.2.6; Koopmans & de Bruyn, 1999). Hence, a somewhat larger decorrelation time scale will give a better estimate of the time delay, because it is less sensitive to these modulations (i.e. it averages over these modulations). For very large decorrelation time scales, however, one becomes sensitive to the fact that both light curves show a long-term gradient. The gradient introduces a systematic difference in the flux-density level between points on the two different light curves, if they are separated by a large time interval. This artificially increases the dispersion with increasing $`\delta `$, as can be seen in Fig.5. This effect also seems to increase the width of the 95% statistical confidence interval. The intermediate decorrelation time scales ($`\delta `$=2 and 4) therefore seem to give a better estimate of the time delay, as it avoids most of these problems.
For the median value of the time delay we take the average of $`\delta `$=2 and $`\delta `$=4, which seem to give the most stable solution (see also Sect.3.1.3 and Fig.6). For the 68% and 95% statistical confidence regions, we conservatively take their combined maximum ranges. Hence,
$`\begin{array}{ccccc}\mathrm{\Delta }t_{\mathrm{B}\mathrm{A}}\hfill & =\hfill & 47_6^{+5}\hfill & \mathrm{d}\hfill & \text{ (68\%)}\hfill \\ & =\hfill & 47_9^{+12}\hfill & \mathrm{d}\hfill & \text{ (95\%)},\hfill \end{array}`$
which we take as our best estimate of the time delay between the images in the gravitational lens system B1600+434 and the statistical confidence intervals.
#### 3.1.3 Systematic uncertainties in the time delay
Several systematic uncertainties remain, which we will investigate below:
1. The first has to do with the six outliers on the light-curves (Sect.2.7). Although we gave them significantly larger errorbars, they can still affect the determination of the time delay.
2. The choice of the decorrelation time scale seems to influence the median time delay (Sect.3.1.2). Larger decorrelation time scales give larger median time delays.
3. The flux density ratio that we used in our analysis (Sect.3.1) is preliminary and the error was assumed to be Gaussian, which might not be the case.
To address the first two points, we ran Monte-Carlo simulations (500 redistributions), for decorrelation time scales of $`\mathrm{\Delta }`$=3,5,…,25 d, using all epochs shown in Fig.4. We repeated this without the six outliers (Sect.2.7). The results are shown in Fig.6. The values for the median time delays range between 43 to 51 d, for the assumed range of decorrelation time-scales. The width of the 95% statistical confidence interval seems to minimize in the range $`\mathrm{\Delta }`$$``$10–15 d, as was already noted in the previous section.
To estimate the effect of a wrongly chosen intrinsic flux-density ratio, we also ran models for $`r_{\mathrm{AB}}`$=1.202 and 1.222, which are the assumed intrinsic flux-density ratio plus-minus twice its estimated error. The first value decreases the median time delay systematically by about 8 days, whereas the latter value increases it by about 7 days. We thus take the range of about $``$8 to +7 days as a good indication of the maximum systematic error range.
#### 3.1.4 The PRH method
We also applied the PRH method (Press, Rybicki & Hewitt 1992), to address the question if the time delay obtained in Sect.3.1.2 might critically depend on the method that was used. We use the implementation of this method in the ESO-MIDAS (version 96NOV) data-reduction package. We tried several different analytical functions to describe the structure functions of the observed light curves and find no strong dependence of the final result (i.e. differences of $``$1 day in the time delay) on the precise functional form of the structure function. We exclude the six outliers (Sect.2.7) from the analysis and find a delay of $`\mathrm{\Delta }t_{\mathrm{B}\mathrm{A}}`$=$`48_4^{+2}`$ d (2-$`\sigma `$), where the formal error is derived by varying the time-delay until $`\mathrm{\Delta }\chi ^2`$ has increased by 4.0. This error is smaller than that derived in Sect.3.1.2, because it does not include all the uncertainties that we took into account in the Monte-Carlo simulations. However, we see that both the minimum-dispersion method and the PRH method give consistent results within their respective 1-$`\sigma `$ error regions. Thence, we conclude that the resulting value of $`\mathrm{\Delta }t_{\mathrm{B}\mathrm{A}}`$ does not strongly depend on the method that is used, at least in the case of B1600+434.
We furthermore find a minimum-$`\chi ^2`$ value of $``$134 at $`\mathrm{\Delta }t_{\mathrm{B}\mathrm{A}}`$=48 d for 137 degrees of freedom, which is statistically highly plausible. This lends strong credibility to our error analysis (Sect.2.7), suggesting that the error on the individual flux-density points are indeed 0.7–0.8% for most epochs.
### 3.2 The Hubble parameter and slope of the radial mass profile of the lens-galaxy dark-matter halo
To estimate a tentative value for the Hubble parameter (H<sub>0</sub>), we use the mass models for B1600+434 from Koopmans et al. (1998). In doing this, one should keep in mind the difficulties and degeneracies that were mentioned in Sect.1. The values derived here should therefore be regarded as indicative, as long as no better constraints on the radial mass profile of the lens galaxy are obtained.
Spectroscopic observations of the lens system on 1998 April 20 with the W.M. Keck-II telescope, have recently confirmed the assumption that the nearby companion galaxy (G2) has the same redshift ($`z`$=0.41; Fassnacht et al., in preparation) as the lensing galaxy (G1). For a lens galaxy with an oblate isothermal dark-matter halo, the relation between the time delay and Hubble parameter was then found to be $`\text{H}_0=50\times [(54_9^{+11}\text{ d})/\mathrm{\Delta }t_{\mathrm{B}\mathrm{A}}]`$ km s<sup>-1</sup> Mpc<sup>-1</sup>, for $`\mathrm{\Omega }_\mathrm{m}`$=1 and $`\mathrm{\Omega }_\mathrm{\Lambda }`$=0. The errors indicate the maximum range of the isothermal-model time delays from Koopmans et al. (1998). Recently, the time-delay dependence of H<sub>0</sub> found for this isothermal mass model was corroborated by Maller et al. (1999), who did a similar analysis of B1600+434, using a deep NICMOS-F160W HST exposure. Combining this relation with the median time delay (Sect.3.1.2), the Hubble parameter then becomes $`\text{H}_0=57_{11}^{+14}`$ km s<sup>-1</sup> Mpc<sup>-1</sup> (95%), for $`\mathrm{\Omega }_\mathrm{m}`$=1 and $`\mathrm{\Omega }_\mathrm{\Lambda }`$=0. A maximum systematic error between $``$15 to +26 km s<sup>-1</sup> Mpc<sup>-1</sup> is estimated from the combination of model and systematic time-delay errors. This error does not include the uncertainty in the slope of the radial mass profile. In fact, for a Modified Hubble Profile (MHP) halo mass model we find a significantly higher value for the Hubble parameter, H<sub>0</sub>=74$`{}_{15}{}^{}{}_{}{}^{+18}`$ km s<sup>-1</sup> Mpc<sup>-1</sup> (95%), with a maximum systematic error between $``$22 to +22 km s<sup>-1</sup> Mpc<sup>-1</sup>. For a flat universe with $`\mathrm{\Omega }_0`$=0.3 and $`\mathrm{\Omega }_\mathrm{\Lambda }`$=0.7, these values of H<sub>0</sub> increase by 5.4%. At this points B1600+434 is not a GL system from which H<sub>0</sub> can be constrained reliably.
In Wucknitz & Refsdal (1999) it was shown that in general the time delay is a strong function of the slope of the radial mass profile of the mass distribution of the lens galaxy. Because the lens-image properties can be recovered to first order for a range of different slopes, the expected time delay from a GL system as function of H<sub>0</sub> is also a strong function of this slope. This degeneracy between the time delay and the slope makes an accurate determination of H<sub>0</sub> from GL systems difficult, especially in two-image lens systems like B1600+434. Hence, if H<sub>0</sub> can be determined with greater accuracy ($``$10% error), the measured time delay can be used to constrain the slope of the radial mass profile of the dark-matter halo.
## 4 Conclusions
We have monitored the CLASS gravitational lens B1600+ 434 at 8.5 GHz with the VLA in A and B-arrays, during the period from February to October 1998. The light curves show a nearly linear decrease of 18-19% in the flux density of both lens images over this period. However, image A also shows rapid variability (up to 11% peak-to-peak) on scales of days to weeks, whereas image B shows significantly less short-term variability (upto 6% peak-to-peak). The short-term variability occurs over an observing period that is much longer than any conceivable time delay.
In Koopmans & de Bruyn (1999) it is shown that the short-term variability is predominantly of external origin. Two plausible explanations of this external variability are suggested: scintillation caused by the ionized component of the Galactic ISM or radio microlensing of a core-jet structure by massive compact objects in the lens galaxy. Both possibilities are examined in more detail in Koopmans & de Bruyn (1999). A comparison of the result from the VLA 8.5-GHz, presented in this paper, and multi-frequency (1.4 and 5-GHz) WSRT monitoring data with the expected dependence of scintillation on frequency (e.g. Narayan 1992; Taylor & Cordes 1993; Rickett et al. 1995), shows that the scintillation hypothesis underestimates the short-term rms variability at 1.4 GHz by a factor $``$8. Within the uncertainties, the microlensing hypothesis predicts the correct frequency-dependence of the short-term rms variability as function of frequency. The radio-microlensing hypothesis therefore seems most viable at present.
From the VLA 8.5-GHz light curves, we determined a median time delay of $`\mathrm{\Delta }\mathrm{t}_{\mathrm{B}\mathrm{A}}`$=47$`{}_{9}{}^{}{}_{}{}^{+12}`$ days (95% statistical confidence) between the lens images. A maximum systematic error between $``$8 and +7 d is estimated. We used the minimum-dispersion method from Pelt et al. (1996), but find the same time-delay from the PRH-method from Press et al. (1992).
Combining this with the isothermal lens mass models from Koopmans et al. (1998), the Hubble parameter would become H<sub>0</sub>=57$`{}_{11}{}^{}{}_{}{}^{+14}`$ km s<sup>-1</sup>Mpc<sup>-1</sup> (95%) for $`\mathrm{\Omega }_\mathrm{m}`$=1 and $`\mathrm{\Omega }_\mathrm{\Lambda }`$=0. A maximum systematic error between $``$15 and +26 km s<sup>-1</sup> Mpc<sup>-1</sup> is estimated. Similarily, the MHP mass models would give H<sub>0</sub>=74$`{}_{15}{}^{}{}_{}{}^{+18}`$ km s<sup>-1</sup> Mpc<sup>-1</sup> (95%), with a maximum systematic error between $``$22 and +22 km s<sup>-1</sup> Mpc<sup>-1</sup>, for the same cosmological model. We hope to improve on the determination of this time delay with an ongoing three-frequency VLA monitoring campaign (June 1999 to Feb. 2000). Because of the degeneracy between the slope of the radial mass profile and the expected time delay between the lens images as function of H<sub>0</sub>, the above-given values of H<sub>0</sub> should be regarded as indicative.
If H<sub>0</sub> can be determined accurately from independent methods and no extra constraints on the lens model can be found, it is more interesting to use that observed time delay to constrain the slope of radial mass profile of the dark-matter halo around the lensing edge-on spiral galaxy in B1600+434.
###### Acknowledgements.
We like to thank Chris Moore for useful discussions and several good suggestions to improve the manuscript. We thank Phillip Helbig, Peter Wilkinson and Ian Browne for carefully reading the manuscript and giving suggestions for improvement. We thank Jaan Pelt and Mark Neeser for helping to implement the PRH method in MIDAS. LVEK and AGdeB acknowledge the support from an NWO program subsidy (grant number 781-76-101). This research was supported in part by the European Commission, TMR Program, Research Network Contract ERBFMRXCT96-0034 ‘CERES’. The National Radio Astronomy Observatory is a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc. The Westerbork Synthesis Radio Telescope (WSRT) is operated by the Netherlands Foundation for Research in Astronomy (ASTRON) with the financial support from the Netherlands Organization for Scientific Research (NWO). This research has made use of data from the University of Michigan Radio Astronomy Observatory which is supported by funds from the University of Michigan. |
warning/0001/physics0001058.html | ar5iv | text | # Calculation of positron binding to silver and gold atoms
## I Introduction
Positron binding by neutral atoms has not been directly observed yet. However, intensive theoretical study of the problem undertaken in the last few years strongly suggests that many atoms can actually form bound states with a positron (see, e.g. ). Most of the atoms studied so far were atoms with a relatively small value of the nuclear charge $`Z`$. It is important to extend the study to heavy atoms. The main obstacle in this way is the rapid rise of computational difficulties with increasing number of electrons. However, as we show in this paper, an inclusion of relativistic effects is also important. The role of these effects in positron binding to atoms has not been truly appreciated. Indeed, one can say that due to strong Coulomb repulsion a positron cannot penetrate to short distances from the nucleus and remains non-relativistic. However, the positron binding is due to interaction with electrons which have large relativistic corrections to their energies and wave functions. The binding energy is the difference between the energies of a neutral atom and an atom bound with a positron. This difference is usually small. On the other hand, relativistic contributions to the energies of both systems are large and there is no reason to expect they are the same and cancel each other. Therefore, some relativistic technique is needed to study positron binding by heavy atoms.
For both light and heavy atoms the main difficulty in calculations of positron interaction comes from the strong electron-positron Coulomb attraction. This attraction leads to virtual positronium (Ps) formation . One can say that it gives rise to a specific short-range attraction between the positron and the atom, in addition to the usual polarizational potential which acts between a neutral target and a charged projectile . This attraction cannot be treated accurately by perturbations and some all-order technique is needed. In our earlier works we used the Ps wave function explicitly to approximate the virtual Ps-formation contribution to the positron-atom interaction and predicted $`e^+`$Mg, $`e^+`$Zn, $`e^+`$Cd and few other bound states. The same physics may also explain the success of the stochastic variational method in positron-atom bound state calculations (see, e.g. and Refs. therein). In this approach the wave function is expanded in terms of explicitly correlated Gaussian functions which include factors $`\mathrm{exp}(\alpha r_{ij}^2)`$ with inter-particle distances $`r_{ij}`$. Using this method Ryzhikh and Mitroy obtained positron bound states for a whole range of atoms (Be, Mg, Zn, Cu, Ag, Li, Na, K, etc.). This method is well suited for few-particle systems. Its application to heavier systems is done by considering the Hamiltonian of the valence electrons and the positron in the model potential of the ionic core. However, for heavier atoms, e.g., Zn, the calculation becomes extremely time consuming , and its convergence cannot be ensured.
Another non-perturbative technique is the configuration interaction (CI) method widely used in standard atomic calculations. This method was applied to the positron-copper bound state in . In this work the single-particle orbitals of the valence electron and positron are chosen as Slater-type orbitals, and their interaction with the Cu<sup>+</sup> core is approximated by the sum of the Hartree-Fock and model polarization potentials. The calculation shows slow convergence with respect to the number of spherical harmonics included in the CI expansion, $`L_{\mathrm{max}}=10`$ being still not sufficient to extrapolate the results reliably to $`L_{\mathrm{max}}\mathrm{}`$.
In their more recent work the same authors applied the CI method to a number of systems consisting of an atom and a positron. These include PsH, $`e^+`$Cu, $`e^+`$Li, $`e^+`$Be, $`e^+`$Cd and CuPs. In spite of some improvements to the method they still regard it as a “tool with which to perform preliminary investigations of positron binding” .
In our previous paper we developed a different version of the CI method for the positron-atom problem . The method is based on the relativistic Hartree-Fock method (RHF) and a combination of the CI method with many body perturbation theory (MBPT). This method was firstly developed for pure electron systems and its high effectiveness was demonstrated in a number of calculations . In the paper it was successfully applied to the positron binding by copper. There are several important advances in the technique compared to the standard non-relativistic CI method which make it a very effective tool for the investigation of positron binding by heavy atoms.
1. The method is relativistic in the sense that the Dirac-Hartree-Fock operator is used to construct an effective Hamiltonian for the problem and to calculate electron and positron orbitals.
2. $`B`$-splines in a cavity of finite radius $`R`$ were used to generate single-particle basis sets for an external electron and a positron. The $`B`$-spline technique has the remarkable property of providing fast convergence with respect to the number of radial functions included into the calculations . Convergence can be further controlled by varying the cavity radius $`R`$ while the effect of the cavity on the energy of the system is taken into account analytically . Convergence was clearly achieved for the $`e^+`$Cu system in Ref. and for the $`e^+`$Ag and $`e^+`$Au systems as presented below.
3. We use MBPT to include excitations from the core into the effective Hamiltonian. This corresponds to the inclusion of the correlations between core electrons and external particles (electron and positron) and of the effect of screening of the electron-positron interaction by core electrons. These effects are also often called the polarization of the core by the external particles. We include them in a fully ab initio manner up to the second order of the MBPT.
In the present paper we apply this method to the problem of positron binding by silver and gold atoms. Using a similar technique we also calculate electron affinities for both these atoms. Calculations for negative ions serve as a test of the technique used for positron-atom binding. We also study the role of the relativistic effects in neutral silver and gold, silver and gold negative ions and silver and gold interacting with a positron. This is done by varying the value of the fine structure constant $`\alpha `$ towards its non-relativistic limit $`\alpha =0`$.
## II Theory
A detailed description of the method was given in Ref. . We briefly repeat it here emphasizing the role of the relativistic effects. We use the relativistic Hartree-Fock method in the $`V^{N1}`$ approximation to obtain the single-particle basis sets of electron and positron orbitals and to construct an effective Hamiltonian.
The two-particle electron-positron wave function is given by the CI expansion,
$$\mathrm{\Psi }(𝐫_e,𝐫_p)=\underset{i,j}{}C_{ij}\psi _i^e(𝐫_e)\psi _j^p(𝐫_p),$$
(1)
where $`\psi _i^e`$ and $`\psi _j^p`$ are the electron and positron orbitals respectively. The expansion coefficients $`C_{ij}`$ are determined by the diagonalization of the matrix of the effective CI Hamiltonian acting in the Hilbert space of the valence electron and the positron,
$`H_{\mathrm{eff}}^{\mathrm{CI}}`$ $`=`$ $`\widehat{h}_e+\widehat{h}_p+\widehat{h}_{ep},`$ (2)
$`\widehat{h}_e`$ $`=`$ $`c𝜶𝒑+(\beta 1)mc^2{\displaystyle \frac{Ze^2}{r_e}}+V_d^{N1}\widehat{V}_{exch}^{N1}+\widehat{\mathrm{\Sigma }}_e,`$ (3)
$`\widehat{h}_p`$ $`=`$ $`c𝜶𝒑+(\beta 1)mc^2+{\displaystyle \frac{Ze^2}{r_p}}V_d^{N1}+\widehat{\mathrm{\Sigma }}_p,`$ (4)
$`\widehat{h}_{ep}`$ $`=`$ $`{\displaystyle \frac{e^2}{|𝐫_e𝐫_p|}}+\widehat{\mathrm{\Sigma }}_{ep},`$ (5)
where $`\widehat{h}_e`$ and $`\widehat{h}_p`$ are the effective single-particle Hamiltonians of the electron and positron, and $`\widehat{h}_{ep}`$ is the effective electron-positron two-body interaction. Apart from the relativistic Dirac operator, $`\widehat{h}_e`$ and $`\widehat{h}_p`$ include the direct and exchange Hartree-Fock potentials of the core electrons, $`V_d^{N1}`$ and $`V_{exch}^{N1}`$, respectively. The additional $`\widehat{\mathrm{\Sigma }}`$ operators account for correlations involving core electrons. $`\mathrm{\Sigma }_e`$ and $`\mathrm{\Sigma }_p`$ are single-particle operators which can be considered as a self-energy part of the correlation interaction between an external electron or positron and core electrons. These operators are often called “correlation potentials” due to the analogy with the non-local exchange Hartree-Fock potential. $`\mathrm{\Sigma }_{ep}`$ represents the screening of the Coulomb interaction between external particles by core electrons (see for a detailed discussion).
To study the role of the relativistic effects we use the form of the operators $`h_e`$ and $`h_p`$ in which the dependence on the fine structure constant $`\alpha `$ is explicitly shown. Single-particle orbitals have the form
$`\psi (𝐫)_{njlm}={\displaystyle \frac{1}{r}}\left(\begin{array}{c}f_n(r)\mathrm{\Omega }(𝐫/r)_{jlm}\\ i\alpha g_n(r)\stackrel{~}{\mathrm{\Omega }}(𝐫/r)_{jlm}\end{array}\right).`$ (8)
Then the RHF equations
$`(h_iϵ_n)\psi _n^i=0,(i=e,p)`$take the following form
$`f_n^{}(r)+{\displaystyle \frac{\kappa _n}{r}}f_n(r)[2+\alpha ^2(ϵ_n\widehat{V})]g_n(r)=0`$ (9)
$`g_n^{}(r){\displaystyle \frac{\kappa _n}{r}}g_n(r)+(ϵ_n\widehat{V})f_n(r)=0,`$ (10)
where $`\kappa =(1)^{l+j+1/2}(j+1/2)`$ and $`V`$ is the effective potential which is the sum of the Hartree-Fock potential and correlation potential $`\mathrm{\Sigma }`$:
$`\widehat{V}`$ $`=`$ $`{\displaystyle \frac{Ze^2}{r_e}}+V_d^{N1}\widehat{V}_{exch}^{N1}+\widehat{\mathrm{\Sigma }}_e,\text{- for an electron},`$ (11)
$`\widehat{V}`$ $`=`$ $`{\displaystyle \frac{Ze^2}{r_p}}V_d^{N1}+\widehat{\mathrm{\Sigma }}_p,\text{- for a positron}.`$ (12)
The non-relativistic limit can be achieved by reducing the value of $`\alpha `$ in (9) to $`\alpha =0`$.
The relativistic energy shift in atoms with one external electron can also be estimated by the following equation
$`\mathrm{\Delta }_n={\displaystyle \frac{E_n}{\nu }}(Z\alpha )^2\left[{\displaystyle \frac{1}{j+1/2}}C(Z,j,l)\right],`$ (13)
where $`E_n`$ is the energy of an external electron, $`\nu `$ is the effective principal quantum number ($`E_n=0.5/\nu ^2`$ a.u.). The coefficient $`C(Z,j,l)`$ accounts for many-body effects. Note that formula (13) is based on the specific expression for the electron density in the vicinity of the nucleus and therefore is not applicable for a positron.
## III Silver and gold negative ions
We calculated electron affinities of silver and gold atoms mostly to test the technique used for positron-atom binding. The calculation of a negative ion Ag<sup>-</sup> or Au<sup>-</sup> is a two-particle problem technically very similar to positron-atom binding. The effective Hamiltonian of the problem has a form similar to (2)
$`H_{\mathrm{eff}}^{\mathrm{CI}}`$ $`=`$ $`\widehat{h}_e(r_1)+\widehat{h}_e(r_2)+\widehat{h}_{ee},`$ (14)
$`\widehat{h}_{ee}`$ $`=`$ $`{\displaystyle \frac{e^2}{|𝐫_e𝐫_p|}}+\widehat{\mathrm{\Sigma }}_{ee},`$ (15)
where $`\widehat{\mathrm{\Sigma }}_{ee}`$ represents the screening of the Coulomb interaction between external electrons by core electrons (see Refs. for detailed discussion). Electron affinity is defined when an electron can form a bound state with an atom. In this case the difference between the energy of a neutral atom and the energy of a negative ion is called the electron affinity to this atom. Energies of Ag, Ag<sup>-</sup>, Au, Au<sup>-</sup> obtained in different approximations and corresponding electron affinities are presented in Table I together with experimental data. The energies are given with respect to the cores (Ag<sup>+</sup> and Au<sup>+</sup>). Like in the case of Cu<sup>-</sup> the accuracy of the Hartree-Fock approximation is very poor. The binding energies of the $`5s`$ electron in neutral Ag and the $`6s`$ electron in neutral Au are underestimated by about 21% and 23% respectively, while the negative ions are unbound. Inclusion of either core-valence correlations ($`\mathrm{\Sigma }`$) or valence-valence correlations (CI) does produce binding but the accuracy is still poor. Only when both these effects are included the accuracy for the electron affinities improves significantly becoming 20% for Ag<sup>-</sup> and 11% for Au<sup>-</sup>. Further improvement can be achieved by introducing numerical factors before $`\widehat{\mathrm{\Sigma }}_e`$ to fit the lowest $`s,p`$ and $`d`$ energy levels of the neutral atoms. These factors simulate the effect of higher-order correlations. Their values are $`f_s=0.88`$, $`f_p=0.97`$, $`f_d=1.08`$ for the Ag atom and $`f_s=0.81`$, $`f_p=1`$, $`f_d=1.04`$ for the Au atom in the $`s`$, $`p`$ and $`d`$ channels, respectively. As is evident from Table I, the fitting of the energies of neutral atoms also significantly improves electron affinities. It is natural to assume that the same procedure should work equally well for the positron-atom problem.
Results of other calculations of the electron affinities of silver and gold are presented in Table II together with the experimental values.
## IV Positron binding to silver and gold and the role of relativistic effects
As for the case of copper we have performed calculations for two different cavity radii $`R=30a_0`$ and $`R=15a_0`$. For a smaller radius convergence with respect to the number of single-particle basis states is fast. However, the effect of the cavity on the converged energy is large. For a larger cavity radius, convergence is slower and the effect of the cavity on the energy is small. When the energy shift caused by the finite cavity radius is taken into account both calculations come to the same value of the positron binding energy. Table III illustrates the convergence of the calculated energies of $`e^+`$Ag and $`e^+`$Au with respect to the maximum value of the angular momentum of single-particle orbitals. Energies presented in the table are two-particle energies (in a.u.) with respect to the energies of Ag<sup>+</sup> and Au<sup>+</sup>. The number of radial orbitals $`n`$ in each partial wave is fixed at $`n=16`$. Fig. 1 shows the convergence of the calculated energy with respect to $`n`$ when maximum momentum of the single-particle orbitals was fixed at $`L=10`$. The cavity radius in both cases was $`R=30a_0`$. Table III and Fig. 1 show that even for a larger cavity radius, convergence was clearly achieved. Table III also shows the convergence in different approximations, namely with and without core-valence correlations ($`\mathrm{\Sigma }`$). One can see that while inclusion of $`\mathrm{\Sigma }`$ does shift the energy, the convergence is not affected.
Table IV shows how positron binding by silver and gold is formed in different approximations. This table is very similar to Table I for the negative ions except there is no RHF approximation for the positron binding. Indeed, the RHF approximation for the negative ions means a single-configuration approximation: $`5s^2`$ for Ag<sup>-</sup> and $`6s^2`$ for Au<sup>-</sup>. These configurations strongly dominate in the two-electron wave function of the negative ions even when a large number of configurations are mixed to ensure convergence. In contrast, no single configuration strongly dominates in the positron binding problem. Therefore we present our results in Table IV starting from the standard CI approximation. In this approximation positron is bound to both silver and gold atoms. However, the inclusion of core-valence correlations through the introduction of the $`\mathrm{\Sigma }_e`$, $`\mathrm{\Sigma }_p`$ and $`\mathrm{\Sigma }_{ep}`$ operators shifts the energies significantly. In the case of gold, the $`e^+`$Au system becomes unbound when all core-valence correlations are included.
As was discussed in our previous paper the dominating factor affecting the accuracy of the calculations is higher-order correlations which mostly manifest themself via the value of the $`\mathrm{\Sigma }`$ operator. An introduction of the fitting parameters as described in the previous section can be considered as a way to simulate the effect of higher-order correlations. Also, the energy shift caused by the fitting can be considered as an estimation of the uncertainty of the calculations. This shift is 0.00240 a.u. in the case of silver and 0.00023 a.u. in the case of gold (see Table IV). Note that these values are considerably smaller than energy shifts for the silver and gold negative ions (0.00854 a.u. and 0.00921 a.u. respectively, see Table I). This is because of the cancellation of the effects of the variation of $`\mathrm{\Sigma }_e`$ and $`\mathrm{\Sigma }_p`$. In particular, for gold it is accidentally very small. One can see that even if the value of 0.00240 a.u. is adopted as an upper limit of the uncertainty of the calculations, the $`e^+`$Ag system remains bound while the $`e^+`$Au system remains unbound. However, the actual accuracy might be even higher. We saw that the fitting procedure significantly improves the accuracy of the calculations for the silver and gold negative ions. It is natural to assume that the same procedure works equally well for the positron binding problem. The final result for the energy of positron binding by the silver atom as presented in Table IV is 0.00434 a.u. This result does not include the effect of the finite cavity size. When this effect is taken into account, by means of the procedure described in Ref. , the binding energy becomes 0.00452 a.u. or 123 meV. If we adopt the value of 0.00240 a.u as an estimation of the uncertainty of the result, then the accuracy we can claim is about 20%.
The calculation of the positron binding by copper , silver and gold reveal an interesting trend. All three atoms have very similar electron structure. However the positron binding energy for silver (123 meV) is considerably smaller than that for copper (170 meV ) while gold atoms cannot bind positrons at all. We believe that this trend is caused by relativistic effects. An argument that the positron is always non-relativistic does not look very convincing because electrons also contribute to the binding energy. Relativistic effects are large for heavy atoms and electron contributions to the positron binding energy could be very different in the relativistic and non-relativistic limits. Indeed, we demonstrated in Ref. that the relativistic energy shift considerably changes the values of the transition frequencies in Hg<sup>+</sup> ion and sometimes even changes the order of the energy levels. If we use formula (13) with the contribution of the many-body effects $`C=0.6`$, as suggested in Ref. , to estimate the relativistic energy shift for neutral Au then the result is -0.037 a.u. This is about an order of magnitude larger than the energy difference between Au and $`e^+`$Au. If the relativistic energy shift in $`e^+`$Au is different from that in Au then the positron binding energy may be strongly affected.
To study the role of the relativistic effects in positron binding in more detail we performed the calculations for Ag, Ag<sup>-</sup>, $`e^+`$Ag, Au, Au<sup>-</sup> and $`e^+`$Au in the relativistic and non-relativistic limits. The latter corresponds to the zero value of the fine structure constant $`\alpha `$ (see Section II). The results are presented in Table V. One can see that the actual relativistic energy shift for neutral Au is even bigger than is suggested by formula (13) with $`C=0.6`$. The shift is 0.0805 a.u. which corresponds to $`C=0.08`$. Formula (13) with $`C=0.08`$ also reproduces the relativistic energy shift for neutral Ag. The relativistic energy shift for an atom with a positron is of the same order of magnitude but a little different in value. This difference turned out to be enough to affect the positron binding energy significantly. In particular, the $`e^+`$Au system which is unbound in relativistic calculations becomes bound in the non-relativistic limit with binding energy 0.0080 a.u or 218 meV. In the case of silver, the positron binding energy is considerably higher in the non-relativistic limit. It is 0.0073 a.u. or 199 meV. It is interesting to compare this value with the value of 150 meV obtained by Mitroy and Ryzhikh using the non-relativistic stochastic variational method . Since the convergence was achieved in both calculations the remaining difference should probably be attributed to the different treatment of the core-valence correlations. We use many-body perturbation theory for an accurate calculation of the $`\mathrm{\Sigma }`$ operator which accounts for these correlations. Mitroy and Ryzhikh use an approximate semi-empirical expression for the $`\mathrm{\Sigma }`$ operator which is based on its long-range asymptotic behavior.
Note that the relativistic energy shift for negative ions is also large. However electron affinities are less affected. This is because electron affinities are many times larger than positron binding energies and therefore less sensitive to the energy shift. Apart from that there is a strong cancellation between relativistic energy shifts in the negative ion and neutral atom. This means in particular that the calculation of the electron affinities cannot serve as a test of a non-relativistic method chosen for the positron binding problem. However, it is still a good test of the relativistic calculations. Note also that our calculated relativistic energy shifts for neutral and negative silver and gold are in very good agreement with calculations performed by Schwerdtfeger and Bowmaker by means of relativistic and non-relativistic versions of the quadratic configuration interaction method (see Table VI and Ref. ).
The authors are grateful to G. F. Gribakin for many useful discussions. |
warning/0001/hep-ph0001014.html | ar5iv | text | # 1. Introduction
## 1. Introduction
Under the assumption of CPT invariance, the observed CP violation in neutral K decays demonstrates T violation in weak decays. However, T violation and CP violation are different physics concepts. CP violation requires the partial rate difference of the particle and its antiparticle while T violation needs the partial difference between the decay process and its time reversed process. The evidence for T violation based on CP violation and CPT invariance is indirect. The test of T violation should be done independently of CP violation.
Recently, CPLEAR collaboration gives the first direct observation of T violation in the difference of probability between $`K^0\overline{K^0}`$ and $`\overline{K^0}K^0`$ in the limit of CPT symmetry and the validity of $`\mathrm{\Delta }S=\mathrm{\Delta }Q`$ rule . Moreover, KTeV observed another evidence of T violation in the planar-angle asymmetry in the $`K_L\pi ^+\pi ^{}e^+e^{}`$ decay . Although the validity of T violation in these two decays was questioned by , we expect the future experiments can exclude some questions about T violation. Up to now, the study of T violation is mainly in K decays. It is well-known that B meson decays can provide another good place for testing CP violation. If the CPT invariance holds, T violation should be exactly equal to CP violation. Since CP violation in neutral B decays can be large ($`𝒪(1)`$), the large T violation may happen in B decays. So it is necessary to study T violation in B decays.
For the weak decay process $`if`$, T violation is defined by $`\mathrm{\Delta }_T\frac{\mathrm{\Gamma }(f^{}i^{})\mathrm{\Gamma }(if)}{\mathrm{\Gamma }(f^{}i^{})+\mathrm{\Gamma }(if)}`$ where $`f^{}i^{}`$ is the reversed process and the prime denotes the reverse of spin and momentum. In general, it is difficult to implement the reversed weak decay. One reason is that it is impossible to set up the initial condition such as in nuclear beta decay. The other reason lies in that the weak decay is so weak that it is unable to extract the reversed weak decay from the strong and electromagnetic backgrounds. The reversed reaction of $`B^0\pi ^+\pi ^{}`$ is such an example. The only exception is the neutral particle oscillation induced by weak interaction such as $`K^0\overline{K^0}`$ and $`B^0\overline{B^0}`$. We shall discuss T violation in $`B^0\overline{B^0}`$ oscillation without or with the assumption CPT invariance in this paper. For the latter case, T violation in $`B^0\overline{B^0}`$ oscillation is predicted to be small in the Standard Model (SM).
Another way to observe T violation is to measure the T-odd correlations in the final states of weak decays. A T-odd correlation is one that changes sign under the reverse of all incoming and outgoing three momentum and polarization. One classic example is the triple T-odd correlation $`\sigma _𝐧(𝐤_𝐩\times 𝐤_𝐞)`$ of the nuclear beta decay where $`\sigma _𝐧`$ is the spin of the neutron and $`𝐤_𝐩`$, $`𝐤_𝐞`$ are the three momentum of proton and electron. Whether the T-odd observable is considered as T violation should be viewed with caution. There is a mimicry of T violation caused by final state interaction (often refers to strong interaction) even if the fundamental interactions is time reversal invariant. Wolfenstein calls it ”pseudo Time Reversal Violation (pseudo TRV)”. In , the authors prove that using the unitarity constraint and CPT invariance of final state interaction, the T-odd effect can be identified with a measurement of T violation if the final state interaction effects are small and negligible. According to , the correlation between the meson and lepton plane in $`K_L\pi ^+\pi ^{}e^+e^{}`$ decay is T-odd and thus violates time reversal symmetry. The same method is used in to study the angular distribution of $`BK^{}\pi ^+e^+e^{}`$ and $`B\pi ^{}\pi ^+e^+e^{}`$. Their results show that T violation is small. In the differential angular distribution of the decay $`BV_1V_2`$ where $`V_{1,2}`$ represent two vector mesons, the interference terms contain T-odd correlation contribution. If $`\varphi `$ is the angle between the decay planes of the two vectors, the angular correlations $`sin2\varphi `$ and $`sin\varphi `$ terms are T-odd. In this paper, we intend to study T violation in $`BVV`$ decays using the angular distribution analysis. The final state interaction effects is also taken into account in the T-odd observable.
## 2. T violation in $`B^0\overline{B^0}`$ oscillation
The $`\mathrm{\Delta }B=\pm 2`$ weak decay via Box diagram causes the mixing between $`B^0`$ and $`\overline{B^0}`$. The physical states are the superposition of the flavor states $`|B^0>`$ and $`|\overline{B^0}>`$. The two mass eigenstates in neutral $`B_d`$ system can be generally written by
$`|B_1>={\displaystyle \frac{1}{\sqrt{|p_1^2|+|q_1|^2}}}[p_1|B^0>+q_1|\overline{B^0}>]`$
$`|B_2>={\displaystyle \frac{1}{\sqrt{|p_2^2|+|q_2|^2}}}[p_2|B^0>q_2|\overline{B^0}>]`$ (1)
In the neutral K system, the mixing parameter $`p_i`$, $`q_i`$ are usually represented by small parameters of $`ϵ`$, $`\mathrm{\Delta }`$. This parameterization method is not suitable to apply in neutral B system because CP and T violation in it is predicted to be large in CKM model. We take the exponential parameterization as given in . Thus the mixing parameters $`p_i`$, $`q_i`$ are related by
$`{\displaystyle \frac{q_1}{p_1}}=tg{\displaystyle \frac{\theta }{2}}e^{i\varphi },{\displaystyle \frac{q_2}{p_2}}=ctg{\displaystyle \frac{\theta }{2}}e^{i\varphi }`$ (2)
where $`\theta `$ and $`\varphi `$ are complex phases in general. According to , T violation requires $`\varphi 0`$.
The time evolution of the initially $`|B^0>`$ or $`|\overline{B^0}>`$ after a proper time $`t`$ is
$`|B^0(t)>=g_+(t)|B^0>+\overline{g}_+(t)|\overline{B^0}>`$
$`|\overline{B^0}(t)>=g_{}(t)|\overline{B^0}>+\overline{g}_{}(t)|B^0>`$ (3)
where
$`g_{\pm (t)}`$ $`=`$ $`e^{im_Bt\frac{1}{2}\mathrm{\Gamma }_Bt}[ch({\displaystyle \frac{ixy}{2}}\mathrm{\Gamma }_Bt)\pm cos\theta sh({\displaystyle \frac{ixy}{2}}\mathrm{\Gamma }_Bt)]`$
$`\overline{g}_\pm (t)`$ $`=`$ $`e^{im_Bt\frac{1}{2}\mathrm{\Gamma }_Bt}sin\theta e^{\pm i\varphi }sh({\displaystyle \frac{ixy}{2}}\mathrm{\Gamma }_Bt)`$ (4)
with $`x\frac{\mathrm{\Delta }m_B}{\mathrm{\Gamma }_B},y\frac{\mathrm{\Delta }\mathrm{\Gamma }_B}{2\mathrm{\Gamma }_B}.`$
The T violation in $`B^0\overline{B^0}`$ oscillation is defined as
$$A_T(t)\frac{P_{B^0(t)\overline{B^0}}P_{\overline{B^0}(t)B^0}}{P_{B^0(t)\overline{B^0}}+P_{\overline{B^0}(t)B^0}}=\frac{|e^{i\varphi }|^2|e^{i\varphi }|^2}{|e^{i\varphi }|^2+|e^{i\varphi }|^2}=2\mathrm{I}\mathrm{m}\varphi $$
(5)
From Eq.(5), $`A_T(t)`$ is independent of $`t`$. It is a constant number. This is the same as the case in the CPLEAR experiment. Moreover $`A_T(t)`$ is not related with the CPT violation parameter $`\theta `$. Note that $`A_T(t)`$ is proportional to $`Im\varphi `$, the imaginary part of the mixing parameter $`\varphi `$. As it will be seen later, T violation in the interference of the $`B^0\overline{B^0}`$ mixing and the decay amplitude requires $`Re\varphi 0`$.
The experimental test of T violation $`A_T(t)`$ can be determined in the semileptonic decay and the same-sign dileptonic ratios of B decays through
$`A_T(t)={\displaystyle \frac{\mathrm{\Gamma }(B^0(t)\overline{X}l^{}\nu )\mathrm{\Gamma }(\overline{B^0}(t)Xl^+\nu )}{\mathrm{\Gamma }(B^0(t)\overline{X}l^{}\nu )+\mathrm{\Gamma }(\overline{B^0}(t)Xl^+\nu )}}={\displaystyle \frac{N^{++}N^{}}{N^{++}+N^{}}}=2\mathrm{I}\mathrm{m}\varphi `$ (6)
where $`N^{++}`$, $`N^{}`$ are the same-sign dilepton events. In Eq.(6), The validity of $`\mathrm{\Delta }B=\mathrm{\Delta }Q`$ rule is assumed.
We next turn to discuss the SM expectation of the $`A_T`$. In SM, CPT is invariant, and the origin of CP and T violation lies in the nonzero complex phase in CKM matrix. The mixing parameter $`\theta `$ in Eq.(2) will be equal to $`\frac{\pi }{2}`$. Thus, $`\frac{q_1}{p_1}=\frac{q_2}{p_2}=\frac{q}{p}`$. According to ,
$`|{\displaystyle \frac{q}{p}}|1=|e^{i\varphi }|1={\displaystyle \frac{1}{2}}Im{\displaystyle \frac{\mathrm{\Gamma }_{12}}{M_{12}}}𝒪(10^3)`$ (7)
Thus
$`A_T(t)Im{\displaystyle \frac{\mathrm{\Gamma }_{12}}{M_{12}}}𝒪(10^3)`$ (8)
The above estimate is based on the assumption that the box diagram with a cut is appropriate to calculate $`\mathrm{\Gamma }_{12}`$. The uncertainty from the use of quark diagram to describe $`\mathrm{\Gamma }_{12}`$ could be a factor of 2-3.
## 3. T violation in the angular distribution of $`BVV`$ decay
As discussed in the Introduction, another way to observe T violation is through the T-odd correlation in the final states. Nonleptonic B decays play important role in exploring CP violation such as the decays $`BJ/\psi K_S`$, $`\pi \pi `$, $`\pi K`$, etc. Unlike CP violation, there is no T-odd correlation in $`BPP`$ and $`BVP`$ decay. Because the decay amplitude contains only T-even term: the momentum square for $`BPP`$ decay; and the product of momentum and polarization vector for $`BVP`$ decay. Both of these terms are invariant under time reversal. In $`BVV`$ decays, the angular correlation between the decay planes of two vectors contains T-odd terms thus provides place to search T violation. We take $`BJ/\psi K^{}`$ as an example to discuss the T violation in $`BVV`$ decays.
The differential decay distribution for $`BK^{}J/\psi (K\pi )(l^+l^{})`$ is :
$`{\displaystyle \frac{d^3\mathrm{\Gamma }}{dcos\theta _1dcos\theta _2d\varphi }}`$ $`=`$ $`{\displaystyle \frac{|\stackrel{}{p}|}{16\pi ^2m_B^2}}{\displaystyle \frac{9}{8}}\{{\displaystyle \frac{1}{4}}sin^2\theta _1(1+cos^2\theta _2)(|H_{+1}|^2+|H_1|^2)+cos^2\theta _1sin^2\theta _2|H_0|^2`$
$`{\displaystyle \frac{1}{2}}sin^2\theta _1sin^2\theta _2[cos2\varphi Re(H_{+1}H_1^{})sin2\varphi Im(H_{+1}H_1^{})]`$
$`{\displaystyle \frac{1}{4}}sin2\theta _1sin2\theta _2[cos\varphi Re(H_{+1}H_0^{}+H_1H_0^{})sin\varphi Im(H_{+1}H_0^{}H_1H_0^{})]\}`$
where $`\theta _1`$ is the polar angle of the $`K`$ momentum in the rest frame of the $`K^{}`$ meson with respect to the helicity axis of $`K^{}`$ meson (the negative of the the direction of the $`J/\psi `$ in $`K^{}`$ rest frame) and similarly $`\theta _2`$ is the polar angle of the positive lepton $`l^+`$ ($`e^+`$ or $`\mu ^+`$) momentum in the rest frame of the $`J/\psi `$ with respect to the helicity axis of $`J/\psi `$; $`\varphi `$ is the angle between the planes of the two decays of $`K^{}K\pi `$ and $`J/\psi l^+l^{}`$. In eq.(9), $`\stackrel{}{p}`$ is the three momentum of the vector $`K^{}`$; $`H_i`$ are the helicity amplitude defined in . The angle correlations $`sin2\varphi `$ and $`sin\varphi `$ are T-odd. To confirm this, define the unit vector $`\widehat{𝐩}\frac{\stackrel{}{p_K^{}}}{|\stackrel{}{p_K^{}}|}`$. Thus,
$`sin\varphi `$ $`=`$ $`({\displaystyle \frac{\stackrel{}{p_K}\times \stackrel{}{p_\pi }}{|\stackrel{}{p_K}\times \stackrel{}{p_\pi }|}})\times ({\displaystyle \frac{\stackrel{}{p_{l^+}}\times \stackrel{}{p_l^{}}}{|\stackrel{}{p_{l^+}}\times \stackrel{}{p_l^{}}|}})\widehat{𝐩}`$
$`sin2\varphi `$ $`=`$ $`2({\displaystyle \frac{\stackrel{}{p_K}\times \stackrel{}{p_\pi }}{|\stackrel{}{p_K}\times \stackrel{}{p_\pi }|}})\times ({\displaystyle \frac{\stackrel{}{p_{l^+}}\times \stackrel{}{p_l^{}}}{|\stackrel{}{p_{l^+}}\times \stackrel{}{p_l^{}}|}})\widehat{𝐩}({\displaystyle \frac{\stackrel{}{p_K}\times \stackrel{}{p_\pi }}{|\stackrel{}{p_K}\times \stackrel{}{p_\pi }|}})({\displaystyle \frac{\stackrel{}{p_{l^+}}\times \stackrel{}{p_l^{}}}{|\stackrel{}{p_{l^+}}\times \stackrel{}{p_l^{}}|}})`$ (10)
From the above equation, $`sin2\varphi `$ and $`sin\varphi `$ contain 9 and 5 momentum vectors in the products respectively. Under the time reversal transformation, they change their signs. All these momentums are defined in the rest frame of $`B`$ meson.
Another form of the angular distribution based on the transversity variable is given in . In that form, the CP-even and odd and T-even and old component is obvious. Both the two froms are principally the same except for the adoption of different variable.
The integration over angles $`\theta _1`$ and $`\theta _2`$ yields the $`\varphi `$ angle distribution
$`{\displaystyle \frac{d\mathrm{\Gamma }}{d\varphi }}={\displaystyle \frac{|\stackrel{}{p}|}{16\pi ^2m_B^2}}\{|H_{+1}|^2+|H_1|^2+|H_0|^2cos2\varphi Re(H_{+1}H_1^{})+sin2\varphi Im(H_{+1}H_1^{})\}`$ (11)
From Eq.(11), only one T-odd $`sin2\varphi `$ term is left when integrating over angles $`\theta _1`$ and $`\theta _2`$. The other T-odd $`sin\varphi `$ term can be extracted from the full three-angle distribution or the difference of $`\varphi `$ angle distribution between the same hemisphere events (e.g. $`0<\theta _1,\theta _2<\frac{\pi }{2}`$) and the opposite hemisphere events (e.g. $`0<\theta _1<\frac{\pi }{2},\frac{\pi }{2}<\theta _2<\pi `$). For this case, the full angle distribution is required to be known from the experiment. In this paper, we restrict our discussion in the single angle $`\varphi `$ distribution given by Eq.(11) because it is easier to treat in the experiment.
So, T violation is given by
$`\mathrm{\Delta }_T={\displaystyle \frac{(_0^{\frac{\pi }{2}}_{\frac{\pi }{2}}^\pi +_\pi ^{\frac{3\pi }{2}}_{\frac{3\pi }{2}}^{2\pi })\frac{d\mathrm{\Gamma }}{d\varphi }d\varphi }{(_0^{\frac{\pi }{2}}+_{\frac{\pi }{2}}^\pi +_\pi ^{\frac{3\pi }{2}}+_{\frac{3\pi }{2}}^{2\pi })\frac{d\mathrm{\Gamma }}{d\varphi }d\varphi }}={\displaystyle \frac{2}{\pi }}{\displaystyle \frac{Im(H_{+1}H_1^{})}{|H_{+1}|^2+|H_1|^2+|H_0|^2}}{\displaystyle \frac{2}{\pi }}\beta _2`$ (12)
Up to now, our analysis is model independent. In the remainder of the paper, we will restrict our discussion in the Standard Model. From Eq.(12), T violation observable $`\mathrm{\Delta }_T`$ is proportional to the angular correlation coefficient $`\beta _2`$. This relation is a general result of the decay $`BVV`$. The nonvanishing $`Im(H_+H_1^{})`$ is caused by weak CKM phases or strong final state interaction phases under the condition that they contribute differently to $`H_{+1}`$ and $`H_1`$. First, we discuss the case that the final state interaction is absent. In , the authors had systematically calculated all the $`BVV`$ decays. Their result shows that $`\beta _2`$ is very small. For most $`BVV`$ process, $`\beta _2`$ is less than $`10^4`$. In the special example of $`BK^{}J/\psi `$, the tree and the dominant QCD Penguin diagram have the same CKM phase, thus $`\beta _2`$ is nearly zero.
Second, we consider the nonvanishing $`\beta _2`$ caused only by strong final state interaction. Since the strong interaction is T invariant, the violation induced by final state interaction is not the true T violation but the mimicry of T violation.
Let us introduce
$`H_{||}={\displaystyle \frac{1}{\sqrt{2}}}(H_{+1}+H_1)`$
$`H_{}={\displaystyle \frac{1}{\sqrt{2}}}(H_{+1}H_1)`$ (13)
and define the strong phase difference $`\delta Arg(H_{||}H_{}^{})`$ then
$`\mathrm{\Delta }_T={\displaystyle \frac{2}{\pi }}{\displaystyle \frac{|H_{||}||H_{}|sin\delta }{|H_{||}|^2+|H_{}|^2+|H_0|^2}}`$ (14)
The fact that the T violation mimicry appears to be proportional to $`sin\delta `$ was revealed long time ago (see ). Here we again find this particular characteristic in $`BVV`$ decays. Up to now, definite quantitative analysis of final state interaction has not been accessible yet. The concrete information about strong phase is unknown. But, based on some phenomenological consideration, the final state interaction effects in $`BK^{}J/\psi `$ is estimated to be small. The small width of $`J/\psi `$ makes the strong coupling between $`J/\psi `$ and strong states be small. Moreover, there are few channels with large branching ratios that can transform into the final state $`K^{}J/\psi `$ through strong interaction. From above, it seems that T violation in the angular distribution of $`BK^{}J/\psi `$ is a small effect.
There are three types of CP violation in neutral B system. CP violations in decay amplitude and mixing are small. The most important type is the CP violation in the interference of mixing and decay. This type is promising to give large CP violation. T violation in $`B^0\overline{B^0}`$ oscillation is due to $`B^0\overline{B^0}`$ mixing. T violation in the above discussion of the $`BVV`$ decay which requires weak phases contribute differently to helicity amplitudes belongs to the T violation in decay. Both of them are small. Next, we intend to seek for large T violation in the interference of mixing and decay.
In the Standard Model, the time evolution of $`B^0`$ and $`\overline{B^0}`$ is obtained from Eq.(3)
$`|B^0(t)>=f_+(t)|B^0>+{\displaystyle \frac{q}{p}}f_+(t)|\overline{B^0}>`$
$`|\overline{B^0}(t)>={\displaystyle \frac{p}{q}}f_{}(t)|B^0>+f_+(t)|\overline{B^0}>`$ (15)
where $`f_+(t)=e^{im_Bt\frac{\mathrm{\Gamma }_Bt}{2}}cos\frac{\mathrm{\Delta }m_Bt}{2}`$, $`f_{}(t)=e^{im_Bt\frac{\mathrm{\Gamma }_Bt}{2}}isin\frac{\mathrm{\Delta }m_Bt}{2}`$, $`\frac{q}{p}=e^{i2\beta }`$, and $`\beta `$ is the angle of the unitarity triangle. We have neglected the width difference of two mass eigenstates.
If $`K^0`$ in the decay $`B^0K^0J/\psi `$ is observed to decay to CP eigenstates $`\pi ^0K_S`$, then angular distribution analysis gives the time dependent T violation as
$`\mathrm{\Delta }_T(t)`$ $`=`$ $`{\displaystyle \frac{2}{\pi }}{\displaystyle \frac{Im(H_{||}(t)H_{}^{}(t))}{|H_{||}(t)|^2+|H_{}(t)|^2+|H_0(t)|^2}}`$
$`=`$ $`{\displaystyle \frac{2}{\pi }}{\displaystyle \frac{|H_{||}(0)||H_{}(0)|[sin\delta cos\mathrm{\Delta }m_Bt+cos\delta cos2\beta sin\mathrm{\Delta }m_Bt]e^{\mathrm{\Gamma }_Bt}}{[|H_{||}(0)|^2+|H_{}(0)|^2+|H_0|^2(0)+sin2\beta sin\mathrm{\Delta }m_Bt(|H_{||}(0)|^2+|H_0(0)|^2|H_{}(0)|^2)]e^{\mathrm{\Gamma }_Bt}}}`$
The above time dependent T violation has two contributions. The first term is the mimicry of T violation induced by final state interaction. The second term which contains $`cos2\beta `$ in the absence of final state interaction is due to the interference of mixing and decay. In this case, the small final state interaction effects can be neglected.
The time integrated T violation obtained from Eq.(16) is
$`D_T`$ $`=`$ $`{\displaystyle \frac{2}{\pi }}{\displaystyle \frac{_0^{\mathrm{}}𝑑tIm(H_{||}(t)H_{}^{}(t))}{_0^{\mathrm{}}𝑑t[|H_{||}(t)|^2+|H_{}(t)|^2+|H_0(t)|^2]}}`$
$`=`$ $`{\displaystyle \frac{2}{\pi }}{\displaystyle \frac{|H_{||}(0)||H_{}(0)|[\frac{sin\delta }{1+x^2}+cos\delta cos2\beta \frac{x}{1+x^2}]}{|H_{||}(0)|^2+|H_{}(0)|^2+|H_0|^2(0)+sin2\beta \frac{x}{1+x^2}(|H_{||}(0)|^2+|H_0(0)|^2|H_{}(0)|^2)}}`$
where $`x\frac{\mathrm{\Delta }m_B}{\mathrm{\Gamma }_B}=0.7`$ is obtained from PDG98 .
In order to estimate the time integrated T violation, we use $`sin2\beta =0.5`$ and the parameters given in which based on the models of BSW, Soares and Cheng . Table 1 gives the results of time integrated T violation in $`BK^{}(\pi ^0K_S)J/\psi (l^+l^{})`$ with different models in the absence of final state interaction. From Table 1, one can see that different models give the T violation range from 0.04 to 0.07.
## 4. Conclusion
In this paper, we present a study of T violation in $`B^0\overline{B^0}`$ oscillation and $`BVV`$ decays. T violation in B decays opens another way to test Standard Model and the origin of CP/T violation. In $`B^0\overline{B^0}`$ oscillation, T violation induced by $`B^0\overline{B^0}`$ mixing is about the order of $`10^3`$. This tiny effects is possible to observe in semileptonic decay and dileptonic decay at B-factory and LHC-B. If a large effect is found, it will be new physics beyond the Standard Model. T violation in decay of $`BVV`$ from the interference term $`\beta _2`$ with different weak phases that contribute to helicity amplitudes is small, and this effect can not be extracted from the mimicry induced by final state interaction. Via interference of $`B^0\overline{B^0}`$ mixing and decay, the time integrated T violation in $`B^0K^0J/\psi (K_S\pi ^0)(l^+l^{})`$ decay can reach 4-7% which is experimentally accessible. In this case final state interaction effects can be neglected.
From the CPT theorem, T violation should be exactly equal to CP violation. CP violation in neutral B decays can be as large as $`𝒪(1)`$, while T violation we had found is at the 10% level of CP violation. So the question that how and where to find large T violation in B decays arises.
Note added: After finishing this paper, we became aware that the T-odd correlation in $`BVV`$ decays was pointed out in hep-ph/9911338 . However, the physics motivation of two papers are different.
## Acknowledgment
This work is supported in part by National Natural Science Foundation of China and the Grant of State Commission of Science and Technology of China. |
warning/0001/cond-mat0001440.html | ar5iv | text | # The critical equation of state of three-dimensional 𝑋𝑌 systems
## I Introduction
In the theory of critical phenomena, continuous phase transitions can be classified into universality classes determined only by a few basic properties characterizing the system, such as the space dimensionality, the range of interaction, the number of components and the symmetry of the order parameter. The renormalization-group theory predicts that, within a given universality class, the critical exponents and the scaling functions are the same for all systems. Here we consider the $`XY`$ universality class, which is characterized by a two-component order parameter and effective short-range interactions. The lattice spin model described by the Hamiltonian
$$_L=J\underset{<ij>}{}\stackrel{}{s}_i\stackrel{}{s}_j+\underset{i}{}\stackrel{}{h}_i\stackrel{}{s}_i,$$
(1)
where $`\stackrel{}{s}_i`$ is a two-component spin satisfying $`\stackrel{}{s}_i\stackrel{}{s}_i=1`$, is one of the systems belonging to the $`XY`$ universality class. It may be viewed as a magnetic system with easy-plane anisotropy, in which the magnetization plays the role of order parameter and the spins are coupled to an external magnetic field $`h`$.
The superfluid transition of <sup>4</sup>He, occurring along the $`\lambda `$-line $`T_\lambda (P)`$ (where $`P`$ is the pressure), belongs to the three-dimensional $`XY`$ universality class. Its order parameter is related to the complex quantum amplitude of helium atoms. It provides an exceptional opportunity for an experimental test of the renormalization-group predictions, thanks to the weakness of the singularity in the compressibility of the fluid and to the purity of the sample. Moreover, experiments in a microgravity environment lead to a reduction of the gravity-induced broadening of the transition. Recently a Space Shuttle experiment performed a very precise measurement of the heat capacity of liquid helium to within 2 nK from the $`\lambda `$-transition obtaining an extremely accurate estimate of the exponent $`\alpha `$ and of the ratio $`A^+/A^{}`$ of the specific-heat amplitudes:
$$\alpha =0.01285(38),A^+/A^{}=1.054(1).$$
(2)
These results represent a challenge for theorists because the accuracy of the test of the renormalization-group prediction is now limited by the precision of the theoretical calculations. We mention the best available theoretical estimates for $`\alpha `$: $`\alpha =0.0150(17)`$ obtained using high-temperature expansion techniques , $`\alpha =0.0169(33)`$ from Monte Carlo simulations using finite-size scaling techniques , $`\alpha =0.011(4)`$ from field theory . The close agreement with the experimental data clearly supports the standard renormalization-group description of the $`\lambda `$-transition .
In this paper we address the problem of determining the critical equation of state characterizing the $`XY`$ universality class. The critical equation of state relates the thermodynamical quantities in the neighborhood of the critical temperature, in both phases. It is usually written in the form (see e.g. Ref. )
$$\stackrel{}{H}=\stackrel{}{M}M^{\delta 1}f(x),xtM^{1/\beta },$$
(3)
where $`f(x)`$ is a universal scaling function (normalized in such a way that $`f(1)=0`$ and $`f(0)=1`$). The universal ratios of amplitudes involving quantities defined at zero-momentum (i.e. integrated in the volume), such as specific heat, magnetic susceptibility, etc…, can be obtained from the scaling function $`f(x)`$.
It should be noted that, for the $`\lambda `$-transition in <sup>4</sup>He, Eq. (3) is not directly related to the conventional equation of state that relates temperature and pressure. Moreover, in this case, the field $`\stackrel{}{H}`$ does not correspond to an experimentally accessible external field, so that the function appearing in Eq. (3) cannot be determined directly in experiments. The physically interesting quantities are universal amplitude ratios of quantities formally defined at zero external field.
As our starting point for the determination of the critical equation of state, we compute the first few nontrivial coefficients of the small-field expansion of the effective potential (Helmholtz free energy) in the high-temperature phase. For this purpose, we analyze the high-temperature expansion of an improved lattice Hamiltonian with suppressed leading scaling corrections . If the leading non-analytic scaling corrections are no longer present, one expects a faster convergence, and therefore an improved high-temperature expansion (IHT) whose analysis leads to more precise and reliable estimates. We consider a simple cubic lattice and the $`\varphi ^4`$ Hamiltonian
$$=\beta \underset{x,y}{}\stackrel{}{\varphi }_x\stackrel{}{\varphi }_y+\underset{x}{}\left[\stackrel{}{\varphi }_x^2+\lambda (\stackrel{}{\varphi }_x^21)^2\right],$$
(4)
where $`x,y`$ labels a lattice link, and $`\stackrel{}{\varphi }_x`$ is a real two-component vector defined on lattice sites. The value of $`\lambda `$ at which the leading corrections vanish has been determined by Monte Carlo simulations using finite-size techniques , obtaining $`\lambda ^{}=2.10(6)`$. In Ref. we have already considered the high-temperature expansion (to 20th order) of the improved $`\varphi ^4`$ Hamiltonian (4) for the determination of the critical exponents, achieving a substantial improvement with respect to previous theoretical estimates. IHT expansions have also been considered for Ising-like systems , obtaining accurate determinations of the critical exponents, of the small-field expansion of the effective potential, and of the small-momentum behavior of the two-point function.
We use the small-field expansion of the effective potential in the high-temperature phase to determine approximate representations of the equation of state that are valid in the whole critical regime. To reach the coexistence curve ($`t<0`$) from the high-temperature phase ($`t>0`$), an analytic continuation in the complex $`t`$-plane is required. For this purpose we use parametric representations , which implement in a rather simple way the known analytic properties of the equation of state (Griffith’s analyticity). This approach was successfully applied to the Ising model, for which one can construct a systematic approximation scheme based on polynomial parametric representations and on a global stationarity condition . This leads to an accurate determination of the critical equation of state and of the universal ratios of amplitudes that can be extracted from it . $`XY`$ systems, in which the phase transition is related to the breaking of the continuous symmetry O(2), present a new important feature with respect to Ising-like systems: the Goldstone singularities at the coexistence curve. General arguments predict that at the coexistence curve ($`t<0`$ and $`H0`$) the transverse and longitudinal magnetic susceptibilities behave respectively as
$$\chi _T=\frac{M}{H},\chi _L=\frac{M}{H}H^{d/22}.$$
(5)
In our analysis we will consider polynomial parametric representations that have the correct singular behavior at the coexistence curve.
The most important result of the present paper, from the point of view of comparison with experiments, is the specific-heat amplitude ratio; our final estimate
$$A^+/A^{}=1.055(3)$$
(6)
is perfectly consistent with the experimental estimate (2), although not as precise.
The paper is organized as follows. In Sec. II we study the small-field expansion of the effective potential (Helmholtz free energy). We present estimates of the first few nontrivial coefficients of such expansion obtained by analyzing the IHT series. The results are then compared with other theoretical estimates. In Sec. III, using as input parameters the critical exponents and the known coefficients of the small-field expansion of the effective potential, we construct approximate representations of the critical equation of state. We obtain new estimates for many universal amplitude ratios. These results are then compared with experimental and other theoretical estimates.
## II The effective potential in the high-temperature phase
### A Small-field expansion of the effective potential in the high-temperature phase
The effective potential (Helmholtz free energy) is related to the (Gibbs) free energy of the model. If $`\stackrel{}{M}\stackrel{}{\varphi }`$ is the magnetization and $`\stackrel{}{H}`$ the magnetic field, one defines
$$(M)=\stackrel{}{M}\stackrel{}{H}\frac{1}{V}\mathrm{log}Z(H),$$
(7)
where $`Z(H)`$ is the partition function and the dependence on the temperature is always understood in the notation.
In the high-temperature phase the effective potential admits an expansion around $`M=0`$:
$$\mathrm{\Delta }(M)(0)=\underset{j=1}{\overset{\mathrm{}}{}}\frac{1}{(2j)!}a_{2j}M^{2j}.$$
(8)
This expansion can be rewritten in terms of a renormalized magnetization $`\phi `$
$$\mathrm{\Delta }=\frac{1}{2}m^2\phi ^2+\underset{j=2}{}m^{3j}\frac{1}{(2j)!}g_{2j}\phi ^{2j}$$
(9)
where
$$\phi ^2=\frac{\xi (t,H=0)^2M(t,H)^2}{\chi (t,H=0)},$$
(10)
$`t`$ is the reduced temperature, $`\chi `$ and $`\xi `$ are respectively the magnetic susceptibility and the second-moment correlation length
$`\chi ={\displaystyle \underset{x}{}}\varphi _\alpha (0)\varphi _\alpha (x),`$ (11)
$`\xi ={\displaystyle \frac{1}{6\chi }}{\displaystyle \underset{x}{}}x^2\varphi _\alpha (0)\varphi _\alpha (x),`$ (12)
and $`m1/\xi `$. In field theory $`\phi `$ is the expectation value of the zero-momentum renormalized field. The zero-momentum $`2j`$-point renormalized constants $`g_{2j}`$ approach universal constants (which we indicate with the same symbol) for $`t0`$. By performing a further rescaling
$$\phi =\frac{m^{1/2}}{\sqrt{g_4}}z$$
(13)
in Eq. (9), the free energy can be written as
$$\mathrm{\Delta }=\frac{m^3}{g_4}A(z),$$
(14)
where
$$A(z)=\frac{1}{2}z^2+\frac{1}{4!}z^4+\underset{j=3}{}\frac{1}{(2j)!}r_{2j}z^{2j},$$
(15)
and
$$r_{2j}=\frac{g_{2j}}{g_4^{j1}}j3.$$
(16)
One can show that $`zt^\beta M`$, and that the equation of state can be written in the form
$$Ht^{\beta \delta }\frac{A(z)}{z}.$$
(17)
### B Zero-momentum renormalized couplings by IHT expansion
To compute the high-temperature series of the four-point coupling $`g_4`$ and of the effective-potential parameters $`r_{2j}`$, we rewrite them in terms of the zero-momentum connected $`2j`$-point Green’s functions $`\chi _{2j}`$
$$\chi _{2j}=\underset{x_2,\mathrm{},x_{2j}}{}\varphi _{\alpha _1}(0)\varphi _{\alpha _1}(x_2)\mathrm{}\varphi _{\alpha _j}(x_{2j1})\varphi _{\alpha _j}(x_{2j})_c$$
(18)
($`\chi =\chi _2`$). For generic $`N`$-vector models we have
$$g_4=\frac{3N}{N+2}\frac{\chi _4}{\chi _2^2\xi ^3}$$
(19)
and
$`r_6=`$ $`10{\displaystyle \frac{5(N+2)}{3(N+4)}}{\displaystyle \frac{\chi _6\chi _2}{\chi _4^2}},`$ (20)
$`r_8=`$ $`280{\displaystyle \frac{280(N+2)}{3(N+4)}}{\displaystyle \frac{\chi _6\chi _2}{\chi _4^2}}+{\displaystyle \frac{35(N+2)^2}{9(N+4)(N+6)}}{\displaystyle \frac{\chi _8\chi _2^2}{\chi _4^3}},`$ (21)
$`r_{10}=`$ $`15400{\displaystyle \frac{7700(N+2)}{(N+4)}}{\displaystyle \frac{\chi _6\chi _2}{\chi _4^2}}+{\displaystyle \frac{350(N+2)^2}{(N+4)^2}}{\displaystyle \frac{\chi _6^2\chi _2^2}{\chi _4^4}}`$ (23)
$`+{\displaystyle \frac{1400(N+2)^2}{3(N+4)(N+6)}}{\displaystyle \frac{\chi _8\chi _2^2}{\chi _4^3}}{\displaystyle \frac{35(N+2)^3}{3(N+4)(N+6)(N+8)}}{\displaystyle \frac{\chi _{10}\chi _2^3}{\chi _4^4}}.`$
The formulae relevant for the $`XY`$ universality class are obtained setting $`N=2`$.
Using the $`\varphi ^4`$ lattice Hamiltonian (4), we have calculated $`\chi `$ and $`m_2_xx^2\varphi (0)\varphi (x)`$ to 20th order, $`\chi _4`$ to 18th order, $`\chi _6`$ to 17th order, $`\chi _8`$ to 16th order, and $`\chi _{10}`$ to 15th order, for generic values of $`\lambda `$. The IHT expansion, i.e. with suppressed leading scaling corrections, is achieved for $`\lambda =2.10(6)`$ . In Table I we report the series of $`m_2`$, $`\chi `$, $`\chi _4`$, $`\chi _6`$, $`\chi _8`$ and $`\chi _{10}`$ for $`\lambda =2.10`$. Using Eqs. (19) and (20) one can obtain the HT series necessary for the determination of $`g_4`$ and $`r_{2j}`$. We analyzed the series using the same procedure applied to the improved high-temperature expansions of Ising-like systems in Ref. . In order to estimate the fixed-point value of $`g_4`$ and $`r_{2j}`$, we considered Padé, Dlog-Padé and first-order integral approximants of the series in $`\beta `$ for $`\lambda =2.10`$, and evaluated them at $`\beta _c`$. We refer to Ref. for the details of the analysis. Our estimates are
$`g_4=21.05(3+3),`$ (24)
$`\overline{g}{\displaystyle \frac{5}{24\pi }}g_4=1.396(2+2),`$ (25)
$`r_6=1.951(11+3),`$ (26)
$`r_8=1.36(6+3).`$ (27)
We quote two errors: the first one is related to the spread of the approximants, while the second one gives the variation of the estimate when $`\lambda `$ varies between 2.04 and 2.16. In addition, we obtained a rough estimate of $`r_{10}`$, i.e. $`r_{10}=13(7)`$. For comparison, we anticipate that the analysis of the critical equation of state using approximate parametric representations will lead to the estimate $`r_{10}=10(3)`$. From the estimates of $`g_4^{}`$ and $`r_{2j}`$ one can obtain corresponding estimates for the zero-momentum renormalized couplings, $`g_{2j}=r_{2j}g_4^{j1}`$ with $`j>2`$.
Table II compares our results (denoted by IHT) with the estimates obtained using other approaches, such as the high-temperature expansion (HT) of the standard lattice spin model (1, field-theoretic methods based on the fixed-dimension $`d=3`$ $`g`$-expansion and on the $`ϵ`$-expansion . The fixed-dimension field-theoretic estimates of $`g_4`$ have been obtained from the zero of the Callan-Symanzik $`\beta `$-function, whose expansion is known to six loops . In the same framework $`g_6`$ and $`g_8`$ have been estimated from the analysis of the corresponding four- and three- loop series respectively . The authors of Ref. argue that the uncertainty on their estimate of $`g_6`$ is approximately 0.3%, while they consider their value for $`g_8`$ much less accurate. The $`ϵ`$-expansion estimates have been obtained from constrained analyses of the four-loop series of $`g_4`$ and the three-loop series of $`r_{2j}`$.
## III The critical equation of state
### A Analytic properties of the scaling equation of state
From the analysis of the IHT series we have obtained the first few non-trivial terms of the small-field expansion of the effective potential in the high-temperature phase. This provides corresponding information for the equation of state
$$Ht^{\beta \delta }F(z),$$
(28)
where $`zMt^\beta `$ and, using Eq. (17),
$$F(z)=\frac{A(z)}{z}=z+\frac{1}{6}z^3+\underset{m=3}{}F_{2m1}z^{2m1}$$
(29)
with
$$F_{2m1}=\frac{1}{(2m1)!}r_{2m}.$$
(30)
The function $`H(M,t)`$ representing the external field in the critical equation of state (28) satisfies Griffith’s analyticity: it is regular at $`M=0`$ for $`t>0`$ fixed and at $`t=0`$ for $`M>0`$ fixed. The first region corresponds to small $`z`$ in Eq. (28), while the second is related to large $`z`$, where $`F(z)`$ has an expansion of the form
$$F(z)=z^\delta \underset{n=0}{}F_n^{\mathrm{}}z^{n/\beta }.$$
(31)
To reach the coexistence curve, i.e. $`t<0`$ and $`H=0`$, one should perform an analytic continuation in the complex $`t`$-plane . The spontaneous magnetization is related to the complex zero $`z_0`$ of $`F(z)`$. Therefore, the description of the coexistence curve is related to the behavior of $`F(z)`$ in the neighbourhood of $`z_0`$.
### B Goldstone singularities at the coexistence curve
The physics of the broken phase of $`N`$-vector models (including $`XY`$ systems which correspond to $`N=2`$) is very different from that of the Ising model, because of the presence of Goldstone modes at the coexistence curve. The singularity of $`\chi _L`$ for $`t<0`$ and $`H0`$ is governed by the zero-temperature infrared-stable fixed point . This leads to the prediction
$$f(x)c_f(1+x)^{2/(d2)}\mathrm{for}x1,$$
(32)
where $`xtM^{1/\beta }`$ and $`f(x)`$ is the scaling function introduced in Eq. (3) (as usual, $`x=1`$ corresponds to the coexistence curve). This behavior at the coexistence curve has been verified in the framework of the large-$`N`$ expansion to $`O(1/N)`$ (i.e. next-to-leading order) .
The nature of the corrections to the behaviour (32) is less clear. Setting $`\omega =1+x`$ and $`y=HM^\delta `$, it has been conjectured that $`\omega `$ has the form of a double expansion in powers of $`y`$ and $`y^{(d2)/2}`$ near the coexistence curve , i.e. for $`y0`$
$$\omega 1+x=c_1y+c_2y^{1ϵ/2}+d_1y^2+d_2y^{2ϵ/2}+d_3y^{2ϵ}+\mathrm{}$$
(33)
where $`ϵ=4d`$. This expansion has been derived essentially from an $`ϵ`$-expansion analysis . Note that in three dimensions this conjecture predicts an expansion in powers of $`y^{1/2}`$, or equivalently an expansion of $`f(x)`$ in powers of $`\omega `$ for $`\omega 0`$.
The asymptotic expansion of the $`d`$-dimensional equation of state at the coexistence curve has been computed analytically in the framework of the large-$`N`$ expansion , using the $`O(1/N)`$ formulae reported in Ref. . It turns out that the expansion (33) does not strictly hold for values of the dimension $`d`$ such that
$$2<d=2+\frac{2m}{n}<4,\mathrm{for}0<m<n,m,n\text{I}\text{N}.$$
(34)
In particular, in three dimensions one finds
$$f(x)=\omega ^2\left[1+\frac{1}{N}\left(f_1(\omega )+\mathrm{log}\omega f_2(\omega )\right)+O(N^2)\right],$$
(35)
where the functions $`f_1(\omega )`$ and $`f_2(\omega )`$ have a regular expansion in powers of $`\omega `$. Moreover,
$$f_2(\omega )=O(\omega ^2),$$
(36)
so that the logarithms affect the power expansion only at the next-next-to-leading order. A possible interpretation of the large-$`N`$ analysis is that the expansion (35) holds for all values of $`N`$, so that Eq. (33) is not correct due to the presence of logarithms. The reason of their appearance is however unclear. Neverthless, it does not contradict the conjecture that the behavior near the coexistence curve is controlled by the zero-temperature infrared-stable Gaussian fixed point. In this case logarithms would not be unexpected, as they usually appear in the reduced-temperature asymptotic expansion around Gaussian fixed points (see e.g. Ref. ).
### C Parametric representations
In order to obtain a representation of the critical equation of state that is valid in the whole critical region, one may use parametric representations, which implement in a simple way all scaling and analytic properties. One may parametrize $`M`$ and $`t`$ in terms of $`R`$ and $`\theta `$ according to
$`M`$ $`=`$ $`m_0R^\beta m(\theta ),`$ (37)
$`t`$ $`=`$ $`R(1\theta ^2),`$ (38)
$`H`$ $`=`$ $`h_0R^{\beta \delta }h(\theta ),`$ (39)
where $`h_0`$ and $`m_0`$ are normalization constants. The variable $`R`$ is nonnegative and measures the distance from the critical point in the $`(t,H)`$ plane; it carries the power-law critical singularities. The variable $`\theta `$ parametrizes the displacements along the line of constant $`R`$. The functions $`m(\theta )`$ and $`h(\theta )`$ are odd and regular at $`\theta =0`$ and at $`\theta =1`$. The constants $`m_0`$ and $`h_0`$ can be chosen so that $`m(\theta )=\theta +O(\theta ^3)`$ and $`h(\theta )=\theta +O(\theta ^3)`$. The smallest positive zero of $`h(\theta )`$, which should satisfy $`\theta _0>1`$, represents the coexistence curve, i.e. $`T<T_c`$ and $`H0`$.
The parametric representation satisfies the requirements of regularity of the equation of state. Singularities can appear only at the coexistence curve (due for example to the logarithms discussed in Sec. III B), i.e. for $`\theta =\theta _0`$. Notice that the mapping (37) is not invertible when its Jacobian vanishes, which occurs when
$$Y(\theta )(1\theta ^2)m^{}(\theta )+2\beta \theta m(\theta )=0.$$
(40)
Thus the parametric representations based on the mapping (37) are acceptable only if $`\theta _0<\theta _l`$ where $`\theta _l`$ is the smallest positive zero of the function $`Y(\theta )`$. One may easily verify that the asymptotic behavior (32) is reproduced simply by requiring that
$$h(\theta )\left(\theta _0\theta \right)^2\mathrm{for}\theta \theta _0.$$
(41)
The relation among the functions $`m(\theta )`$, $`h(\theta )`$ and $`F(z)`$ is given by
$`z=\rho m(\theta )\left(1\theta ^2\right)^\beta ,`$ (42)
$`F(z(\theta ))=\rho \left(1\theta ^2\right)^{\beta \delta }h(\theta ),`$ (43)
where $`\rho `$ is a free parameter . Indeed, in the exact parametric equation the value of $`\rho `$ may be chosen arbitrarily but, as we shall see, when adopting an approximation procedure the dependence on $`\rho `$ is not eliminated. In our approximation scheme we will fix $`\rho `$ to ensure the presence of the Goldstone singularities at the coexistence curve, i.e. the asymptotic behavior (41). Since $`z=\rho \theta +O(\theta ^3)`$, expanding $`m(\theta )`$ and $`h(\theta )`$ in (odd) powers of $`\theta `$,
$`m(\theta )`$ $`=`$ $`\theta +{\displaystyle \underset{n=1}{}}m_{2n+1}\theta ^{2n+1},`$ (44)
$`h(\theta )`$ $`=`$ $`\theta +{\displaystyle \underset{n=1}{}}h_{2n+1}\theta ^{2n+1},`$ (45)
and using Eqs. (42) and (43), one can find the relations among $`\rho `$, $`m_{2n+1}`$, $`h_{2n+1}`$ and the coefficients $`F_{2n+1}`$ of the expansion of $`F(z)`$.
One may also write the scaling function $`f(x)`$ in terms of the parametric functions $`m(\theta )`$ and $`h(\theta )`$:
$`x={\displaystyle \frac{1\theta ^2}{\theta _0^21}}\left[{\displaystyle \frac{m(\theta _0)}{m(\theta )}}\right]^{1/\beta },`$ (46)
$`f(x)=\left[{\displaystyle \frac{m(\theta )}{m(1)}}\right]^\delta {\displaystyle \frac{h(\theta )}{h(1)}}.`$ (47)
In App. A we report the definitions of some universal ratios of amplitudes that have been introduced in the literature, and the corresponding expressions in terms of the functions $`m(\theta )`$ and $`h(\theta )`$.
### D Approximate polynomial representations
In order to construct approximate parametric representations we consider polynomial approximations of $`m(\theta )`$ and $`h(\theta )`$. This kind of approximation turned out to be effective in the case of Ising-like sistems . The major difference with respect to the Ising case is the presence of the Goldstone singularities at the coexistence curve. In order to take them into account, at least in a simplified form which neglects the logarithms found in Eq. (35), we require the function $`h(\theta )`$ to have a double zero at $`\theta _0`$ as in Eq. (41). Polynomial schemes may in principle reconstruct also the logarithms, but of course only in the limit of an infinite number of terms.
In order to check the accuracy of the results, it is useful to introduce two distinct schemes of approximation. In the first one, which we denote as (A), $`h(\theta )`$ is a polynomial of fifth order with a double zero at $`\theta _0`$, and $`m(\theta )`$ a polynomial of order $`(1+2n)`$:
$`\mathrm{scheme}(\mathrm{A}):`$ $`m(\theta )=\theta \left(1+{\displaystyle \underset{i=1}{\overset{n}{}}}c_i\theta ^{2i}\right),`$ (49)
$`h(\theta )=\theta \left(1\theta ^2/\theta _0^2\right)^2.`$
In the second scheme, denoted by (B), we set
$`\mathrm{scheme}(\mathrm{B}):`$ $`m(\theta )=\theta ,`$ (51)
$`h(\theta )=\theta \left(1\theta ^2/\theta _0^2\right)^2\left(1+{\displaystyle \underset{i=1}{\overset{n}{}}}c_i\theta ^{2i}\right).`$
Here $`h(\theta )`$ is a polynomial of order $`5+2n`$ with a double zero at $`\theta _0`$. In both schemes the parameter $`\rho `$ is fixed by the requirement (41), while $`\theta _0`$ and the $`n`$ coefficients $`c_i`$ are determined by matching the small-field expansion of $`F(z)`$. This means that, for both schemes, in order to fix the $`n`$ coefficients $`c_i`$ we need to know $`n+1`$ values of $`r_{2j}`$, i.e. $`r_6,\mathrm{}r_{6+2n}`$. Note that for the scheme (B)
$$Y(\theta )=1\theta ^2+2\beta \theta ^2,$$
(52)
independently of $`n`$, so that $`\theta _l=(12\beta )^1`$. Concerning the scheme (A), we note that the analyticity of the thermodynamic quantities for $`|\theta |<\theta _0`$ requires the polynomial function $`Y(\theta )`$ not to have complex zeroes closer to origin than $`\theta _0`$.
In App. B we present a more general discussion on the parametric representations.
### E Results
As input parameters for the determination of the parametric representations, we use the best available estimates of the critical exponents, which are $`\alpha =0.01285(38)`$ (from the experiment of Ref. ), $`\eta =0.0381(3)`$ (from the high-temperature analysis of Ref. ). Moreover we use the following estimates of $`r_{2j}`$: $`r_6=1.96(2)`$ which is compatible with all the estimates of $`r_6`$ reported in Table II, and $`r_8=1.40(15)`$ which takes somehow into account the differences among the various estimates.
The case $`n=0`$ of the two schemes (A) and (B) is the same, and requires the knowledge of $`\alpha `$, $`\eta `$ and $`r_6`$. Unfortunately this parametrization does not satisfy the consistency condition $`\theta _0^2<\theta _l^2=(12\beta )^1`$. Both schemes give acceptable approximations for $`n=1`$, using $`r_8`$ as an additional input parameter. The numerical values of the relevant parameters and the resulting estimates of universal amplitude ratios (see the appendix for their definition) are shown in Table III. The errors reported are related to the errors of the input parameters only. They do not take into account possible systematic errors due to the approximate procedure we are employing. We will return on this point later.
In Figs. 1 and 2 we show respectively the scaling functions $`F(z)`$ and $`f(x)`$, as obtained from the approximate representations given by the schemes (A) and (B) for $`n=1`$, using the input values $`\alpha =0.01285`$, $`\eta =0.0381`$, $`r_6=1.96`$ and $`r_8=1.4`$. The two approximations of $`F(z)`$ are practically indinstinguishable in Fig. 1. This is also numerically confirmed by the estimates of the universal costant $`F_0^{\mathrm{}}`$ (reported in Table III), which is related to the large-$`z`$ behavior of $`F(z)`$:
$$F(z)F_0^{\mathrm{}}z^\delta \mathrm{for}z\mathrm{}.$$
(53)
This agreement is not trivial since the small-$`z`$ expansion has a finite convergence radius given by $`|z_0|=R_4^{1/2}2.8`$. Therefore, the determination of $`F(z)`$ on the whole positive real axis from its small-$`z`$ expansion requires an analytic continuation, which turns out to be effectively performed by the approximate parametric representations we have considered. We recall that the large-$`z`$ limit corresponds to the critical theory $`t=0`$, so that positive real values of $`z`$ describe the high-temperature phase up to $`t=0`$. Instead, larger differences between the approximations given by the schemes (A) and (B) for $`n=1`$ appear in the scaling function $`f(x)`$, especially in the region $`x<0`$ corresponding to $`t<0`$ (i.e. the region which is not described by real values of $`z`$). Note that the apparent differences for $`x>0`$ are essentially caused by the normalization of $`f(x)`$, which is performed at the coexistence curve $`x=1`$ and at the critical point $`x=0`$ requiring $`f(1)=0`$ and $`f(0)=1`$. Although the large-$`x`$ region corresponds to small $`z`$, the difference between the two approximate schemes does not decrease in the large-$`x`$ limit due to their slightly different estimates of $`R_\chi `$ (see Table III). Indeed, for large values of $`x`$, $`f(x)`$ has an expansion of the form
$$f(x)=x^\gamma \underset{n=0}{}f_n^{\mathrm{}}x^{2n\beta }$$
(54)
with $`f_0^{\mathrm{}}=R_\chi ^1`$.
We also considered the case $`n=2`$, using the estimate $`r_{10}=13(7)`$. In this case the scheme (A) was not particularly useful because it turned out to be very sensitive to $`r_{10}`$, whose estimate has a relatively large error. Combining the consistency condition $`\theta _0<\theta _l`$ (which excludes values of $`r_{10}10`$ when using the central values of $`\alpha `$, $`\eta `$, $`r_6`$ and $`r_8`$) with the IHT estimate of $`r_{10}`$, we found a rather good result for $`A^+/A^{}`$, i.e. $`A^+/A^{}=1.053(4)`$. On the other hand, the results for the other universal amplitude ratios considered, such as $`R_\chi `$, $`R_c`$, $`R_4`$, etc…, although consistent, turned out to be much more imprecise than those obtained for $`n=1`$, for example $`R_\chi =1.2(3)`$. This fact may be explained noting that, except for the small interval $`10r_{10}9`$ (this interval corresponds to the central values of the other input parameters), the function $`Y(\theta )`$, cf. Eq. (40), has zeroes in the the complex plane which are closer to the origin than $`\theta _0`$. Therefore, the parametric function $`g_2(\theta )`$ related to the magnetic susceptibility (see App. A 2), and higher-order derivatives of the free-energy, have poles within the disk $`|\theta |<\theta _0`$. On the other hand, in the case $`n=1`$, $`\theta _0^2`$ was closer to the origin than the zeroes of $`Y(\theta )`$ for the whole range of values of the input parameters.
For these reasons, for $`n=2`$, we present results only for the scheme (B). Combining the consistency condition $`\theta _0<\theta _l`$, which restricts the acceptable values of $`r_{10}`$ (for example using the central estimates of the other input parameters it excludes values $`|r_{10}|12`$), with the IHT estimate $`r_{10}=13(7)`$, we arrive at the results reported in Table III (third column of data). The reported value has been obtained using $`r_{10}9`$.
The coefficients $`c_i`$, reported in Table III, turn out to be relatively small in both schemes, and decrease rapidly, supporting our choice of the approximation schemes. The results of the various approximate parametric representations are in reasonable agreement. Their comparison is useful to get an idea of the systematic error due to the approximation schemes. There is a very good agreement for $`A^+/A^{}`$, which is the experimentally most important quantity. Our final estimate is
$$A^+/A^{}=1.055(3).$$
(55)
We mention that approximately one half of the error is due to the uncertainty on the critical exponent $`\alpha `$, which unfortunately can be hardly improved by present theoretical means. The approximation schemes (A) and (B) with $`n=1`$ provide independent results for $`r_{10}`$, leading to the estimate $`r_{10}=10(3)`$, which is agreement with the IHT result $`r_{10}=13(7)`$. The determination of $`c_f`$, cf. Eq. (32), turns out to be rather unstable, indicating that the approximate parametric representation we have constructed are still relatively inaccurate in the region very close to the coexistence curve. The constant $`c_f`$ is very sensitive to the values of the coefficients $`r_{2j}`$. Improved estimates of $`r_{2j}`$ would be important especially for $`c_f`$.
Finally, to further check our results, we applied again the scheme (B) for $`n=2`$, replacing $`r_{10}`$ with the precise experimental estimate $`A^+/A^{}=1.054(1)`$ as input parameter. In practice we fix the coefficient $`c_2`$ in such a way to obtain the experimental estimate of $`A^+/A^{}`$. The idea is to use the quantities known with the highest precision to determine, within our scheme of approximation, the equation of state and the corresponding universal amplitude ratios. The results are reported in the last column of Table III.
From the results of Table III we arrive at the final estimates
$`R_\xi ^+(A^+)^{1/3}f^+=0.353(3),`$ (56)
$`R_c{\displaystyle \frac{\alpha A^+C^+}{B^2}}=0.12(1),`$ (57)
$`R_\chi {\displaystyle \frac{C^+B^{\delta 1}}{(\delta C^c)^\delta }}=1.4(1),`$ (58)
$`R_4{\displaystyle \frac{C_4^+B^2}{(C^+)^3}}=|z_0|^2=7.6(4),`$ (59)
$`F_0^{\mathrm{}}\underset{z\mathrm{}}{lim}z^\delta F(z)=0.0303(3),`$ (60)
$`0<c_f20.`$ (61)
In Table IV we compare our results with the available estimates obtained from other theoretical approches and from experiments (for a review see e.g. Ref. ). Our precision for $`A^+/A^{}`$ is comparable with the estimate reported in Ref. , obtained in the field-theoretic framework of the minimal renormalization without $`ϵ`$-expansion, which is a perturbative expansion at fixed dimension $`d=3`$. The agreement with the experimental result of Ref. is very good.
### F Conclusions
Starting from the small-field expansion of the effective potential in the high-temperature phase, we have constructed approximate representations of the critical equation of state valid in the whole critical region. We have considered two approximation schemes based on polynomial representations that satisfy the general analytic properties of the equation of state (Griffith’s analyticity) and take into account the Goldstone singularities at the coexistence curve. The coefficients of the truncated polynomials are determined by matching the small-field expansion in the high-temperature phase, which has been studied by lattice high-temperature techniques. The schemes considered can be systematically improved by increasing the order of the polynomials. However, such possibility is limited by the number of known coefficients $`r_{2j}`$ of the small-field expansion of the effective potential. We have shown that the knowledge of the first few $`r_{2j}`$ already leads to satisfactory results, for instance for the specific-heat amplitude ratio. Through the approximation schemes we have presented in this paper, the determination of the equation of state may be improved by a better determination of the coefficients $`r_{2j}`$, which may be achieved by extending the high-temperature expansion. We hope to return on this issue in the future.
Finally we mention that the approximation schemes which we have proposed can be applied to other $`N`$-vector models. Physically relevant values are $`N=3`$ and $`N=4`$. The case $`N=3`$ describes the critical phenomena in isotropic ferromagnets . The case $`N=4`$ is interesting for high-energy physics: it should describe the critical behavior of finite-temperature QCD with two flavours of quarks at the chiral-symmetry restoring phase transition .
## A Universal ratios of amplitudes
### 1 Notations
Universal ratios of amplitudes characterize the critical behavior of thermodynamic quantities that do not depend on the normalizations of the external (e.g. magnetic) field, order parameter (e.g. magnetization) and temperature. Amplitude ratios of zero-momentum quantities can be derived from the critical equation of state. We consider several amplitudes derived from the singular behavior of the specific heat
$$C_H=A^\pm |t|^\alpha ,$$
(A1)
the magnetic susceptibility in the high-temperature phase
$$\chi =\frac{1}{2}C^+t^\gamma ,$$
(A2)
the zero-momentum four-point connected correlation function in the high temperature phase
$$\chi _4=\frac{8}{3}C_4^+t^{\gamma 2\beta \delta },$$
(A3)
the second-moment correlation length in the high-temperature phase
$$\xi =f^+t^\nu ,$$
(A4)
and the spontaneous magnetization on the coexistence curve
$$M=B|t|^\beta .$$
(A5)
Using the above normalizations for the amplitudes, the zero-momentum four-point coupling $`g_4`$, cf. Eq. (19), can be written as
$$g_4=\frac{C_4^+}{(C^+)^2(f^+)^3}$$
(A6)
In addition, one can also define amplitudes along the critical isotherm, such as
$$\chi _L=C^c|H|^{\gamma /\beta \delta }.$$
(A7)
### 2 Universal ratios of amplitudes from the parametric representation
In the following we report the expressions of the universal ratios of amplitudes in terms of the parametric representation (37) of the critical equation of state.
The singular part of the free energy per unit volume can be written as
$$=h_0m_0R^{2\alpha }g(\theta ),$$
(A8)
where $`g(\theta )`$ is the solution of the first-order differential equation
$$(1\theta ^2)g^{}(\theta )+2(2\alpha )\theta g(\theta )=Y(\theta )h(\theta )$$
(A9)
that is regular at $`\theta =1`$. The function $`Y(\theta )`$ has been defined in Eq. (40). The longitudinal magnetic susceptibility can be written as
$$\chi _L^1=\frac{h_0}{m_0}R^\gamma g_2(\theta ),g_2(\theta )=\frac{2\beta \delta \theta h(\theta )+(1\theta ^2)h^{}(\theta )}{Y(\theta )}.$$
(A10)
The function $`g_2(\theta )`$ must vanish at $`\theta _0`$ in order to reproduce the predicted Goldstone singularities, according to
$$g_2(\theta )\theta _0\theta \mathrm{for}\theta \theta _0.$$
(A11)
From Eq. (A10) we see that $`g_2(\theta )`$ satisfies this condition if $`h(\theta )(\theta _0\theta )^2`$ for $`\theta \theta _0`$.
From the equation of state one can derive universal amplitude ratios of quantities defined at zero momentum, i.e. integrated in the volume. We consider
$`A^+/A^{}=(\theta _0^21)^{2\alpha }{\displaystyle \frac{g(0)}{g(\theta _0)}},`$ (A12)
$`R_c{\displaystyle \frac{\alpha A^+C^+}{B^2}}=\alpha (1\alpha )(2\alpha )(\theta _0^21)^{2\beta }[m(\theta _0)]^2g(0),`$ (A13)
$`R_4{\displaystyle \frac{C_4^+B^2}{(C^+)^3}}=|z_0|^2=\rho ^2[m(\theta _0)]^2\left(\theta _0^21\right)^{2\beta },`$ (A14)
$`R_\chi {\displaystyle \frac{C^+B^{\delta 1}}{(\delta C^c)^\delta }}=(\theta _0^21)^\gamma [m(\theta _0)]^{\delta 1}[m(1)]^\delta h(1),`$ (A15)
Using Eqs. (42) and (43) one can compute $`F(z)`$ and obtain the small-$`z`$ expansion coefficients of the effective potential $`r_{2j}`$. The constant $`F_0^{\mathrm{}}`$, which is related to the behavior of $`F(z)`$ for $`z\mathrm{}`$, cf. Eq. (31), is given by
$$F_0^{\mathrm{}}\underset{z\mathrm{}}{lim}z^\delta F(z)=\rho ^{1\delta }[m(1)]^\delta h(1).$$
(A16)
Using the relations (46) concerning the scaling function $`f(x)`$, one can easily obtain the constant $`c_f`$, which is related to the behavior of $`f(x)`$ for $`x1`$, cf. Eq. (32),
$$c_f\underset{x1}{lim}(1+x)^2f(x).$$
(A17)
We consider also the universal amplitude ratio $`R_\xi ^+(A^+)^{1/3}f^+`$ which can be obtained from the estimates of $`R_4`$, $`R_c`$ and $`g_4`$:
$$R_\xi ^+(A^+)^{1/3}f^+=\left(\frac{R_4R_c}{g_4}\right)^{1/3}.$$
(A18)
We mention that in the case of the superfluid helium it is customary to define also the hyperuniversal combination
$$R_\xi ^\mathrm{T}(A^{})^{1/3}f_\mathrm{T}^{},$$
(A19)
where $`f_\mathrm{T}^{}`$ is the amplitude of a transverse correlation length $`\xi _\mathrm{T}`$ defined from the stiffness constant $`\rho _s`$, i.e. $`\xi _\mathrm{T}=\rho _s^1`$. $`R_\xi ^\mathrm{T}`$ can be determined directly from experiments below $`T_c`$ (see e.g. Ref. ).
## B General discussion on the parametric representations
A wide family of parametric representations was introduced a long time ago , in forms that can be related to our Eqs. (37). Since the application of parametric representations in practice requires some approximation scheme, one may explore the freedom left in these representations and understand how this freedom may be exploited in order to optimize the approximation. The parametric form of the equation of state forces relations between the two functions $`\overline{m}\rho m(\theta )`$ and $`h(\theta )`$, but it is easy to get convinced that one of the two functions can be chosen arbitrarily. For definiteness, let’s take $`\overline{m}`$ to be arbitrary and find the constraints that must be satisfied by $`h(\theta )`$ as a consequence of the equation of state.
It is convenient for our purposes to establish these constraints by imposing the formal independence of the function $`F(z)`$ from the parametrization adopted for $`\overline{m}\rho m(\theta )`$, which we may symbolically write in the form of a functional equation,
$$\frac{\delta }{\delta \overline{m}}\left[\frac{\rho h(\theta )}{(1\theta ^2)^{\beta \delta }}\right]=0,$$
(B1)
keeping $`z`$ fixed. By expanding $`m(\theta )`$ according to Eq. (44) and treating the coefficients $`\rho `$ and $`c_nm_{2n+1}`$ as variational parameters, we may turn the above equation into a set of partial differential equations (keeping $`z`$ fixed)
$`{\displaystyle \frac{d}{d\rho }}\left[{\displaystyle \frac{\rho h(\theta )}{(1\theta ^2)^{\beta \delta }}}\right]=0,`$ (B2)
$`{\displaystyle \frac{d}{dc_i}}\left[{\displaystyle \frac{\rho h(\theta )}{(1\theta ^2)^{\beta \delta }}}\right]=0,`$ (B3)
which must be satisfied exactly for all $`i`$ by the function $`h(\theta )`$. Simple manipulations lead to the following explicit form:
$$Y(\theta )\left(h+\rho \frac{h}{\rho }\right)=m(\theta )\left[(1\theta ^2)\frac{h}{\theta }+2\beta \delta \theta h\right],$$
(B4)
$$Y(\theta )\frac{h}{c_i}=\theta ^{2i+1}\left[(1\theta ^2)\frac{h}{\theta }+2\beta \delta \theta h\right],$$
(B5)
where $`Y(\theta )`$ is defined in Eq. (52). In turn, by expanding
$$h(\theta ,\rho ,c_i)=\theta +\underset{n=1}{\overset{\mathrm{}}{}}h_{2n+1}(\rho ,c_i)\theta ^{2n+1},$$
(B6)
and substituting into the above equations one obtains an infinite set of linear differential recursive equations for the coefficients $`h_{2n+1}`$, which generalize the relations found in Ref. , where the case $`m(\theta )=\theta `$ was analyzed. A typical approximation to the exact parametric equation of state amounts to a truncation of $`h`$ to a polynomial form. We may in this case refine the approximation by reinterpreting the first recursion equations involving a coefficient $`h_{2t+1}`$ which is forcefully set equal to zero as stationarity conditions, which force the parameters $`\rho `$ and $`c_i`$ into the values minimizing the unwanted dependence of the truncated $`F(z)`$ on the parameters themselves, i.e. on the choice of the function $`\overline{m}\rho m(\theta )`$. The above procedure implies global (i.e. $`\theta `$-independent) stationarity, and as a consequence all physical amplitudes turn out to be stationary with respect to variations of $`\rho `$ and $`c_i`$.
These statements are fairly general, but it is certainly interesting to consider the first few non trivial examples. For the lowest order truncation of $`h`$ we may adopt the parametrization
$$h(\theta )=\theta \left(1\frac{\theta ^2}{\theta _0^2}\right)^p,$$
(B7)
which includes both the Ising model in three dimensions ($`p=1`$) and general $`O(N)`$ symmetric models with Goldstone bosons in $`d`$ dimensions ($`p=2/(d2)`$). It is easy to recognize that the following relationship must then hold:
$$\frac{1}{6}\rho ^2+c_1=\gamma \frac{p}{\theta _0^2}.$$
(B8)
Global stationarity implies that the stability conditions may be extracted from the variation of any physical quantity. In particular we may concentrate on the universal zero of $`F(z)`$, $`z_0`$, noting that $`z_0=z(\theta _0)`$, cf. Eq. (42). As a consequence of the above results, the simplest models can all be described by the parametrization
$$\frac{z_0}{\sqrt{6}}=\frac{\sqrt{(\gamma c_1)\theta _0^2p}}{(1\theta _0^2)^\beta }(1+c_1\theta _0^2).$$
(B9)
Let us first consider the case $`c_1=0`$. The minimization procedure leads to
$`\rho ^2={\displaystyle \frac{6\gamma (\gamma p)}{\gamma 2p\beta }},`$ (B10)
$`\theta _0^2={\displaystyle \frac{\gamma 2p\beta }{(12\beta )\gamma }}.`$ (B11)
Setting $`p=1`$ one immediately recognizes the linear parametric model representation for the Ising model . Unfortunately when $`p=2`$ the solution is not physically satisfactory, because it gives $`\theta _0^2<0`$ for all $`N2`$, and therefore the scheme is useless for models with Goldstone singularities.
Let us now include $`c_1`$. Requiring $`z_0`$ to be stationary with respect to variations of both parameters, we obtain
$`\rho ^2={\displaystyle \frac{6\gamma (\gamma p+1)}{3\gamma 2(p1)\beta }},`$ (B12)
$`\theta _0^2={\displaystyle \frac{3\gamma 2(p1)\beta }{(32\beta )\gamma }},`$ (B13)
$`c_1=\gamma {\displaystyle \frac{2(\gamma p)+(2\beta 1)}{3\gamma 2(p1)\beta }},`$ (B14)
implying also
$$\frac{|z_0|}{\sqrt{6}}=2\left(\frac{\gamma p+1}{32\beta }\right)^{\frac{3}{2}\beta }\left(\frac{\gamma }{2\beta }\right)^\beta .$$
(B15)
In the Ising model the above solution reduces to
$`\rho ^2=2\gamma ,`$ (B16)
$`\theta _0^2={\displaystyle \frac{3}{32\beta }},`$ (B17)
$`c_1={\displaystyle \frac{2}{3}}(\beta +\gamma )1.`$ (B18)
Note that, substituting the physical values of the critical exponents $`\beta `$ and $`\gamma `$ for the $`N=1`$ model, $`c_1`$ turns out to be a very small number ($`c_1=0.04256`$) and the predicted numerical value of $`z_0`$ is 2.8475, consistent within $`1\%`$ with the linear parametric model prediction . It is fair to say that in the Ising case the above solution has a status which is comparable to the linear parametric model, both conceptually and in terms of predicting power. It is therefore possible to take it as the starting point of an alternative approximation scheme whose higher-order truncations might prove quite effective.
Unfortunately when we consider the $`XY`$ system, setting $`p=2`$ and choosing the values of the exponents pertaining to $`N=2`$, the value of $`c_1`$ becomes too large for the approximation to be sensible. Indeed we get $`c_1=0.6762`$ and all testable predictions turn out to be far away from the corresponding physical values. It is however worth exploring the features of this approach because, as we shall show, it has formal properties which might prove useful when considering parametric representations of the equation of state for higher values of $`N`$. Let us indeed consider the function $`g_2(\theta )`$ entering the parametric representation of the magnetic susceptibility. We know that this function will in general show singularities in the complex $`\theta `$ plane corresponding to the zeroes of the function $`Y(\theta )`$. However, when substituting the expressions of $`h(\theta )`$ and $`m(\theta )`$ obtained from the saddle-point evaluation of the parameters $`\rho `$, $`\theta _0`$ and $`c_1`$, after some simple manipulations, we find out that all singularities cancel and
$$g_2(\theta )=\left(1\frac{\theta ^2}{\theta _0^2}\right)^{p1}.$$
(B19)
This fact was already observed in the case $`c_1=0`$ for all values of the truncation order $`t`$. Therefore, the stationarity prescription is a way to ensure a higher degree of regularity in the parametric representation of thermodynamic functions.
Finally, let us observe that the stationary solution can be applied to the large-$`N`$ limit of $`O(N)`$ models in any dimension $`2<d<4`$. In this limit $`\beta =\frac{1}{2}`$, $`\gamma =p=2/(d2)`$. As a consequence we obtain from our previous results $`\rho ^2=12/(d+2)`$, $`\theta _0^2=(d+2)/4`$, $`c_1=0`$, implying also
$`h(\theta )=\theta \left(1{\displaystyle \frac{4}{d+2}}\theta ^2\right)^{\frac{2}{d2}},`$ (B20)
$`g_2(\theta )=\left(1{\displaystyle \frac{4}{d+2}}\theta ^2\right)^{\frac{4d}{d2}}.`$ (B21)
We therefore obtained an exact parametrization of the equation of state in the large-$`N`$ limit for all $`d`$. Thus, for sufficiently large values of $`N`$, the scheme we have defined may be a sensible starting point for the parametric representation of the thermodynamical functions in the critical domain. |
warning/0001/math0001131.html | ar5iv | text | # Quotients of divisorial toric varieties
## Introduction
It is a frequently occuring question in algebraic geometry, if an algebraic group action $`G\times XX`$ admits a categorical quotient, i.e., a regular map $`XY`$ that is universal with respect to $`G`$-invariant regular maps $`XZ`$. For example, moduli functors are often corepresented by categorical quotients. In general, it is a difficult problem to decide whether a categorical quotient exists. Some counterexamples for actions of the multiplicative group $`^{}`$ are presented in .
As these examples show, difficulties already arise with subtorus actions on toric varieties. Such actions have been investigated by several authors, mainly focusing on the much more restrictive concept of a good quotient, see e.g. , and . The description of toric varieties in terms of rational fans relates the problem of constructing quotients to problems of combinatorial convexity. Hence the class of toric varieties serves as a testing ground for more general ideas.
Now, let $`X`$ be a toric variety and let $`H`$ be a subtorus of the big torus of $`X`$. Our approach to categorical quotients for the induced action of $`H`$ on $`X`$ is to consider the problem in suitable subcategories. A first step is to construct a quotient in the category of toric varieties itself: In , we showed that there always exists a toric quotient
$$p:XX/\mathrm{tq}H.$$
This is a toric morphism that is universal with respect to $`H`$-invariant toric morphisms. The essential part of the proof is an explicit algorithm in terms of combinatorial data. The toric quotient is a canonical starting point for quotients in further categories. For example, in we gave an explicit method to decide by means of the toric quotient when a subtorus action on a quasiprojective toric variety admits a categorical quotient in the category of quasiprojective varieties.
In the present article we give a considerable generalization of the results of , namely we solve the analogous problem in the category of divisorial varieties. Recall that an irreducible variety $`X`$ is called divisorial if every point $`xX`$ has an affine neighbourhood of the form $`X\mathrm{Supp}(D)`$ with an effective Cartier divisor $`D`$ on $`X`$, see e.g. and \[8, II.2.2\].
The class of divisorial varieties contains the quasiprojective varieties as well as all $``$-factorial varieties. It has nice functorial properties, see , and moreover it often provides a natural framework to extend statements known to hold for quasiprojective varieties on the one hand and for smooth varieties on the other hand.
A connection to toric geometry is provided by the embedding results of : A variety is divisorial if and only if it admits a closed embedding into a smooth toric prevariety $`Z`$ having an affine diagonal map $`ZZ\times Z`$. The equivariant version of this statement implies in particular that a toric variety is divisorial if and only if it has enough invariant effective Cartier divisors in the sense of T. Kajiwara , see Section 1.
Now, given a divisorial toric variety $`X`$ and a subtorus $`H`$ of the big torus of $`X`$, when does the action of $`H`$ on $`X`$ admit a categorical quotient in the category of divisorial varieties? As mentioned, we start with the toric quotient
$$p:XX/\mathrm{tq}H.$$
A first problem is that in general the toric quotient space $`X/\mathrm{tq}H`$ is not a divisorial variety. To deal with this effect, we construct a toric divisorial reduction. This is a toric morphism
$$q:X/\mathrm{tq}H(X/\mathrm{tq}H)^{\mathrm{tdr}}$$
which is universal with respect to toric morphism to divisorial toric varieties. The question then is, how these toric constructions behave in the essentially larger category of arbitrary divisorial varieties. Our main result gives the following answer, see Corollary 6.3:
###### Theorem.
The action of $`H`$ on $`X`$ admits a categorical quotient in the category of divisorial varieties if and only if the composition $`qp`$ is surjective. Moreover, in the latter case, $`qp`$ is the desired categorical quotient.
The paper is organized as follows: In Section 1 we discuss divisoriality in the context of $`G`$-varieties and provide some general statements used in the subsequent constructions. Sections 2 and 3 are devoted to the construction of the toric divisorial reduction. This is done in the language of combinatorial convexity. The main tool are convex support maps extending the notion of a convex support function on a fan.
Generalizing the corresponding well-known statement on projectivity and support functions, we show that divisoriality of a given toric variety is characterized by the existence of a strictly convex support map on its fan. Moreover, we relate convex support maps to toric morphisms to divisorial toric varieties. This allows the construction of the toric divisorial reduction. Finally, we present some examples in Section 3.
In Sections 4 and 5 we prepare the proof of the main results. The essential task is to reduce arbitrary $`H`$-invariant regular maps to $`H`$-invariant toric morphisms. This is done by the Decomposition Lemma presented in Section 5: Given an $`H`$-invariant regular map $`f:XY`$ to a divisorial variety, we construct a decomposition $`f=hg`$ with an $`H`$-invariant toric morphism $`g`$ followed by a rational map $`h`$ defined near $`g(X)`$.
The ingredients for the proof of this Decomposition Lemma are the abovementioned embedding of $`Y`$ into a certain smooth toric prevariety $`Z`$ provided by and the following lifting result, presented in Section 4: There exist quasiaffine toric varieties $`\stackrel{~}{X}`$ and $`\stackrel{~}{Z}`$ “above” $`X`$ and $`Z`$ respectively such that the map $`f`$ admits a lifting $`\stackrel{~}{f}:\stackrel{~}{X}\stackrel{~}{Z}`$. This basically reduces the decomposition problem to the case of quasiaffine toric varieties.
In Section 6 we give statements and proofs of the main results. Finally, in Section 7 we formulate an open problem on categorical quotients for subtorus actions on toric varieties.
## 1. Divisorial $`G`$-varieties
Throughout the whole article, we work over a fixed algebraically closed field $`𝕂`$. So a prevariety is a reduced irreducible scheme of finite type over $`𝕂`$, and a variety is a separated prevariety. We say that a prevariety $`X`$ is of affine intersection, if its diagonal morphism $`XX\times X`$ is affine.
As usual, when we speak of a $`G`$-(pre-)variety where $`G`$ is an algebraic group, we mean an algebraic (pre-)variety $`X`$ together with a $`G`$-action given by a regular map $`G\times XX`$. For the basic notions on toric varieties and prevarieties, we refer to and .
In this section, we provide some general facts on group actions on divisorial varieties. Following Borelli , we call a prevariety $`X`$ divisorial if every point $`xX`$ has an affine open neighbourhood of the form $`U=X\mathrm{Supp}(D)`$ with an effective Cartier divisor $`D`$ on $`X`$.
###### Remark 1.1.
1. Quasiprojective varieties are divisorial.
2. Locally closed subspaces of divisorial prevarieties are divisorial.
3. Every divisorial prevariety $`X`$ is of affine intersection.
4. Every $``$-factorial prevariety of affine intersection is divisorial.
A geometric quotient for the action of a reductive group $`G`$ on a variety $`X`$ is an affine regular map $`p:XY`$ such that the fibres of $`p`$ are precisely the $`G`$-orbits and the canonical homomorphism $`𝒪_Yp_{}(𝒪_X)^G`$ is bijective. The analogous notion in the setting of prevarieties, i.e. for possibly non-separated $`X`$ and $`Y`$, is called a geometric prequotient.
In the sequel, we shall make use of the following characterization of divisoriality in terms of geometric quotients and closed embeddings, see \[14, Theorem 3.1\]:
###### Theorem 1.2.
A variety $`X`$ is divisorial if and only if one of the following statements holds:
1. $`X`$ is a geometric quotient of a quasiaffine variety by a free algebraic torus action.
2. $`X`$ admits a closed embedding into a smooth toric prevariety of affine intersection.
Here a torus action is called free if every orbit map is a locally closed embedding. The above result has the following equivariant version, see \[14, Theorem 3.4\]:
###### Theorem 1.3.
Let $`X`$ be a normal divisorial $`T`$-variety where $`T`$ is an algebraic torus acting effectively.
1. There is a quasiaffine variety $`\widehat{X}`$ with a regular action of a torus $`T\times H`$ such that $`H`$ acts freely with a $`T`$-equivariant geometric quotient $`\widehat{X}X`$.
2. There is a $`T`$-equivariant closed embedding $`XZ`$ into a smooth toric prevariety $`Z`$ of affine intersection where $`T`$ acts as a subtorus of the big torus.
A first consequence is that divisorial varieties with torus actions always have many invariant effective Cartier divisors. For a toric variety this means that it is divisorial if and only if it has enough invariant effective Cartier divisors in the sense defined by T. Kajiwara, see .
###### Proposition 1.4.
Let $`T`$ be an algebraic torus, and let $`X`$ be a normal algebraic $`T`$-variety $`X`$. Then $`X`$ is divisorial if and only if there exist $`T`$-invariant effective Cartier divisors $`D_1,\mathrm{},D_r`$ on $`X`$ such that the sets $`X\mathrm{Supp}(D_i)`$ are affine and cover $`X`$.
###### Proof.
We may assume that $`T`$ acts effectively. Let $`X`$ be divisorial. By Theorem 1.3, there is a $`T`$-equivariant closed embedding of $`X`$ into a smooth toric prevariety $`Z`$ of affine intersection where $`T`$ acts as a subtorus of the big torus. Hence $`X`$ inherits the desired property from $`Z`$. The reverse implication is trivial. ∎
As the example of the rational nodal curve with standard $`𝕂^{}`$-action shows, the assumption of normality is essential in the above statement. Our next result states that divisoriality is inherited by geometric quotients for torus actions:
###### Proposition 1.5.
Let $`T`$ be an algebraic torus and suppose that $`X`$ is a normal $`T`$-variety with geometric quotient $`p:XY`$. Then $`X`$ is divisorial if and only if $`Y`$ is divisorial.
###### Proof.
We may assume that the torus $`T`$ acts effectively on $`X`$. If the quotient variety $`Y`$ is divisorial, then we obtain the desired effective Cartier divisors on $`X`$ by pulling back suitable divisors from $`Y`$. Conversely, suppose that $`X`$ is divisorial. Then, by Theorem 1.3, we may assume in the proof that $`X`$ is a quasiaffine $`T`$-variety.
Given $`yY`$, we have to find an affine open neighbourhood of $`y`$ that is the complement of the support of an effective Cartier divisor on $`Y`$. Choosing any $`T`$-equivariant affine closure of $`X`$, we find a function $`f𝒪(X)`$, homogeneous with respect to some character $`\chi _f\mathrm{X}(T)`$, such that for $`D:=\mathrm{div}(f)`$ the $`T`$-invariant set $`U:=XV(f)=X\mathrm{Supp}(D)`$ is an affine neighbourhood of the fibre $`p^1(y)`$.
By $`T`$-closedness of $`p:XY`$, the set $`V:=p(U)`$ is an open neighbourhood of $`yY`$. Moreover, as a geometric quotient space of the affine $`T`$-variety $`U`$, the set $`V`$ is again affine. Thus, to prove the assertion, we only have to show that $`p(\mathrm{Supp}(D))`$ is the support of an effective Cartier divisor $`E`$ on $`Y`$. We construct local equations for such an $`E`$.
First we claim that every point $`zY`$ has an affine neighbourhood $`V_zY`$ such that on $`U_z:=p^1(V_z)`$ there is an invertible function $`h_z𝒪(U_z)`$ that is homogeneous with respect to some positive multiple $`m_z\chi _f`$. To check this, start with any affine neighbourhood $`V_zY`$ of $`z`$ and choose a point $`xp^1(z)`$. Consider the sublattice $`N\mathrm{X}(T)`$ of characters occuring as weights of homogeneous functions $`g𝒪(U_z)`$ with $`g(x)=1`$.
The sublattice $`N`$ is of full rank in $`\mathrm{X}(T)`$: Otherwise we found a nontrivial one-parameter-subgroup $`\lambda :𝕂^{}T`$ such that $`\chi \lambda =1`$ holds for all $`\chi N`$. It follows that $`\lambda (𝕂^{})`$ is contained in the isotropy group $`T_x`$. On the other hand, the $`T`$-action on $`U_z`$ is effective and closed. Hence $`T_x`$ is finite, a contradiction. Thus $`N`$ is of full rank. In particular, some positive multiple $`m_z\chi _f`$ lies in $`N`$ and our claim follows.
Now cover $`Y`$ by finitely many $`V_z`$ as in the above claim. Then we may assume that all the invertible functions $`h_z𝒪(U_z)`$ are homogeneous with respect to the same multiple $`m\chi _f`$. Every function $`g_z:=f^m/h_z`$ is $`T`$-invariant, regular on $`U_z`$ and vanishes precisely on $`\mathrm{Supp}(D)U_z`$. Since it is $`T`$-invariant, $`g_z`$ may be viewed as a regular function on $`V_z=p(U_z)`$ and there its zero set is just
$$p(\mathrm{Supp}(D)U_z)=p(\mathrm{Supp}(D))V_z.$$
Since every $`g_z/g_z^{}`$ is an invertible regular function on $`V_zV_z^{}`$ it follows that the $`g_z`$ are local equations for the desired Cartier divisor $`E`$ on $`Y`$. ∎
As T. Kajiwara has shown, every toric variety $`X`$ with enough invariant effective Cartier divisors arises as a geometric quotient of a quasiaffine toric variety $`\widehat{X}`$ by an algebraic subgroup of the big torus of $`\widehat{X}`$, see \[15, Theorem 1.9\]. In view of the above results, we can enhance Kajiwara’s statement as follows:
###### Corollary 1.6.
A toric variety $`X`$ is divisorial if and only if there is a quasiaffine toric variety $`\widehat{X}`$ and a toric morphism $`p:\widehat{X}X`$ such that $`\mathrm{ker}(p)`$ is a subtorus of the big torus of $`\widehat{X}`$ and $`p`$ is a geometric quotient for the action of $`\mathrm{ker}(p)`$ on $`\widehat{X}`$.
###### Proof.
If $`X`$ is divisorial, then Theorem 1.3 gives the desired quotient presentation. The converse follows from Proposition 1.5. ∎
Finally, we consider translates of divisorial open subsets with respect to an action of a connected group. If the complement of the subset is small enough, the union of such translates is again divisorial:
###### Lemma 1.7.
Let $`G`$ be a connected linear algebraic group, and let $`X`$ be a normal $`G`$-variety. If $`UX`$ is a divisorial open subset with $`\mathrm{codim}(XU)2`$, then also $`GU`$ is divisorial.
###### Proof.
We may assume that $`X=GU`$ holds. Let $`D_1^U,\mathrm{},D_r^U`$ be Cartier divisors on $`U`$ such that the sets $`U_i:=U\mathrm{Supp}(D_i^U)`$ form an affine cover of $`U`$. By closing components, each $`D_i^U`$ extends to a Weil divisor $`D_i`$ on $`X`$.
We claim that $`X\mathrm{Supp}_{D_i}=U_i`$. To see this, let $`A_i:=XU_i`$. Since $`U_i`$ is affine, $`A_i`$ is of pure codimension one. Clearly $`\mathrm{Supp}(D_i^U)A_i`$ and hence $`\mathrm{Supp}(D_i)A_i`$. Thus $`\mathrm{Supp}(D_i)`$ is a union of irreducible components of $`A_i`$. Moreover we have
$$XU=X(U_i\mathrm{Supp}(D_i^U))=A_i\mathrm{Supp}(D_i^U).$$
Since $`XU`$ has codimension at least two, it follows that for each irreducible component $`A_i^{}`$ of $`A_i`$ its intersection with $`\mathrm{Supp}(D_i^U)`$ is dense in $`A_i^{}`$. This implies $`A_i=\mathrm{Supp}(D_i)`$ and our claim is proved. In particular, we have
$$X=GU=G\underset{i=1}{\overset{r}{}}X\mathrm{Supp}(D_i)=\underset{i=1}{\overset{r}{}}\underset{gG}{}X\mathrm{Supp}(gD_i).$$
Thus it suffices to show that for each $`D_i`$ some multiple is Cartier on $`X`$. This is done as follows: The restriction $`D_i^{}`$ of $`D_i`$ to the regular locus $`X^{}X`$ is Cartier. Since $`X^{}`$ is $`G`$-invariant, we may apply $`G`$-linearization, i.e., replacing $`D_i`$ with a suitable multiple we achieve that $`𝒪_{D_i^{}}`$ is a $`G`$-sheaf, see e.g. \[17, Proposition 2.4\].
We claim that this structure of a $`G`$-sheaf extends canonically to $`𝒪_{D_i}`$. For an open set $`VX`$ let $`V^{}:=VX^{}`$. Given a section $`s𝒪_{D_i}(V)`$, we define its translates $`gs`$ as follows: Translate the restriction $`s^{}𝒪_{D_i}(V^{})`$ to a section $`gs^{}𝒪_{D_i}(gV^{})`$ and then extend $`gs^{}`$ to the desired section $`gs𝒪_{D_i}(gV)`$.
Using the $`G`$-sheaf structure on $`𝒪_{D_i}`$ we see that locally $`𝒪_{D_i}`$ is generated by a single function. That means $`D_i`$ is a Cartier divisor. ∎
## 2. Support maps
Projectivity of a given toric variety is characterized by the existence of a strictly convex support function on its fan, see e.g. . Generalizing the notion of a support function here we introduce the concept of a support map on a fan and define convexity properties for such maps. The main result of this section states that for a given fan existence of a strictly convex support map is equivalent to divisoriality of the associated toric variety.
For a lattice $`N`$, we denote the associated rational vector space by $`N_{}`$. A cone in $`N`$ is a polyhedral (not necessarily strictly) convex cone $`\sigma N_{}`$. A quasifan in $`N`$ is a finite set $`\mathrm{\Lambda }`$ of cones in $`N`$ such that for $`\sigma \mathrm{\Lambda }`$ also every face of $`\sigma `$ belongs to $`\mathrm{\Lambda }`$ and for $`\sigma ,\sigma ^{}\mathrm{\Lambda }`$ the intersection $`\sigma \sigma ^{}`$ is a face of both, $`\sigma `$ and $`\sigma ^{}`$. A fan is a quasifan containing only strictly convex cones.
The support of a quasifan $`\mathrm{\Lambda }`$ is the union of all its cones and is denoted by $`|\mathrm{\Lambda }|`$. A map of quasifans $`\mathrm{\Lambda }`$ in a lattice $`N`$ and $`\mathrm{\Lambda }^{}`$ in a lattice $`N^{}`$ is a lattice homomorphism $`NN^{}`$ such that the associated linear map $`N_{}N_{}^{}`$ maps the cones of $`\mathrm{\Lambda }`$ into cones of $`\mathrm{\Lambda }^{}`$.
For the definition of support maps, fix a lattice $`N`$ and a quasifan $`\mathrm{\Delta }`$ in $`N`$. We say that a map $`N_{}^k`$ is linear on a subset $`AN_{}`$ if its restriction to $`A`$ is the restriction of a linear map.
###### Definition 2.1.
A support map on $`\mathrm{\Delta }`$ is a map $`h:|\mathrm{\Delta }|^k`$ that is linear on every cone $`\sigma \mathrm{\Delta }`$.
For a support map $`h:|\mathrm{\Delta }|^k`$, let $`\gamma `$ be the cone in $`\widehat{N}:=N\times ^k`$ generated by the graph $`\mathrm{\Gamma }_h`$ of $`h`$, and let $`𝔉(\gamma )`$ denote the quasifan consisting of all faces of $`\gamma `$. The filled graph of $`h`$ is the minimal subquasifan $`\mathrm{\Lambda }_h`$ of $`𝔉(\gamma )`$ with $`\mathrm{\Gamma }_h|\mathrm{\Lambda }_h|`$. So, $`\mathrm{\Lambda }_h`$ is generated by the cones $`\delta \gamma `$ whose relative interior $`\delta ^{}`$ meets $`\mathrm{\Gamma }_h`$.
###### Definition 2.2.
The support map $`h:|\mathrm{\Delta }|^k`$ is called convex, if the projection $`P:\widehat{N}_{}N_{}`$ is injective on the support $`|\mathrm{\Lambda }_h|`$.
This notion of convexity includes the classical concept of a convex support function on a complete fan as defined for example in \[12, p. 67\]:
###### Remark 2.3.
Let $`h:|\mathrm{\Delta }|`$ be a support map on a fan $`\mathrm{\Delta }`$. If there are linear forms $`u_\sigma `$, $`\sigma \mathrm{\Delta }`$, on $`N`$ such that for any pair $`\sigma ,\tau \mathrm{\Delta }`$ we have
$$h|_\sigma =u_\sigma |_\sigma ,h|_\tau u_\sigma |_\tau $$
then $`h`$ is a convex support map on $`\mathrm{\Delta }`$. Conversely, if $`\mathrm{\Delta }`$ is complete and $`h`$ is convex then $`h`$ or $`h`$ satisfies the above condition.
On noncomplete fans, the concept of convexity for a support function via the above inequalities is more restrictive than our concept:
###### Example 2.4.
Consider the fan $`\mathrm{\Delta }`$ in $`^2`$ generated by the two maximal cones
$$\sigma _1:=\mathrm{cone}((1,0),(1,1)),\sigma _2:=\mathrm{cone}((0,1),(1,1))$$
and the support map $`h:|\mathrm{\Delta }|`$ determined by
$$h(v_1,v_2):=\{\begin{array}{cc}2v_1+2v_2& \text{if }(v_1,v_2)\sigma _1,\hfill \\ v_1+v_2& \text{if }(v_1,v_2)\sigma _2.\hfill \end{array}$$
Then $`h`$ is convex: The convex hull $`\gamma `$ of the graph $`\mathrm{\Gamma }_h`$ is a strictly convex cone with four rays, namely
$$\gamma =\mathrm{cone}((1,0,2),(1,1,0),(0,1,1),(1,1,0)).$$
Moreover, the maximal cones of $`\mathrm{\Lambda }_h`$ are precisely the two faces of $`\gamma `$ above $`\sigma _1`$ and $`\sigma _2`$ respectively.
However neither the function $`h`$ nor the function $`h`$ satisfies the inequalities of Remark 2.3, because we have:
$$h((0,1))=1<2,h((1,1))=0>2.$$
In order to define the notion of strict convexity, we have to note some observations on convex support maps. The first one is:
###### Lemma 2.5.
If the support map $`h:|\mathrm{\Delta }|^k`$ is convex, then the projected cones $`P(\delta )`$, $`\delta \mathrm{\Lambda }_h`$, form a quasifan $`\mathrm{\Sigma }_h`$ in the lattice $`N`$.
###### Proof.
The projection $`P`$ is injective on any given $`\delta \mathrm{\Lambda }_h`$, and hence induces a bijection between the faces of $`\delta `$ and the faces of $`P(\delta )`$. Moreover, given $`\delta _1,\delta _2\mathrm{\Lambda }_h`$, injectivity of $`P`$ on $`|\mathrm{\Lambda }_h|`$ implies
$$P(\delta _1)P(\delta _2)=P(\delta _1\delta _2).$$
Since $`\delta _1\delta _2`$ is a face of both $`\delta _i`$, the above consideration yields that $`P(\delta _1\delta _2)`$ is a common face of $`P(\delta _1)`$ and $`P(\delta _2)`$. ∎
If $`h:|\mathrm{\Delta }|^k`$ is a convex support map, then we call $`\mathrm{\Sigma }_h`$ the quasifan associated to $`h`$. We need the following properties of this quasifan:
###### Lemma 2.6.
Let $`\mathrm{\Sigma }_h`$ be the quasifan associated to a convex support map $`h:|\mathrm{\Delta }|^k`$. Then we have:
1. Every cone of $`\mathrm{\Delta }`$ is contained in a cone of $`\mathrm{\Sigma }_h`$.
2. Every cone $`\sigma \mathrm{\Sigma }_h`$ is generated by the cones $`\tau \mathrm{\Delta }`$ with $`\tau \sigma `$.
###### Definition 2.7.
We say that a convex support map $`h:|\mathrm{\Delta }|^k`$ is strictly convex if its associated quasifan $`\mathrm{\Sigma }_h`$ equals $`\mathrm{\Delta }`$.
Using Remark 2.3, one verifies that on a complete fan $`\mathrm{\Delta }`$, our notion of strict convexity for a support map $`h:|\mathrm{\Delta }|`$ concides with the usual one, as defined in \[12, p. 67\]. Again, for noncomplete fans the notions differ:
###### Example 2.8.
The convex support map $`h:|\mathrm{\Delta }|`$ of Example 2.4 is even strictly convex.
We now come to the announced main result of this section, namely the characterization of divisoriality of a toric variety via existence of a strictly convex support map:
###### Proposition 2.9.
For a fan $`\mathrm{\Delta }`$ in a lattice $`N`$, the following statements are equivalent:
1. $`\mathrm{\Delta }`$ admits a strictly convex support map.
2. The toric variety $`X`$ associated to $`\mathrm{\Delta }`$ is divisorial.
In the proof of this statement, we make use of the following wellknown characterization of existence of geometric quotients for subtorus actions in terms of fans, see e.g. \[13, Theorem 5.1\]:
###### Proposition 2.10.
Let $`\widehat{\mathrm{\Delta }}`$ be a fan in a lattice $`\widehat{N}`$ with associated toric variety $`\widehat{X}`$, let $`P:\widehat{N}N`$ be a surjective lattice homomorphism, and let $`H`$ be the subtorus of the big torus of $`\widehat{X}`$ corresponding to $`\mathrm{ker}(P)`$. The following statements are equivalent:
1. $`P`$ is injective on the support $`|\widehat{\mathrm{\Delta }}|`$.
2. The action of $`H`$ on $`\widehat{X}`$ has a geometric quotient.
If one of these statements holds, then the quotient variety $`\widehat{X}/H`$ is the toric variety determined by the fan $`\{P(\sigma );\sigma \widehat{\mathrm{\Delta }}\}`$ in $`N`$.
###### Proof of Proposition 2.9.
Assume first that the fan $`\mathrm{\Delta }`$ admits a strictly convex support map $`h:|\mathrm{\Delta }|^k`$. Then since $`\mathrm{\Delta }=\mathrm{\Sigma }_h`$, all cones of $`\mathrm{\Sigma }_h`$ are strictly convex. As before, let $`\widehat{N}:=N\times ^k`$. By convexity of $`h`$, the projection $`P:\widehat{N}_{}N_{}`$ is an injection on $`|\mathrm{\Lambda }_h|`$. In particular, all cones of $`\mathrm{\Lambda }_h`$ are strictly convex. That means that $`\mathrm{\Lambda }_h`$ is a fan.
The toric variety $`\widehat{X}`$ associated to $`\mathrm{\Lambda }_h`$ is quasiaffine, and the projection $`P:\widehat{N}N`$ gives rise to a toric morphism $`p:\widehat{X}X`$. According to Proposition 2.10, this toric morphism $`p`$ is a geometric quotient for the subtorus action on $`\widehat{X}`$ corresponding to $`\mathrm{ker}(P)\widehat{N}`$. Thus, Corollary 1.6 yields that $`X`$ is divisorial.
Suppose now that the toric variety $`X`$ determined by the fan $`\mathrm{\Delta }`$ is divisorial. By Corollary 1.6, there is a quasiaffine toric variety $`\widehat{X}`$ and a toric morphism $`p:\widehat{X}X`$ such that $`H:=\mathrm{ker}(p)`$ is a subtorus of the big torus of $`\widehat{X}`$ and $`p`$ is a geometric quotient for the action of $`H`$ on $`\widehat{X}`$.
Let $`p:\widehat{X}X`$ arise from a map $`P:\widehat{N}N`$ of fans $`\widehat{\mathrm{\Delta }}`$ and $`\mathrm{\Delta }`$. Since $`H=\mathrm{ker}(p)`$ is connected, the map $`P`$ is surjective and we obtain a section $`N\widehat{N}`$ for $`P`$. So we may assume that $`\widehat{N}=N\times ^k`$ holds and that $`P`$ is the projection onto the first factor. By the above Proposition 2.10, the projection $`P`$ is injective on $`|\widehat{\mathrm{\Delta }}|`$. Thus, for each $`\widehat{\sigma }\widehat{\mathrm{\Delta }}`$, the restriction
$$P|_{\widehat{\sigma }}:\widehat{\sigma }\sigma :=P(\widehat{\sigma })$$
admits a uniquely determined linear inverse of the form $`g_\sigma =(\mathrm{id}_N_{},h_\sigma )`$. The maps $`h_\sigma :\sigma ^k`$ patch together to a support map $`h`$ on $`\mathrm{\Delta }`$. By construction, $`\mathrm{\Lambda }_h=\widehat{\mathrm{\Delta }}`$ and $`\mathrm{\Sigma }_h=\mathrm{\Delta }`$. So $`h`$ is the desired strictly convex support map on $`\mathrm{\Delta }`$. ∎
In the remainder of this section we show that convex support maps in a canonical way define toric morphisms to divisorial toric varieties. Let $`\mathrm{\Delta }`$ be a fan in a lattice $`N`$, and let $`h:|\mathrm{\Delta }|^k`$ be a convex support map.
There is a universal method to construct a fan from the associated quasifan $`\mathrm{\Sigma }_h`$: Let $`\sigma _{\mathrm{min}}\mathrm{\Sigma }_h`$ denote its minimal cone. This is a linear subspace of $`N_{}`$. Let $`N_0:=\sigma _{\mathrm{min}}N`$, set $`N_h:=N/N_0`$, and denote by $`F_h:NN_h`$ the projection. The quotient fan of $`\mathrm{\Sigma }_h`$ is the fan
$$\mathrm{\Delta }_h:=\{F_h(\sigma );\sigma \mathrm{\Sigma }_h\}.$$
The projection $`F_h:NN_h`$ is a map of the quasifans $`\mathrm{\Sigma }_h`$ and $`\mathrm{\Delta }_h`$. Moreover, $`F_h`$ is universal in the sense that every map of quasifans from $`\mathrm{\Sigma }_h`$ to a fan $`\mathrm{\Delta }^{}`$ factors uniquely through $`F_h`$.
Now, let $`X`$ and $`X_h`$ denote the toric varieties associated to the fans $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }_h`$ respectively. Our precise statement is the following:
###### Proposition 2.11.
The toric variety $`X_h`$ is divisorial, and the projection $`F_h`$ induces a toric morphism $`f_h:XX_h`$.
###### Proof.
By Lemma 2.6 i) and the universal property of the quotient fan $`\mathrm{\Delta }_h`$, the projection $`F_h:NN_h`$ is a map of the fans $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }_h`$ and hence induces a toric morphism $`f_h:XX_h`$. So we only have to show that $`X_h`$ is divisorial. In view of Proposition 2.9, we look for a strictly convex support map an $`\mathrm{\Delta }_h`$.
The first step is to construct a strictly convex support map $`g`$ on the quasifan $`\mathrm{\Sigma }_h`$ associated to $`h`$: Consider a cone $`\sigma \mathrm{\Sigma }_h`$. Then, as earlier denoting by $`P:\widehat{N}N`$ the projection, we have $`\sigma =P(\delta )`$ for some cone $`\delta \mathrm{\Lambda }_h`$.
By convexity of $`h`$, the restriction $`P:\delta \sigma `$ has an inverse of the form $`(\mathrm{id},g_\sigma )`$. The maps $`g_\sigma `$ patch together to a support map $`g`$ on $`\mathrm{\Sigma }_h`$, and $`g`$ extends $`h`$. Moreover, $`\mathrm{\Gamma }_g`$ equals $`\mathrm{\Lambda }_h`$ and hence the quasifan associated to $`g`$ coincides with $`\mathrm{\Sigma }_h`$.
Note that $`\mathrm{\Sigma }_g=\mathrm{\Sigma }_h`$ does not change if we add a global linear function to $`g`$. So we may assume that the support function $`g`$ vanishes on the minimal cone of $`\mathrm{\Sigma }_g`$. But then we can push down $`g`$ to a strictly convex support function on the quotient fan $`\mathrm{\Delta }_h`$. ∎
## 3. Toric divisorial reduction
Fix a toric variety $`X`$. In , we presented a universal way to reduce $`X`$ to a quasiprojective toric variety. In this section we give an analogous construction, that reduces to divisorial toric varieties.
###### Definition 3.1.
A toric divisorial reduction of $`X`$ is a toric morphism $`r:XX^{\mathrm{tdr}}`$ to a divisorial toric variety $`X^{\mathrm{tdr}}`$ such that every toric morphism $`f:XZ`$ to a divisorial toric variety $`Z`$ has a unique factorization $`f=\stackrel{~}{f}r`$ with a toric morphism $`\stackrel{~}{f}:X^{\mathrm{tdr}}Z`$.
###### Theorem 3.2.
Every toric variety admits a toric divisorial reduction.
The proof is given below. We need the following statement on the pullback of a convex support map:
###### Lemma 3.3.
Let $`F:NN^{}`$ be a map of fans $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }^{}`$ in lattices $`N`$ and $`N^{}`$ respectively. If $`h^{}:|\mathrm{\Delta }^{}|^k`$ is a convex support map on $`\mathrm{\Delta }^{}`$, then $`h:=h^{}F`$ is a convex support map on $`\mathrm{\Delta }`$ and $`F`$ is a map of the associated quasifans $`\mathrm{\Sigma }_h`$ and $`\mathrm{\Sigma }_h^{}`$.
###### Proof.
Clearly $`h`$ is a support map on $`\mathrm{\Delta }`$. To prove convexity of $`h`$, we consider the filled graphs $`\mathrm{\Lambda }_h`$, $`\mathrm{\Lambda }_h^{}`$ and the map
$$\widehat{F}:=F\times \mathrm{id}_^k:N\times ^kN^{}\times ^k.$$
We claim that $`\widehat{F}`$ is a map of the quasifans $`\mathrm{\Lambda }_h`$ and $`\mathrm{\Lambda }_h^{}`$. To verify this, note first that $`\widehat{F}`$ maps the graph $`\mathrm{\Gamma }_h`$ to $`\mathrm{\Gamma }_h^{}`$. Let $`\delta \mathrm{\Lambda }_h`$. We have to show that the minimal face $`\delta ^{}`$ of $`\mathrm{conv}(\mathrm{\Gamma }_h^{})`$ containing $`\widehat{F}(\delta )`$ belongs to $`\mathrm{\Lambda }_h^{}`$. Let
$$G:=\mathrm{id}_{|\mathrm{\Delta }|}\times h,G^{}:=\mathrm{id}_{|\mathrm{\Delta }^{}|}\times h^{}.$$
By definition of $`\mathrm{\Lambda }_h`$, the relative interior $`\delta ^{}`$ of $`\delta `$ contains a point of the graph of $`h`$, i.e. a point of the form $`G(v)`$ for some $`v|\mathrm{\Delta }|`$. By the choice of $`\delta ^{}`$ this means $`\widehat{F}(G(v))(\delta ^{})^{}`$. On the other hand, by definition of $`G`$, $`G^{}`$ and $`\widehat{F}`$ we have
$$\widehat{F}(G(v))=G^{}(F(v))\mathrm{\Gamma }_h^{}.$$
Hence $`\mathrm{\Gamma }_h^{}(\delta ^{})^{}\mathrm{}`$. This implies $`\delta ^{}\mathrm{\Lambda }_h^{}`$, and our claim is proved.
For convexity of $`h`$, we have to show that the projection $`P:N\times ^kN`$ is injective on $`|\mathrm{\Lambda }_h|`$. Suppose $`w_i=(v_i,t_i)|\mathrm{\Lambda }_h|`$ are two points such that $`P(w_1)`$ equals $`P(w_2)`$, that means $`v_1=v_2`$. Then we have
$$P^{}(\widehat{F}(w_1))=P^{}(\widehat{F}(w_2)),$$
where $`P^{}:N^{}\times ^kN^{}`$ is the projection. Since $`\widehat{F}`$ is a map of the quasifans $`\mathrm{\Lambda }_h`$ and $`\mathrm{\Lambda }_h^{}`$ and $`P^{}`$ is injective on $`|\mathrm{\Lambda }_h^{}|`$, this implies $`\widehat{F}(w_1)=\widehat{F}(w_2)`$. In particular, we have $`t_1=t_2`$ and thus $`w_1=w_2`$.
Finally, the fact that $`F`$ is a map of the quasifans $`\mathrm{\Sigma }_h^{}`$ and $`\mathrm{\Sigma }_h`$ follows immediately from the fact that $`\widehat{F}`$ is a map of the quasifans $`\mathrm{\Lambda }_h^{}`$ and $`\mathrm{\Lambda }_h`$. ∎
###### Proof of Theorem 3.2.
Let $`X`$ be a toric variety arising from a fan $`\mathrm{\Delta }`$ in a lattice $`N`$. First we show that any given toric morphism $`f:XZ`$ from $`X`$ to a divisorial variety $`Z`$ factors uniquely through one of the toric morphisms $`f_h`$ arising from a convex support map on $`\mathrm{\Delta }`$ as in Proposition 2.11.
To see this, consider the map of fans $`F:\mathrm{\Delta }\mathrm{\Delta }^{}`$ associated to the given toric morphism $`f`$ and choose a strictly convex support map $`h^{}`$ on $`\mathrm{\Delta }^{}`$. Lemma 3.3 tells us that by pulling back $`h^{}`$ via $`F`$, we obtain a convex support map $`h`$ on $`\mathrm{\Delta }`$. Moreover, $`F`$ defines a map of quasifans from $`\mathrm{\Sigma }_h`$ to $`\mathrm{\Sigma }_h^{}=\mathrm{\Delta }^{}`$.
Now, the map of fans $`F`$ factors as a map of fans through the projection $`F_h:NN_h`$, i.e., $`F`$ induces a map from the quotient fan $`\mathrm{\Delta }_h`$ of $`\mathrm{\Sigma }_h`$ to $`\mathrm{\Delta }^{}`$. Obviously, the corresponding toric morphism is the desired factorization of $`f:XZ`$ through $`f_h:XX_h`$.
Now let us take a closer look at the toric morphisms $`f_h:XX_h`$ arising from convex support maps. Recall that the morphism $`f_h`$ is already determined by the quasifan $`\mathrm{\Sigma }_h`$ associated to $`h`$. By Lemma 2.6 ii), each such quasifan has the property that all cones are generated by cones of $`\mathrm{\Delta }`$. Consequently there exist only finitely many of such quasifans, say $`\mathrm{\Sigma }_1,\mathrm{},\mathrm{\Sigma }_r`$.
Let $`f_i:XY_i`$ denote the toric morphisms to divisorial toric varieties determined by $`\mathrm{\Sigma }_i`$, and consider their product $`f:=f_1\times \mathrm{}\times f_r`$. Let $`Y`$ denote the closure of the image $`f(X)`$ in $`Y_1\times \mathrm{}\times Y_r`$. The normalization $`\stackrel{~}{Y}`$ of $`Y`$ is again a divisorial toric variety, and $`f`$ lifts to a toric morphism to $`\stackrel{~}{Y}`$. In $`\stackrel{~}{Y}`$ we choose the smallest open toric subvariety $`Y^{}`$ containing the image of $`f`$, and restricting $`f`$, we obtain a toric morphism $`r:XY^{}`$.
By construction, for every $`i`$ we have a unique factorization of $`f_i`$ through $`r`$, namely $`f_i=\mathrm{pr}_ir`$, where $`\mathrm{pr}_i:Y^{}Y_i`$ denotes the restriction of the projection on the $`i`$-th factor. This proves that $`r`$ is the desired toric divisorial reduction. ∎
We conclude this section with some examples. Note that any two-dimensional toric variety is simplicial and hence divisorial. So the minimal dimension for interesting examples is $`3`$.
###### Example 3.4.
If a toric variety does not admit nontrivial effective Cartier divisors, see e.g. \[12, p. 25\], then its toric divisorial reduction is a point.
###### Example 3.5.
Consider the following eight vectors in $`^3`$:
$$\begin{array}{ccccccc}v_1:=(2,2,1),\hfill & & v_2:=(2,2,1),\hfill & & v_3:=(2,2,1),\hfill & & v_4:=(2,2,1),\hfill \\ v_5:=(1,1,1),\hfill & & v_6:=(1,1,1),\hfill & & v_7:=(1,1,1),\hfill & & v_8:=(2/3,1/3,1).\hfill \end{array}$$
Let $`\mathrm{\Delta }`$ denote the fan in $`^3`$ with maximal cones
$$\begin{array}{ccc}\sigma _1:=\mathrm{cone}(v_1,v_2,v_5,v_6),\hfill & & \sigma _2:=\mathrm{cone}(v_2,v_3,v_6,v_7),\hfill \\ \sigma _3:=\mathrm{cone}(v_3,v_4,v_7,v_8),\hfill & & \sigma _4:=\mathrm{cone}(v_1,v_4,v_5,v_8),\hfill \\ \sigma _5:=\mathrm{cone}(v_5,v_6,v_7,v_8).\hfill & & \end{array}$$
Intersection of $`\mathrm{\Delta }`$ with the plane $`x_3=1`$.
The identity on $`^3`$ defines a map of fans from $`\mathrm{\Delta }`$ to the fan of faces $`𝔉(\sigma )`$ of the cone $`\sigma :=\mathrm{cone}(v_1,v_2,v_3,v_4)`$. We claim that the corresponding toric morphism $`r:X_\mathrm{\Delta }X_\sigma `$ is the toric divisorial reduction of $`X_\mathrm{\Delta }`$.
To see this, consider a convex support map $`h:|\mathrm{\Delta }|^k`$, and its associated quasifan $`\mathrm{\Sigma }_h`$. Lemma 2.6 implies that we have only two possibilities, namely $`\mathrm{\Sigma }_h=𝔉(\sigma )`$ or $`\mathrm{\Sigma }_h=\mathrm{\Delta }`$. Thus, to verify our claim, we only have to exclude the latter possibility, i.e., we have to show that $`h`$ cannot be strictly convex.
Otherwise, let $`\delta _5\mathrm{\Lambda }_h`$ be the maximal cone above $`\sigma _5`$ and choose a linear form $`\lambda :N_{}\times ^k`$ that is nonnegative on $`\gamma :=\mathrm{conv}(\mathrm{\Gamma }_h)`$ and fulfills $`\delta _5=\gamma \lambda ^{}`$. Pulling back $`\lambda `$ via $`\mathrm{id}_N\times h`$, we obtain a nonnegative support function $`g`$ on $`\mathrm{\Delta }`$ vanishing precisely on $`\sigma _5`$. Note that
$$g(v_1)=g(v_2)=g(v_3).$$
Moreover, we have the relations
$$v_4=17v_328v_7+12v_8,v_4=5v_116v_5+12v_8.$$
Applying $`g`$, we obtain $`17g(v_3)=5g(v_1)`$. This contradicts $`g(v_1)=g(v_3)`$. So, $`h`$ cannot be strictly convex and our claim is proved.
###### Example 3.6.
We describe a toric variety with a nonsurjective toric divisorial reduction. Similarly to the preceding example, consider the vectors
$$\begin{array}{ccccc}v_1:=(2,2,1,0),\hfill & & v_2:=(2,2,1,0),\hfill & & v_3:=(2,2,1,0),\hfill \\ v_4:=(2,2,1,0),\hfill & & v_5:=(1,1,1,0),\hfill & & v_6:=(1,1,1,0),\hfill \\ v_7:=(1,1,1,0),\hfill & & v_8:=(2/3,1/3,1,0).\hfill & & \end{array}$$
in $`^4`$ and let furthermore $`e_4`$ be the fourth canonical base vector. Let $`\mathrm{\Delta }`$ denote the fan in $`^4`$ with maximal cones
$$\begin{array}{ccc}\sigma _1:=\mathrm{cone}(v_1,v_2,v_5,v_6),\hfill & & \sigma _2:=\mathrm{cone}(v_2,v_3,v_6,v_7),\hfill \\ \sigma _3:=\mathrm{cone}(v_3,v_4,v_7,v_8),\hfill & & \sigma _4:=\mathrm{cone}(v_1,v_4,v_5,v_8),\hfill \\ \sigma _5:=\mathrm{cone}(v_5,v_6,v_7,v_8).\hfill & & \sigma _6:=\mathrm{cone}(v_5,v_6,e_4)\hfill \end{array}$$
Intersection of $`\mathrm{\Delta }`$ with the hyperplane $`x_3=1`$.
The identity on $`^4`$ defines a map of fans from $`\mathrm{\Delta }`$ to the fan of faces $`𝔉(\sigma )`$ of the cone $`\sigma :=\mathrm{cone}(v_1,v_2,v_3,v_4,e_4)`$. We claim that the corresponding toric morphism $`r:X_\mathrm{\Delta }X_\sigma `$ is the toric divisorial reduction of $`X_\mathrm{\Delta }`$. Note that this map is not surjective.
Let us verify the claim. If $`h`$ is a convex support map it follows that $`|\mathrm{\Sigma }_h|\sigma `$. The restriction of $`h`$ to the support of the subfan $`\mathrm{\Delta }^{}`$ of $`\mathrm{\Delta }`$ generated by the cones $`\sigma _1,\mathrm{},\sigma _5`$ defines a convex support map $`h^{}`$ of $`\mathrm{\Delta }^{}`$. So by the previous example, $`\mathrm{\Sigma }_h^{}=𝔉(\sigma ^{})`$, where $`\sigma ^{}`$ denotes the cone generated by $`v_1,\mathrm{},v_4`$.
Now Lemma 3.3 implies that the smallest cone $`\tau `$ in $`\mathrm{\Sigma }_h`$ containing $`\sigma _5`$ also contains all of $`\sigma ^{}`$. That means by Lemma 2.6 that either $`\tau =\sigma ^{}`$ or $`\tau =\sigma `$. In any case, since $`\sigma ^{}`$ is a face of $`\sigma `$ we obtain $`\sigma ^{}\mathrm{\Sigma }_h`$.
Next consider the smallest cone $`\tau ^{}\mathrm{\Sigma }_h`$ containing $`\sigma _6`$. We have $`v_5,v_6\sigma _6`$, so the cone $`\tau ^{}`$ meets $`\sigma ^{}`$ in its relative interior. Since $`\mathrm{\Sigma }_h`$ is a quasifan, we can conclude that $`\sigma ^{}`$ is in fact a face of $`\tau ^{}`$. Because $`e_4\tau `$ this implies $`\tau ^{}=\sigma `$, and we obtain $`\mathrm{\Sigma }_h=𝔉(\sigma )`$.
## 4. A Lifting Lemma
Here we relate regular maps between divisorial toric prevarieties to regular maps between quasiaffine toric varieties. For maps of projective spaces, this is a classical observation:
###### Example 4.1.
Let $`f:_n_m`$ be a regular map of projective spaces. Then $`f`$ is of the form
$$[z_0,\mathrm{},z_n][f_0(z_0,\mathrm{},z_n),\mathrm{},f_m(z_0,\mathrm{},z_n)]$$
with homogeneous polynomials $`f_i`$ that are pairwise of the same degree. In other words, there is a lifting
The main result of this section is the following generalization of the above lifting statement:
###### Lemma 4.2.
Let $`f:X_1X_2`$ be a regular map of divisorial toric prevarieties such that $`f(X_1)`$ intersects the big torus of $`X_2`$. Then there exists a commutative diagram
where $`\widehat{X}_1`$, $`\widehat{X}_2`$ are quasiaffine toric varieties, $`q_i:\widehat{X}_iX_i`$ are geometric prequotients for free subtorus actions on $`\widehat{X}_i`$ and $`\widehat{f}:\widehat{X}_1\widehat{X}_2`$ is a regular map.
###### Proof.
We use the ideas and methods presented in \[14, Section 2\]. Choose effective $`T_i`$-invariant Cartier divisors $`D_1^i,\mathrm{},D_{r_i}^i`$ on $`X_i`$ such that the complements $`X_i\mathrm{Supp}(D_j^i)`$ form an affine cover of $`X_i`$. Let $`W_i\mathrm{CDiv}(X_i)`$ denote the subgroup generated by $`D_1^i,\mathrm{},D_{r_i}^i`$. The pullback via $`f`$ gives rise to a group homomorphism
$$\psi :W_2\mathrm{CDiv}(X_1),Df^{}(D).$$
Enlarge $`W_1`$ by adding the image $`\psi (W_2)`$. Note that the line bundles determined by the divisors of $`W_i`$ are $`T_i`$-linearizable, see \[17, p. 67, Remark\]. We shall regard $`\psi `$ in the sequel as a homomorphism from $`W_2`$ to $`W_1`$. Consider the $`𝒪_{X_i}`$-algebras
$$𝒜_i:=\underset{DW_i}{}𝒪_D(X_i)$$
and their associated relative spectra $`\widehat{X}_i:=\mathrm{Spec}(𝒜_i)`$. By \[14, Remark 2.1\], the inclusion $`𝒪_{X_i}𝒜_i`$ gives rise to a geometric prequotient $`q_i:\widehat{X}_iX_i`$ for the free action of the algebraic torus $`H_i:=\mathrm{Spec}(𝕂[W_i])`$ on $`\widehat{X}_i`$ induced by the $`W_i`$-grading of $`𝒜_i`$.
Since $`W_1`$ and $`W_2`$ define ample groups of line bundles in the sense of \[14, Definition 2.2\], each $`\widehat{X}_i`$ is in fact a quasiaffine variety. Moreover, by \[14, Proposition 2.3\], the variety $`\widehat{X}_i`$ carries a regular action of the algebraic torus $`T_i`$ commuting with the action of $`H_i`$ such that $`q_i:\widehat{X}_iX_i`$ becomes $`T_i`$-equivariant. It follows that $`\widehat{X}_i`$ is a toric variety with big torus $`\widehat{T}_i=T_i\times H_i`$.
We still have to construct the lifting $`\widehat{f}:\widehat{X}_1\widehat{X}_2`$. As to this, note that for every affine open subset $`UX_2`$, we obtain a homomorphism of $`W_i`$-graded algebras by setting
$$𝒜_2(U)𝒜_1(f^1(U)),𝒪_D(U)hf^{}(h)𝒪_{\psi (D)}(U)(DW_2).$$
Note that on the homogeneous component $`𝒜_2(U)_0`$, this is just the comorphism of the map $`f`$. By definition of $`\widehat{X}_i`$ and the maps $`q_i:\widehat{X}_iX_i`$, each of the above homomorphisms gives rise to a lifting
$$\widehat{f}_U:q_1^1(f^1(U))q_2^1(U)$$
of the restriction $`f:f^1(U)U`$. By construction, the maps $`\widehat{f}_U`$ patch together to the desired lifting $`\widehat{f}:\widehat{X}_1X_2`$ of $`f:X_1X_2`$. ∎
The following observation will be needed later to obtain equivariance properties for the lifting $`\widehat{f}:\widehat{X}_1\widehat{X}_2`$ constructed in the above Lemma.
###### Lemma 4.3.
For $`i=1,2`$, let $`T_i`$ be algebraic tori and let $`Y_i`$ be irreducible $`T_i`$-varieties such that $`T_2`$ acts freely on $`Y_2`$. If $`f:Y_1Y_2`$ is regular and maps the orbits of $`T_1`$ into orbits of $`T_2`$, then there is a homomorphism $`\phi :T_1T_2`$ such that $`f(tx)=\phi (t)f(x)`$ holds for all $`(t,x)T_1\times Y_1`$.
###### Proof.
By Sumihiro’s Theorem \[20, Corollary 2\], we may assume that $`Y_2`$ is affine. Thus, there is an algebraic quotient $`Y_2Y`$ for the action of $`T_2`$ on $`Y_2`$. Since $`T_2`$ acts freely, the quotient map $`Y_2Y`$ is equivariantly locally trivial. Thus, shrinking $`Y`$, we may even assume that $`Y_2=T_2\times Y`$ holds. In particular, one has $`f=(f_1,f_2)`$ with regular maps $`f_1:Y_1T_2`$ and $`f_2:Y_1Y`$. So, we obtain a regular map
$$\mathrm{\Phi }:T_1\times Y_1T_2,(t,x)f_1(tx)f_1(x)^1.$$
For fixed $`xY_1`$, the map $`t\mathrm{\Phi }(t,x)`$ maps the neutral element of $`T_1`$ to the neutral element of $`T_2`$ and hence is necessarily a homomorphism of the tori $`T_1`$ and $`T_2`$. By rigidity of tori \[9, III.8.10\], the map $`\mathrm{\Phi }`$ does not depend on $`x`$. So there is a homomorphism $`\phi :T_1T_2`$ with $`\mathrm{\Phi }(t,x)=\phi (t)`$ for all $`(t,x)T_1\times Y_1`$. Clearly, $`\phi `$ is as desired. ∎
A different aspect of the lifting problem is discussed extensively in : Given two quotient presentations $`\widehat{X}_iX_i`$ of toric varieties in the sense of and a regular map $`f:X_1X_2`$, when can this map be lifted to a regular map $`F:\widehat{X}_1\widehat{X}_2`$?
## 5. Decomposition of regular maps
Let $`X`$ be a toric variety with big torus $`T`$ and consider the action of a closed subgroup $`HT`$ on $`X`$. Here we provide the key to relate $`H`$-invariant regular maps $`XY`$ to $`H`$-invariant toric morphisms:
###### Lemma 5.1.
Let $`f:XY`$ be an $`H`$-invariant regular map to a divisorial variety $`Y`$. Then there exists a dominant $`H`$-invariant toric morphism $`g:XX^{}`$ to a divisorial toric variety $`X^{}`$, an open subset $`UX^{}`$ with $`g(X)U`$ and a regular map $`h:UY`$ such that $`f=hg`$.
###### Proof.
First we reduce the problem to the case that $`H`$ is connected. Suppose that $`g:XX^{}`$ and $`h:UY`$ satisfy the assertion for the identity component $`H^0`$ of $`H`$. Then $`g`$ induces an action of the finite abelian group $`\mathrm{\Gamma }:=H/H^0`$ on $`X^{}`$. Let $`p:X^{}X^{\prime \prime }`$ be the geometric quotient for this action. Note that $`p`$ is a toric morphism. Using Corollary 1.6, we see that the variety $`X^{\prime \prime }`$ is again divisorial.
By appropriate shrinking, we achieve that $`U`$ is $`\mathrm{\Gamma }`$-invariant. Since $`p`$ is geometric, $`p(U)`$ is open in $`X^{\prime \prime }`$ and the restriction $`p:Up(U)`$ is again a geometric quotient for the action of $`\mathrm{\Gamma }`$. Since $`h`$ is $`\mathrm{\Gamma }`$-invariant, we have $`h=h^{}p`$ for some regular map $`h^{}:p(U)Y`$. It follows that $`f=h^{}(pg)`$ is the desired decomposition. Consequently, it suffices to give the proof for connected $`H`$.
The next simplification provides the link to the toric setting: As mentioned before, we can realize $`Y`$ as a closed subvariety of a smooth toric prevariety $`Z`$ of affine intersection, see 1.2. Let $`Z^{}Z`$ denote the minimal orbit closure of the big torus of $`Z`$ such that $`f(X)Z^{}`$ holds. Then $`Z^{}`$ is again a smooth toric prevariety of affine intersection, but in $`Z^{}`$ the image $`f(X)`$ intersects the big torus.
Now, for the moment regard $`f`$ as a map from $`X`$ to $`Z^{}`$ and suppose that $`g:XX^{}`$ and $`h:UZ^{}`$ satisfy the assertion for $`f:XZ^{}`$. Taking closures in $`U`$ and $`Z^{}`$ respectively, we obtain
$$h(U)h\left(\overline{g(X)}\right)\overline{h(g(X))}=\overline{f(X)}Y.$$
That means $`h`$ is in fact a map from $`U`$ to $`Y`$. Thus $`X^{}`$, $`g`$, $`h`$ and $`U`$ also provide the desired data for the original $`f:XY`$. Consequently, we can assume in the sequel that $`Y`$ is a smooth toric prevariety of affine intersection and that $`f(X)`$ intersects the big torus of $`Y`$. But then according to Lemma 4.2 there is a commutative diagram
where $`\widehat{X}`$, $`\widehat{Y}`$ are quasiaffine toric varieties and the vertical maps are geometric prequotients for free actions of subtori $`H_X`$ and $`H_Y`$ of the big tori of $`\widehat{X}`$ and $`\widehat{Y}`$ respectively. We may even assume that $`\widehat{X}=X`$ holds:
Let $`H^{}:=p^1(H)`$ and suppose that the $`H^{}`$-invariant regular map $`f^{}:=fp`$ admits a decomposition of the form $`f^{}=h^{}g^{}`$ with a dominant $`H^{}`$-invariant toric morphism $`g^{}:\widehat{X}X^{}`$ and a regular map $`h^{}:UY`$ defined on an open neighbourhood $`U`$ of the image of $`g^{}`$.
Then, by the universal property of $`p`$, there is a toric morphism $`g:XX^{}`$ with $`g^{}=gp`$. Clearly this morphism is dominant. Moreover, since $`p`$ is surjective, it is $`H`$-invariant and $`g(X)U`$ holds. Consequently, $`f=h^{}g`$ is a decomposition as wanted. So it suffices to prove the assertion for the case that $`\widehat{X}=X`$ and $`H_X=1`$ hold and $`p`$ is the identity map.
Now we consider the regular map $`\widehat{f}:X\widehat{Y}`$ as a map from an $`H`$-variety to an $`H_Y`$-variety. Since $`q\widehat{f}=f`$ is $`H`$-invariant, every $`H`$-orbit is mapped by $`\widehat{f}`$ into a fiber of $`q`$. On the other hand, the fibers of $`q`$ are precisely the $`H_Y`$-orbits. So we can apply Lemma 4.3 and conclude that $`\widehat{f}`$ is $`H`$-equivariant with respect to a homomorphism $`HH_Y`$.
Choosing a locally closed toric embedding $`\widehat{Y}^s`$, we obtain a homomorphism $`H_Y^s`$, and the induced map $`\widehat{f}:X^s`$ is $`H`$-equivariant with respect to the homomorphism $`HH_Y^s`$. So the components of $`\widehat{f}`$ are $`H`$-homogeneous regular functions. By writing the components of $`\widehat{f}`$ as linear combinations of character functions of the big torus $`TX`$, and using the summands to define a toric morphism $`g^{}:X^r`$, we obtain a decomposition of $`\widehat{f}`$ in the form $`\widehat{f}=sg^{}`$, with a linear map $`s:^r^s`$. Note that $`g^{}`$ induces an action of $`H`$ on $`^r`$ making $`s:^r^s`$ into an $`H`$-equivariant map.
Let $`W`$ be the normalization of the closure of $`g^{}(X)`$ in $`^r`$. Then $`W`$ is an affine toric variety with big torus $`g^{}(T)`$. We can lift $`g^{}`$ to a dominant toric morphism $`\widehat{g}:XW`$, and pull back $`s`$ to a regular map $`\widehat{s}:W^s`$. Both, $`\widehat{g}`$ and $`\widehat{s}`$, are again equivariant for the induced $`H`$-action on $`W`$. The set $`V:=\widehat{s}^1(\widehat{Y})`$ is $`H`$-invariant and open in $`W`$. Moreover, we have $`\widehat{g}(X)V`$. So far, we are in the following situation:
Since $`\widehat{s}:W^s`$ is an affine map, also its restriction $`\widehat{s}:V\widehat{Y}`$ is affine. Thus $`q\widehat{s}:VY`$ is an affine $`H`$-invariant regular map. Existence of an affine $`H`$-invariant map $`VY`$ already implies existence of a good quotient $`p:VV//H`$ for the action of $`H`$, see e.g. \[19, Prop. 3.12\]. So we obtain the following commutative diagram of regular maps:
Note that $`g:=p\widehat{g}:XV//H`$ is $`H`$-invariant and $`V//H`$ is divisorial, because $`Y`$ is divisorial and $`h`$ is an affine morphism. So the decomposition $`f=hg`$ is almost as wanted. To complete the proof it suffices to show that we can embed $`V//H`$ as an open subset into a divisorial toric variety $`X^{}`$ such that $`g`$ viewed as a morphism from $`X`$ to $`X^{}`$ is toric.
For this last step we argue as follows: Note that we constructed $`V`$ as an open $`H`$-invariant subset of the toric variety $`W`$. In , J. Świȩcicka shows that “maximal” open subsets with a good quotient by a given subtorus in a toric variety are in fact toric subvarieties.
More precisely, according to \[21, Corollary 2.4\], $`V`$ is contained in an open toric subvariety $`V^{}W`$ with a good toric quotient $`p^{}:V^{}V^{}//H`$ such that the induced map $`V//HV^{}//H`$ is an open inclusion. Of course, we can choose $`V^{}`$ in such a manner that $`V^{}//H=T^{}(V//H)`$ holds, where $`T^{}`$ denotes the big torus of $`V^{}//H`$. We set $`X^{}:=V^{}//H`$ and $`U:=V//H`$ and arrive at the following commutative diagram:
The morphism $`XV^{}`$ sending $`x`$ to $`\widehat{g}(x)`$ is a dominant toric morphism because $`\widehat{g}:XW`$ is one. Hence the same is true for $`g=p^{}\widehat{g}:XX^{}`$. Moreover, because $`\widehat{g}(X)V`$ holds, we conclude that the big torus $`T^{}`$ of $`X^{}`$ is contained in $`U`$. It follows that the complement $`X^{}U`$ is of codimension at least 2 in $`X^{}`$. Thus Lemma 1.7 yields that the toric variety $`X^{}`$ is also divisorial. This ends the proof. ∎
## 6. Divisorial reduction and categorical quotients
In this section we come to the main results of this article. Recall from that a categorical quotient for a $`G`$-variety $`X`$ is a $`G`$-invariant regular map $`XY`$ such that any $`G`$-invariant regular map $`XZ`$ factors uniquely through $`XY`$. Clearly, this notion can be restricted to any subcategory of the category of algebraic varieties, as soon as the $`G`$-variety $`X`$ belongs to this subcategory.
We give an answer to the problem of existence of categorical quotients for subtorus actions in the divisorial category. Our method of proof in fact solves the existence problem of a more general universal object: Consider a toric variety $`X`$ with big torus $`T`$ and the action of a subtorus $`HT`$.
###### Definition 6.1.
An $`H`$-invariant divisorial reduction of $`X`$ is a regular map $`r:XY`$ to a divisorial variety $`Y`$ such that every $`H`$-invariant regular map $`f:XZ`$ to a divisorial variety $`Z`$ admits a unique factorization $`f=\stackrel{~}{f}r`$ with a regular map $`\stackrel{~}{f}:YZ`$. If $`H=1`$, then we simply speak of a divisorial reduction.
A candidate for such a reduction is constructed in two steps. First, recall from that there is a toric quotient for the action of $`H`$ on $`X`$, that means a toric morphism
$$p:XX/\mathrm{tq}H$$
which is a categorical quotient for the action of $`H`$ on $`X`$ in the category of toric varieties. In a second step, construct the toric divisorial reduction of the toric quotient space as described in Section 3:
$$q:X/\mathrm{tq}H(X/\mathrm{tq}H)^{\mathrm{tdr}}.$$
###### Theorem 6.2.
For a toric variety $`X`$, the following statements are equivalent:
1. $`X`$ admits an $`H`$-invariant divisorial reduction.
2. The composition $`qp:XZ`$ is surjective.
Moreover, if one of these statements holds, then $`qp`$ is the $`H`$-invariant divisorial reduction.
Applying this result to divisorial toric varieties $`X`$, we obtain the following solution for the above quotient problem:
###### Corollary 6.3.
The action of a subtorus $`H`$ on a divisorial toric variety $`X`$ admits a categorical quotient in the category of divisorial varieties if and only if the composition of $`XX/\mathrm{tq}H`$ and $`X/\mathrm{tq}H(X/\mathrm{tq}H)^{\mathrm{tdr}}`$ is a surjective map.
A further special case of Theorem 6.2 is the case of a trivial torus $`H=1`$. Here we obtain the following:
###### Corollary 6.4.
A toric variety admits a divisorial reduction if and only if its toric divisorial reduction is surjective.
###### Proof of Theorem 6.2.
Assume first that $`qp`$ is surjective. We show that a given $`H`$-invariant regular map $`f:XZ`$ to a divisorial variety $`Z`$ factors through $`qp`$. Lemma 5.1 yields a decomposition $`f=hg`$ with an $`H`$-invariant dominant toric morphism $`g:XX^{}`$ to a divisorial toric variety $`X^{}`$.
By the universal properties of $`p`$ and $`q`$, the toric morphism $`g`$ has a factorization $`g=g^{}(qp)`$. By surjectivity of $`qp`$, the map $`h`$ is defined on a neighbourhood of the image of $`g^{}`$. Hence $`f=(hg^{})(qp)`$ is the desired factorization. Thus $`qp`$ is the $`H`$-invariant divisorial reduction of $`X`$.
Conversely, suppose that $`X`$ has an $`H`$-invariant divisorial reduction $`r:XY`$. Since the normalization of a divisorial variety is again divisorial, we can conclude that $`Y`$ is normal. Moreover, the universal property of $`r:XY`$ implies that $`r`$ is surjective, and that $`Y`$ inherits a set-theoretical action of the big torus $`TX`$ making $`r`$ equivariant. Note that a priori it is not clear that this action is regular, so we cannot treat $`Y`$ as a toric variety.
Let $`Z:=(X/\mathrm{tq}H)^{\mathrm{tdr}}`$. We shall compare the $`H`$-invariant divisorial reduction $`r:XY`$ with the toric morphism $`qp:XZ`$. On the one hand, because of the universal property of $`r`$, the map $`qp`$ factors uniquely through $`r`$. So there is a unique regular map $`\alpha :YZ`$ with $`qp=\alpha r`$.
On the other hand, Lemma 5.1 provides a decomposition $`r=hg`$ with a dominant toric morphism $`g:XX^{}`$ to a divisorial toric variety $`X^{}`$ and a rational map $`h`$ from $`X^{}`$ to $`Y`$ that is defined on the image of $`g`$. By the universal properties of $`p`$ and $`q`$, we have $`g=g^{}qp`$ with a toric morphism $`g^{}:ZX^{}`$. So we arrive at the following commutative diagram:
Note that $`g^{}(q(p(X)))=g(X)`$ is contained in the domain of definition of the rational map $`h`$. Since $`r`$ is surjective, we have $`q(p(X))=\alpha (Y)`$ and we obtain that $`h`$ is defined on $`g^{}(\alpha (Y))`$. It follows that $`(hg^{})\alpha `$ is the identity on $`Y`$. This shows that $`\alpha `$ is injective. Moreover, on the big torus of $`Z`$, the map $`\alpha (hg^{})`$ is the identity.
Consequently $`\alpha :YZ`$ is a birational injection. Since $`Z`$ is normal, Zariski’s main theorem tells us that $`\alpha `$ is in fact an open embedding. Since the image $`\alpha (Y)`$ is invariant under the induced set-theoretical action of $`T`$ on $`Y`$, the map $`\alpha `$ is an isomorphism. In particular, $`r:XY`$ is surjective. ∎
We conclude this section with some examples. The above results in many situations give positive answers to the problem of existence of quotients. A typical case are toric varieties defined by fans with convex support:
###### Corollary 6.5.
Let $`X`$ be a toric variety arising from a fan with convex support. Then $`X`$ admits a divisorial reduction.
###### Proof.
Let the toric divisorial reduction $`q:XX^{}`$ arise from a map $`Q:NN^{}`$ of fans $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }^{}`$. Then $`\sigma :=Q(|\mathrm{\Delta }|)`$ is a convex cone in $`N^{}`$ and $`\sigma |\mathrm{\Delta }^{}|`$. Intersecting the cones of $`\mathrm{\Delta }^{}`$ with $`\sigma `$, we obtain a further fan in $`N^{}`$, namely
$$\mathrm{\Delta }^{\prime \prime }:=\underset{\tau ^{}\mathrm{\Delta }^{}}{}𝔉(\tau ^{}\sigma ).$$
Let $`X^{\prime \prime }`$ be the associated toric variety. The identity map $`NN^{}`$ defines an affine toric morphism $`g:X^{\prime \prime }X^{}`$. In particular, $`X^{\prime \prime }`$ is divisorial. Moreover, $`Q:NN^{}`$ is also a map of the fans $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }^{\prime \prime }`$. The corresponding toric morphism $`q^{}:XX^{\prime \prime }`$ is surjective because $`Q(|\mathrm{\Delta }|)`$ equals $`|\mathrm{\Delta }^{\prime \prime }|`$. Consider the decomposition
$$\text{}.$$
The universal property of the toric divisorial reduction implies that $`g:X^{\prime \prime }X^{}`$ is an isomorphism. Hence $`q:XX^{}`$ is surjective and the assertion follows from Corollary 6.4. ∎
###### Corollary 6.6.
Let $`X`$ be a divisorial toric variety arising from a fan with convex support. Then every subtorus action on $`X`$ admits a categorical quotient in the category of divisorial varieties.
###### Proof.
Let the toric quotient $`p:XX^{}`$ arise from a map $`P:NN^{}`$ of fans $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }^{}`$. By \[2, Remark 2.5\], each cone $`\sigma ^{}\mathrm{\Delta }^{}`$ is generated by images $`P(\sigma )`$ of certain $`\sigma \mathrm{\Delta }`$. Thus also $`\mathrm{\Delta }^{}`$ has convex support and $`p:XX^{}`$ is surjective. So, Corollaries 6.3 and 6.5 give the claim. ∎
However, Corollary 6.3 also provides counterexamples to existence of quotients. There can be different reasons for nonsurjectivity of $`qp`$, as the following examples show:
###### Example 6.7.
For the toric variety $`X`$ described in Example 3.6 the toric divisorial reduction is not surjective. Hence $`X`$ does not admit a divisorial reduction. Moreover by Cox’s construction, see , $`X`$ is a good quotient of an open subset $`\widehat{X}𝕂^9`$ by a five dimensional subtorus $`H(𝕂^{})^9`$. So, the action of $`H`$ on $`\widehat{X}`$ admits no categorical quotient in the category of divisorial varieties.
###### Example 6.8.
Let $`\mathrm{\Delta }`$ be the fan in $`^4`$ having the following maximal cones:
$$\begin{array}{ccc}\sigma _1:=\mathrm{cone}((1,0,0,0),(0,1,0,0)),\hfill & & \sigma _2:=\mathrm{cone}((0,0,1,0),(0,0,0,1))\hfill \end{array}$$
The associated toric variety $`X`$ is an open toric subset of $`𝕂^4`$. Define a projection $`P:^4^3`$ by
$$\begin{array}{ccc}P((1,0,0,0)):=(1,0,0),\hfill & & P((0,1,0,0)):=(0,1,0),\hfill \\ P((0,0,1,0)):=(0,0,1),\hfill & & P((0,0,0,1)):=(1,1,0).\hfill \end{array}$$
By , the toric morphism $`p:X𝕂^3`$ defined by $`P`$ is the toric quotient for the action of the subtorus $`H:=\mathrm{ker}(p)`$ on $`X`$. Since $`p`$ is not surjective, the action of $`H`$ on $`X`$ has no categorical quotient in the category of divisorial varieties.
## 7. An open problem
In this article we have solved the problem of existence of categorical quotients for subtorus actions on toric varieties in the divisorial category. For the analogous question in the category of all algebraic varieties we have partial results.
For example, the toric quotient $`p:XX/\mathrm{tq}H`$ is a categorical quotient in the category of algebraic varieties if the subtorus $`H`$ is of codimension at most two , or if the map $`p`$ satisfies a certain curve lifting property and $`X/\mathrm{tq}H`$ is of expected dimension .
However, the general question still remains open. Therefore we pose it here as a problem:
###### Problem 7.1.
Give necessary and sufficient conditions for subtorus actions on toric varieties to admit a categorical quotient in the category of algebraic varieties. |
warning/0001/astro-ph0001517.html | ar5iv | text | # A Young Globular Cluster in the Galaxy NGC 69461footnote 11footnote 1Based on observations made with the Nordic Optical Telescope, operated on the island of La Palma jointly by Denmark, Finland, Iceland, Norway, and Sweden, in the Spanish Observatorio del Roque de los Muchachos of the Instituto de Astrofisica de Canarias.
## 1 Introduction
During a recent search for massive young clusters in the nearby spiral galaxy NGC 6946, Larsen & Richtler (1999) found a circular bubble containing numerous small clusters and a bright compact cluster that resembled a young globular. Some of the small clusters in the bubble are organized into arc-like shapes. This feature was found and sketched first by Hodge (1967) as the only positive detection in his search for multiple-arc structures similar to that in the Constellation III region of the LMC. The diffuse object inside the bubble was noted later by Efremov (1999), who suggested it was a cluster causally related to the arcs in the same way that the LMC massive cluster NGC 1978 might be related to Constellation III.
The mass and radius of the brightest cluster in this region make it similar to the globular clusters in the outer halo of the Milky Way (see review in VandenBerg, Stetson & Bolte 1996). Globular cluster formation is very rare in normal galaxy disks (the Milky Way disk has none), yet other young globulars have been found in dwarf galaxies like the LMC (van den Bergh 1991; Richtler 1993), and in interacting galaxies (Whitmore & Schweizer 1995) and starbursts (Maoz et al. 1996; Holtzman et al. 1992; Meurer et al. 1995). Because of this, there is growing anticipation that we will soon be able to understand the origin of the old population of globular clusters in the Milky Way and other galaxies by studying their younger versions. In this respect, we can use the young globular in NGC 6946 for closer study.
Here we consider the properties of the bright compact cluster in NGC 6946 in view of the environment in which it formed. The observations are presented in section 2, an analysis of the energetics and pressures of the region are in section 3, and a discussion of possible mechanisms for the formation of the globular are in section 4.
## 2 Observations
Observations of clusters in NGC 6946 and other galaxies were made with the Nordic Optical Telescope by Larsen & Richtler (1999), who discussed the results in general terms. The bubble surrounding the brightest cluster in the galaxy caught our attention because of its peculiar structure and possible implications for understanding globular cluster formation.
Figure 1 shows a color<sup>2</sup><sup>2</sup>2Black and White for Astro-ph to limit file size image of this galaxy in B, V, and R+H$`\alpha `$, with an insert magnifying the bubble region. The structure of the star formation here is unlike anything else in the galaxy. It is located at the end of a short spiral arm, downstream from a dust lane, with a circular bubble outlined by an arc of clusters in the west and a dark region, suggesting dust, just to the right. Red patchy emission from H$`\alpha `$ surrounds the bubble as if star formation is continuing there in dense residual gas. The brightest and largest concentration of starlight emission inside the bubble is the globular-like cluster. The overall dimensions of the cavity are $`600`$ pc east-west and $`550`$ pc north-south, assuming a distance of 5.5 Mpc (Tully 1988).
The extinction in the region is shown better by a plot of V$``$I in figure 2, which covers 28”$`\times `$33” around the bubble. Many of the dark patches are irregularly shaped, like clouds in the midst of diffuse and clustered stellar emission. A large dark region near the center of the bubble, north of the brightest cluster, could be the remains of the molecular cloud from which the cluster formed. Its dimensions are $`275\times 80`$ pc. We measured the extinction in this region in two ways: the brightness deficiency in the B band, and the V-I color excess, both compared to the surrounding light. These two measurements are independent and they give the same result, $`A_B1`$ mag.. This implies that the feature is really dust and not just a lack of stars, and that the dust is mostly in front of the emission. With this average extinction for the molecular cloud feature and its size given above, the mass becomes $`2\times 10^5`$ M.
Another V$``$I dust feature halfway out to the western edge resembles a cometary globule, and this, along with the scalloped structures along the western edge, point back to the brightest cluster. Other dust patches and partial arcs are also evident in the V-I figure, inside the bubble and beyond the bubble boundary. The dark region just to the north of the globular cluster is probably an artifact of color variations in the point spread profiles of the bright cluster image.
Figure 3 shows a V-band image of the same bubble region on the left and a de-reddened V-band image on the right, made from the observed ratios of V and I band intensities and corrected in V for the corresponding extinction, assuming all the dust is foreground and the background color is uniform. Many of the dark regions and bands in figure 2 fill in with what appear to be tiny clusters, and some of the arc-like structures go away. This technique cannot correct for extremely heavy obscuration because then there is no emission for the determination of V-I color excess. The dark region beyond the bubble in the west, seen in figure 2, may have this problem (the spots indicate that the V/I intensity ratio has overflowed).
A low-contrast I-band image with cluster ages in millions of years is shown in figure 4. The ages were determined from their integral colors, U-B and B-V, following the S-method and calibration in Girardi et al. (1995). For the brightest clusters, the ages range from 10 to 30 Myr with no overall pattern inside the bubble. The colors of the globular cluster (U-B=-0.74, B-V=0.01, V-I=0.52; Larsen 1999), give it an age of $`15`$ million years. The bright spot in the southeast outside the bubble is a foreground star.
A color-magnitude diagram for the same objects is shown in figure 5. The lines are stellar isochrones for stars at the indicated ages, with each line spanning a wide range of stellar masses from 0.6 M, which is outside the figure, to the masses corresponding to the endpoints of stellar evolution (from Bertelli et al. 1994). The diamonds are the objects inside the bubble; the source at $`V=12.9`$ is the globular cluster. This figure demonstrates that the small sources inside the bubble are too bright to be stars if they are at the distance of NGC 6946. Even the brightest stars in M33, at M$`{}_{V}{}^{}9`$ mag (Humphreys & Sandage 1980), are not as bright as the objects in the bubble.
The ages given in figure 4 are slightly sensitive to unknown absorption. We can estimate the error from figure 10 in Girardi et al. (1995), which suggests that a color excess of E(B-V) = 0.2 mag. changes an age of 15 My to 11 My. Thus the ages determined from the S-method are probably accurate to no better than $`\pm 30`$%, and may be systematically large. We can also estimate the ages from the extinction-free $`Q`$ parameter ($`Q=(UB)0.72(BV)`$; van den Bergh 1968), calibrating Q vs. log(age) using the population synthesis models by Bruzual & Charlot (1996), and assuming all the emission is behind the dust and not mixed with it. These ages also tend to be slightly less than the values given in figure 4. For example, the globular cluster age becomes $`14`$ My from its extinction-free $`Q`$ value. This second check on the ages is also useful because the Girardi et al. ages were calibrated for the LMC, where the metallicity is relatively low. Considering measurement errors and extinction uncertainties, the ages in figure 4 are probably accurate only to $`\pm 10`$ My. In what follows, we assume an age for the whole region of $`30`$ My; it could range between 20 and 40 My, but this distinction is not important. The energy required to make the bubble scales as the inverse square of the age.
The absolute magnitude ($`V=12.9`$ mag) of the globular cluster makes its mass $`6.9\times 10^5`$ M, from the calibration in Bruzual & Charlot (1993, 1996; see Leitherer et al. 1996), assuming a Salpeter IMF between 0.1 M and 150 M. If we assume a more realistic flattening of the IMF below $`0.3`$ M (Festin 1997; Hillenbrand 1997; Luhman & Rieke 1998; Lada, Lada & Muench 1998; Hillenbrand & Carpenter 1999), then the mass is $`5\times 10^5`$ M. Other clusters in the same region are at least 2 magnitudes fainter than the globular, although they have about the same ages, and all the other clusters elsewhere in this galaxy are at least 1 magnitude fainter than the globular.
The radial intensity profile for the globular cluster is shown in figure 6, along with the profile for the bright star in the southwest. The globular is clearly more extended than the star. The half-light radius is 11 pc from a fit to a King light profile, corrected for telescope resolution (Larsen & Richtler 1999; Larsen 1999). This radius is highly uncertain because the seeing resolution of 0.6” corresponds to 15 pc. The radial profiles of the fainter objects in the field are somewhat extended too, making them appear fuzzy in the image and more like clusters than stars. However, the S/N for these profiles is too poor to give cluster sizes. We intend to verify the cluster nature of the fainter objects using HST images.
The large-scale environment of the bubble is shown in figure 7, from an I-band image. The bubble is clearly defined in the midst of spiral arm stars and surrounding dust. There is a thin dust arc just outside the bubble edge in the southwest, which is probably part of the dust shell surrounding the cavity. A larger dust feature is at the edge of the bubble in the northwest, connecting to the spiral arm dust lane that trails off further in this direction (cf. Fig 1). Essentially all of the star formation in the vicinity of the globular cluster is confined to the bubble cavity; the surrounding dust hides only a small amount of additional star formation in comparison to what has already occurred, considering the weak H$`\alpha `$ outside the bubble (Fig. 1). This suggests that much of the star formation inside the bubble, including the arc of clusters at the edge, is causally related and possibly triggered or synchronized by a series of high pressure events.
## 3 The Origin of the Bubble
### 3.1 Multiple stellar winds and supernovae
Figures 27 indicate that the bubble in NGC 6946 contains a single bright cluster and a concentration of $`50`$ smaller sources, many of which are probably smaller clusters.
Most of these clusters are concentrated within a $`300`$ pc diameter region inside a $`600`$ pc cavity. A few other clusters lie along the rim of the bubble and more may be outside, producing the H$`\alpha `$ emission evident in figure 1. The densest concentration of sources is within $`60`$ pc of the globular, as seen in figure 3. We count about 25 objects there. This concentration is packed so tightly that many of the luminous sources blend together.
The total luminosity of all of the bright clusters in the bubble suggests that a large number of supernovae exploded during its $`30`$ My lifetime. These supernova and the winds of the associated massive stars could have dispersed most of the dense gas that formed the clusters and left only the $`10^5`$ M dust cloud and other dust clouds seen in figure 2, now located $`100`$ pc or more from the globular. The cometary and shell-like debris inside the cavity could have been shaped by the supernova and stellar winds too.
A disturbance $`R300`$ pc in radius that has been driven for $`t30`$ million years by continuous pressures from winds and supernovae will have a mass $`M\pi R^2\mathrm{\Sigma }`$ for ambient surface density $`\mathrm{\Sigma }`$ and a speed $`v=0.6R/t`$ (Weaver et al. 1977). The kinetic energy is $`E=0.5Mv^20.5\pi R^2\mathrm{\Sigma }\left(0.6R/t\right)^2=2.1\times 10^{51}A_V`$ erg. Here we have rewritten the ambient mass column density of the galaxy in terms of the equivalent V-band extinction through the disk, $`\mathrm{\Sigma }=4.5\times 10^3`$A<sub>V</sub> gm cm<sup>-2</sup>, using the conversion in Bohlin, Savage & Drake (1978) with 10% He by number.
Typically the energy required to make a bubble from multiple supernovae and stellar winds is $`10`$ times the kinetic energy of the expansion because of energy loss from radiation in the shocks and decompression of the cavity (MacLow & McCray 1988). For example, the product of the wind luminosity and the time in the Weaver et al. (1977) model (cf. their equation 21) is 5.1 times the instantaneous kinetic energy. For a factor of 10 from MacLow & McCray, the source energy becomes $`10^{52.3}A_V`$ ergs, scaling with the inverse square of the age of the bubble.
We expect $`A_V3`$ before star-formation began for the average of the dense part of this spiral arm gas, based on direct measurements of extinction in spiral arm shocks (Elmegreen 1980) and on extinction, HI, and CO maps of this galaxy. The average extinction on $`100`$ pc scales surrounding the bubble currently exceeds 2 magnitudes in B band (Trewhella 1998), while the summed atomic (Boulanger & Viallefond 1992) and molecular (Tacconi & Young 1989) gas column densities over a 1 kpc region correspond to $`A_V=1.1+1.4=2.5`$ mag.
If $`A_V`$ equaled 3 mag in the vicinity of the bubble before star-formation began, then the expansion required the equivalent of $`60`$ supernovae, each with a typical energy of $`10^{51}`$ ergs. If the bubble age is only 15 My, the age of the globular cluster, then the equivalent number of supernova would be $`240`$. This is easy for a cluster with $`5\times 10^5`$ M, because there would be 2700 stars more massive than the turnoff mass at 15 My (which is 14 M), assuming a flattened Salpeter IMF out to 150 $`M_{}`$ and stellar evolution models in Bertelli et al. (1994). With so many expected supernovae, the expansion age could be even younger than 15 My, or $`A_V`$ larger than 3. Presumably some of these supernovae made the smaller shells inside the bubble too, producing the overall frothy appearance in figure 1. In this case these supernovae may have been ejected from the globular.
The average energy density in the bubble is $`10^{52.3}A_V/\left(4\pi R^3/3\right)10^{4.6}A_V`$ cm<sup>-3</sup> K. The energy density in the core, at the peak of the input rate, should have been at least 10 times this average. For $`A_V3`$, this peak is $`1.2\times 10^6`$ cm<sup>-3</sup> K, which is similar to that around OB associations in the Solar neighborhood. Thus the pressure that made the bubble was not unreasonably large, in spite of the high density of stars in the center. This is because the overall star formation rate in the $`600`$ pc region inside the bubble is typical for spiral arms. The only difference in the bubble here is that most of this star formation went into a single dense cluster, rather than a giant, extended star complex, like Gould’s Belt.
### 3.2 Alternative Scenarios
An alternative source of pressure for the bubble should also be considered for this massive central cluster. It seems possible that such a cluster would have sufficient opportunity to produce several rare hypernovae (Paczyński 1998) from the collapse of extremely massive or otherwise rare single or binary stars. For example, a Salpeter IMF that flattens below 0.3 M in a cluster that contains $`5\times 10^5`$ $`M_{}`$ would be likely to produce $`8`$ stars with masses between 100 M and 150 M if the maximum mass were 150 M. Such massive stars could produce hypernovae of the type described by Paczyński (1998). If the IMF continued up to arbitrarily high masses, then the largest star would have an unlikely mass of $`900`$ M in a $`5\times 10^5`$ M cluster, just from sampling statistics alone (without any consideration of whether such massive stars are physically possible). Of course, stars this massive have never been found, but neither has a fixed upper limit to the stellar mass in a cluster: the largest star just keeps getting larger for larger clusters (e.g., Massey & Hunter 1998). Thus the largest star could conceivably contain several hundred M. If such a star, or any other rare object exploded, possibly making a gamma ray burst, it could have an energy larger than a single supernovae by a factor of $`100`$ (Rees & Mézáros 1998) and be able to produce one of the large arcs of gas and dust individually (see also Efremov 2000).
Although there is no direct evidence for unusual pressure sources like hypernovae, there are several remarkable properties of the bubble that are not usually part of the standard, multiple supernovae/stellar wind scenario (Tenorio-Tagle & Bodenheimer 1988).
First, the bubble is not elliptical or elongated along the major axis of the galaxy. This implies that it is not a circular ring in the midplane (unless an intrinsic ellipticity fortuitously cancels the projection effect), but more like a sphere with some thickness perpendicular to the plane. Yet its diameter (600 pc) is much larger than the likely disk thickness. This leads us to wonder how such a circular bubble could have formed in the first place. Its uniform shape would seem to require a very uniform external medium for the pressure to push around, yet disturbances perpendicular to the plane generally propagate into decreasing densities. Perhaps the shocked gas in the density wave really does have such a large thickness, and what we were calling a GMC in Section 2 is really the near side of the bubble, far above the plane.
Second, the clusters that lie along the bubble rim to the north, west, and south are somewhat uniformly spaced, as if they were triggered by the regular gravitational collapse of a swept-up shell (e.g. McCray & Kafatos 1987). However, the ages of these rim clusters are about the same as the ages of all the other clusters inside the bubble, including the globular cluster, so there is no evident time sequence or source of pressure for this triggering. Efremov (2000) suggested that there had to be some invisible older generation ($`50`$ My) of clusters or super-explosions that made all of these features at about the same time, i.e., the overall bubble, the outer cluster arcs, and possibly the globular cluster. Alternatively, the photometric ages could be wrong.
Third, the interior of the bubble contains dust (and presumably gas) clouds, including a comet-shaped cloud half-way out to the western edge. How can the bubble outside these clouds be so circular? Wouldn’t the comet and other interior clouds block the expansion along their paths if the pressure source was near the bubble center? Does this imply that the bubble formed before the comet cloud and that there were two epochs of high pressure in this region? A similar situation seems to occur in the Constellation III region of the LMC, where the HI cavity is much larger than a triggered arc of stars (Efremov & Elmegreen 1998). The solution here might be simply that the dark debris in the interior of the bubble is far above the plane on the near side, and not on the line of sight between the central pressure source and the western edge.
Fourth, there is apparently some star formation beyond the bubble, in the H$`\alpha `$-emitting regions. This makes us wonder if the western arc of clusters is not really an arc of stars but an arc of clear-viewing, seen through a dust screen made up of dense material with its own arc shape, just inside the stellar arc, and other dense material with an arc-shaped inner border, just beyond the western edge (cf. Fig. 2). What is the source of the inner gas shell in this interpretation? Can multiple supernovae and stellar winds make such a shell inside a pre-existing bubble cavity, or does it require a more concentrated giant explosion as in the hypernova scenario discussed above? Did this additional explosion also occur after the bubble formed, as suggested above for the comet cloud?
All in all, the high concentration of star formation in this region, forming what appears to be a dense cluster of clusters, is unusual for non-interacting galaxies. The presence of background stars to the west of the bubble in figure 7 suggests that the obscuration there is not so large that it could hide a significant amount of additional star formation, at least not comparable to what is in the bubble itself. This means that the western edge of the bubble is only partly the result of an end to the clearing of the cavity; it is also from a lack of bright clusters beyond. If there was any triggering of star formation, especially to make the arcs along the bubble edge, then this triggering has stopped near the current limit of the bubble. The sharp western edge of the bubble may then be explained if this is a triggered shell of clusters.
The sharp edge and the perfect circular shape of the western rim of the bubble is the property shared also by a number of other stellar arcs in a few galaxies, first of all in the Constellation III region of the LMC. The circular and sharp outer edges indicate triggered star formation in the partial spherical shells, the occurrence of the multiple arcs near each other being connected with the common origin of the pressure sources. Most probably they were ejected from a rich nearby cluster (Efremov, 2000). As mentioned above, the exact nature of suggested triggering is difficult to understand, as is the source of pressure, but super-explosions cannot be overlooked.
Thus we see that even though the energy and pressure inside the bubble are somewhat normal for an aggregate of stars like this, the morphology of the whole region is baffling. Perhaps higher resolution and longer wavelength observations will help clarify these questions.
## 4 Speculations about Globular Cluster Formation
### 4.1 General Discussion of the Observed Environment for Globular Cluster Formation
Why should a globular cluster form in this region of NGC 6946 and no where else? This is a special location at the abrupt end or break-point in a secondary spiral arm. When star formation began, there was a gross asymmetry in the self-gravitational force of the spiral arm gas, with most of the acceleration pointing inwards along the arm. Such asymmetry would cause catastrophic collapse in a large mass of gas, possibly like the end of a tidal tail (Barnes & Hernquist 1992; Elmegreen et al. 1993). The total atomic (Boulanger & Viallefond, 1992) and molecular (Tacconi & Young 1989) mass observed in a $``$kpc region around the globular cluster is $`10^8`$ M, which is $`10`$ times higher than the typical cloud mass from normal density wave triggering (Elmegreen 1994).
Not all disk globular clusters form at the ends of spiral arms, however, and the other arm endpoints in NGC 6946 contain no globular clusters, even though they often contain massive star-forming regions (cf. Fig. 1). Something peculiar happened here to make the dense, massive cloud that was necessary for the globule to have formed.
What is unusual about this region? The density of the globular cluster alone is not particularly high. The average half-light density is only $`200`$ stars pc<sup>-3</sup>, which is much smaller than the $`10^4`$ stars pc<sup>-3</sup> in the Orion Trapezium cluster (Prosser et al. 1994; McCaughrean & Stauffer 1994) and Mon R2 (Carpenter et al. 1997). The corresponding $`H_2`$ density would have been $`6\times 10^3`$ cm<sup>-3</sup> with a 25% star-formation efficiency (e.g., Elmegreen et al. 2000), and this is similar to that in a GMC core. The main peculiarity with the globular cluster is the large mass at this density. For an efficiency of 25%, the gas mass to make a $`5\times 10^5`$ M cluster would have been $`2\times 10^6`$ M, which is comparable to that of the largest GMCs in the Milky Way but here concentrated inside several tens of parsecs. Such a high mass concentration means that the virial velocity and pressure of the cluster-forming core were large. A $`M=2\times 10^6`$ M cloud with a half-light radius of $`R11`$ pc has a virial speed $`\left(0.2GM/R\right)^{1/2}12`$ km s<sup>-1</sup> and a pressure $`P0.1GM^2/R^4`$, which converts to $`P/k_B=6\times 10^8`$ K cm<sup>-3</sup> for Boltzmann constant $`k_B`$.
The pressure would have been slightly lower in the surrounding gas. If we consider the core of luminous sources surrounding the globular cluster in figure 3, which has a $`60`$ pc radius, and estimate that the total mass in that core increased linearly with radius as for an $`r^2`$ isothermal density profile, then the mass there would be $`5`$ times larger than the globular core mass, or $`10^7`$ M. The average pressure would be $`2\times 10^7`$ K cm<sup>-3</sup> and the virial speed still $`12`$ km s<sup>-1</sup>.
These virial speeds are comparable to the thermal speed of an HII region, suggesting that the O-type stars which formed in the clusters would have had some difficulty destroying the cloud cores. The core pressures are also large enough to avoid complete ionization of the gas. These are welcome checks on the most basic theory of globular cluster formation: the cloud-binding energy densities and escape speeds have to be large enough to keep the gas around to continue forming stars, in spite of all the O star ionization, until the critical efficiency to form a bound cluster is reached (Elmegreen & Efremov 1997).
The pressure in the 60 pc core region was $`10`$ times higher than the typical pressure in a local OB association and molecular cloud core and $`10^3`$ times higher than the ambient pressure in the Solar neighborhood. The pressure in the 11 pc core was $`100`$ times higher than in a local OB association/GMC core. Most likely the ambient pressure in the NGC 6946 region was $`10`$ times the ambient pressure in the Solar neighborhood, and the extra factor of $`1001000`$ for the cloud cores over the ambient value was from the self-gravity of the material surrounding it. Such overpressures are typical for star-forming clouds in the Solar neighborhood. What might have given the pre-cluster interstellar medium a pressure of $`2\times 10^5`$ cm<sup>-3</sup>K over a region perhaps a kiloparsec in size?
Considering the location of this region of star formation inside NGC 6946, we should probably concentrate on the effects of a spiral density wave shock. For an incoming average HI density of 1 cm<sup>-3</sup>, a shock speed of $`35`$ km s<sup>-1</sup> would have reached the required $`2\times 10^5`$ K cm<sup>-3</sup> pressure over a large region. Such a shock speed is not unusual for a galaxy like this, and when combined with an asymmetric collapse from the arm end, could have conspired to make the pressure in a massive virialized cloud somewhat high.
A second constraint on the formation of the globular cluster is its small photometric age. Obviously it had to form very quickly, probably within only several million years, to have such a young age with no obvious dust or irregular structure remaining inside of it. The only dense dusty material that we see nearby is a $`10^5`$ M cloud some 100 pc away, which could be the remains of the cloud core that formed it. Such a quick formation time follows from the most basic model (e.g., Fig. 4 in Elmegreen & Efremov 1997) if clouds typically convert their gas into stars in only several dynamical time scales of the cloud core. At the pre-cluster density in the globular core, the dynamical time was $`\left(G\rho \right)^{1/2}0.8`$ million years. This is fast enough to form the globular and clear away most of the debris within its $`15`$ My lifetime. Such a short time scale also implies that supernovae in the cluster did not enrich the gas much before it was cleared away. This helps us understand how the oldest globulars in the Milky Way halo can have very low metallicities even though they must have had massive stars.
### 4.2 A Coalescence Model for the Formation of the Globular Cluster
These aspects of the environment of the globular cluster in NGC 6946 are all consistent with what we believe to be necessary for the formation of such a massive object with a normal IMF, and they are also consistent with the location of this cluster inside the galaxy. They do not indicate how the globular actually formed, however. It need not have formed by the monotonic collapse of a single cloud core, for example, as in the standard scenario for cluster formation. Instead, it could have formed by the agglomeration and accretion of smaller clusters in the vicinity, provided the timescale for such accretion was short enough. Indeed, there are several peculiarities of star formation here that suggest such accretion might have actually happened.
First, there is an extremely high concentration of small luminous sources that may be individual clusters surrounding the globular cluster. This is shown best by the extinction-corrected image in figure 3. We count $`25`$ small objects within a radius of 60 pc, and estimate that there are probably an equal number in the foreground and background of the bright globular cluster image. This would make the current cluster density $`n_{cl}6\times 10^5`$ pc<sup>-3</sup> if all these objects are small clusters.
Second, the mass function of all the known clusters in this region is discontinuous. The globular cluster is $`10`$ times more massive than any of the other clusters. This is unusual for disk cluster systems, which tend to have a continuous power-law mass distribution (van den Bergh & Lafontaine 1984; see the review in Elmegreen et al. 2000). Discontinuous mass functions, with a single largest object having a mass much larger than any other object, can arise from a runaway accretion process as the gravitational cross section of the largest object increases with mass. This could have been important for old globular clusters because their mass distribution differs from that of disk clusters today (van den Bergh 1995). In fact the old clusters have a characteristic mass similar to that of the globular cluster in NGC 6946, so formation processes like agglomeration, which break away from the power law distributions of normal disk clusters, might be important.
Third, the photometric age of the globular cluster, $`15`$ My, is about the average age of all the other clusters in the bubble, which is 14.4 My not counting the globular cluster itself, or 12.5 My if we omit the 58 My-old cluster in figure 4. This equality is to be expected in a coalescence model: the photometric age of the globular would not be its true age from some single star formation event, but the luminosity-weighted average age of all of the clusters it accreted. The true age could be the maximum age of the neighbors, because this marks the beginning of all star formation in the region and the time when the accretion process might have begun. This third point is uncertain, however, because of the statistical errors in the age measurements.
Finally, the globular cluster has a low average density, lower than the densities of embedded clusters that typically form near the Sun. Accretion can do this because it puffs up a cluster with the kinetic energy of the accreting pieces, giving it a lower average density than any single piece. We discuss this effect in more detail below.
We first check the accretion model by calculating the accretion rate for the small objects that are currently in the dense region surrounding the globular. This rate is $`n_{cl}v\sigma `$ for surrounding cluster density $`n_{cl}`$ estimated above, cluster velocity dispersion $`v`$, and capture cross section $`\sigma `$. The cluster velocity dispersion is taken to be 12 km s<sup>-1</sup> from the virial dispersion in the cluster core, derived above. The capture cross section is
$$\sigma \pi \left(R_{GC}+R_{cl}\right)^2\left(1+\frac{2GM_{GC}}{R_{GC}v^2}\right)$$
(1)
for core radii of the globular and incident clusters, $`R_{GC}11`$ pc and $`R_{cl}`$, globular cluster mass, $`M_{GC}5\times 10^5`$ M, and gravitational focusing factor ($`3.7`$) in large parentheses above. For $`R_{cl}<<R_{GC}`$, $`\sigma =1400`$ pc<sup>2</sup>. Then the accretion rate is $`n_{cl}v\sigma 1`$ field cluster per My. If the true age of the globular cluster is not the average photometric age of all the accreted clusters but the total age range in the whole region, which is twice this average value, then the globular had $`30`$ My to accrete smaller clusters from the dense swarm around it. Moreover, the initial field cluster density was probably higher before the accretion began, so the accretion rate could have been higher then too. Thus the globular could conceivably have grown by accretion over a $`30`$ My period from an initial size comparable to that of one of smaller clusters in the area.
The potential errors involved in this calculation should be reiterated because the observations are not good enough yet to be certain about the coalescence rate. First, the globular cluster radius is close to the seeing limit so the King-profile fit used to determine $`R_{GC}11`$ pc is inaccurate. The cluster could be smaller. Second, the small sources surrounding the globular in figure 3 could be image noise, or, if not, then some of the fainter ones could be luminous stars instead of clusters. Third, the ages of the clusters are not well known, so we cannot be certain that the photometric age of the globular is the average of all the others in the neighborhood. Fourth, the gas mass that formed the globular is not observed directly, but only inferred from the need for cluster self-binding, which implies a high star formation efficiency. The first three of these uncertainties could be clarified by higher resolution observations. There is no direct way to know what the previous cloud was like.
Nevertheless, the possibility that the globular cluster formed by the accretion of smaller clusters is intriguing, and, if true, could change our way of thinking about globular cluster formation. Thus we seek another piece of evidence from the low density of the globular cluster today. Such a density might be expected in a coalescence model because of the conversion of neighbor cluster orbital energy into globular cluster binding energy.
Consider two clusters with masses $`M_1`$ and $`M_2`$ and gravitational binding energies $`\alpha GM_1^2/R_1`$ and $`\alpha GM_2^2/R_2`$ for constant $`\alpha `$ depending on the mass distributions. If these clusters coalesce after an initial release at infinite distance and zero velocity, then, because of energy conservation, the final cluster will have a total energy $`0.5\alpha G\left(M_1^2/R_1+M_2^2/R_2\right)`$. After relaxation, this will equal the total energy of the new virialized cluster, $`0.5\alpha G\left(M_1+M_2\right)^2/R_3`$ for new radius $`R_3`$. For initial cluster densities $`n_1`$ and $`n_2`$, the final cluster density, $`n_3`$, becomes
$$n_3=\frac{\left(n_1^{1/3}M_1^{5/3}+n_2^{1/3}M_2^{5/3}\right)^3}{\left(M_1+M_2\right)^5}.$$
(2)
If we consider that the globular cluster formed by the successive accretion of $`N`$ identical smaller clusters, each of density $`n_0`$, and write the ratio of the current globular density to this initial density, $`X=n_3/n_0`$, then $`X`$ decreases iteratively with $`N`$ as
$$X_N=\frac{\left(1+X_{N1}^{1/3}N^{5/3}\right)^3}{\left(1+N\right)^5}.$$
(3)
Starting with $`X_{N=0}=1`$, this gives a sequence of $`X_N=1/(N+1)^2=1,`$ 0.25, 0.111, and so on. Thus after the successive addition of $`N=9`$ small clusters, a large cluster, 10 times as massive, forms with $`1`$% of the initial density. This factor of $`10`$ for mass is about the ratio of the globular cluster luminosity to the luminosities of the bright neighboring clusters (2.5 mag difference in brightness), and $`1`$% is about the ratio of the globular cluster stellar density to the density in a trapezium-type cluster, which is presumably similar to the cluster densities in the NGC 6946 neighborhood.
These simple checks suggest that the globular cluster in NGC 6946 could have had enough time, $`30`$ My, and currently has about the right mass and density, for it to have formed by the coalescence of many of the smaller objects, perhaps normal clusters, that seem to be nearby. The photometric age of the globular would then be the average age of these neighbor clusters, $`15`$ My, as observed within the errors. During this same total time of 30 My, the pressures from thousands of stellar winds and hundreds of supernovae in the dense star-forming region made the 600 pc bubble and pushed the residual dense gas to $`100`$ pc distance inside the bubble, leaving a $`2\times 10^5`$ M cloud and several cometary-shaped clouds.
## 5 Conclusions
A young globular cluster has been found inside a bubble of gas filled with numerous smaller clusters at the end of a short spiral arm in the nearby galaxy NGC 6946. The photometric age of the globular cluster is probably between 5 My and 30 My; the best fit is $`15`$ My. The mass of the globular is $`5\times 10^5`$ M, and the half-light radius was estimated to be $`11`$ pc by a fit to King model. This radius is less than the seeing limit (16 pc) so the actual radius could be smaller. The mass and size of the whole star-forming region correspond to such a high velocity dispersion and pressure in the previous cloud core that disruption by OB stars and stellar winds would have been difficult before the supernova era began. This observation satisfies the most fundamental requirement of bound cluster formation, that the star formation process be able to continue unimpeded until a high efficiency is reached. The total time for the cluster to form was probably very short, considering the $`1`$ My dynamical time in the dense part of the cloud core. This time is consistent with the young age and relaxed structure of the globular cluster today.
The bubble surrounding the globular cluster could have been made by the wind and supernova pressures from evolved massive stars in the cluster. No additional pressure sources are needed, like the “super-supernova” postulated by Hodge (1967). Nevertheless, such unusual explosions might be expected for such a massive cluster, and they may have contributed to some of the arc-like structures in the gas and dust. Some of these cluster arcs seem to be real, and then they could be triggered star formation. The average pressure in the bubble is typical for the pressures in OB associations near the Sun, and the bubble size is not unusual for giant disturbances in galaxies. The near-perfect circular shape of the bubble is unusual, though, considering its size in comparison to the likely scale height of the disk.
This circular shape and the absence of clusters outside the western parts of the rim of the ’bubble’ suggest another possible interpretation: that this is in fact not a cavity in absorption (a bubble), but a spherical shell of triggered star clusters, like some other arc-shaped regions of star formation in other galaxies (Efremov 2000). This would remove the difficulties with the clearing model. At any event, such a spherical feature is unusual, and more so if the similar ages for all of the clusters inside the shell are confirmed. Then Hodge’s (1967) suggestion that there was a super-explosion in this region will be more tenable. The only alternative might be an infall of a group of clouds.
Also unusual is the formation of a Gould’s Belt worth of star mass inside a half-light radius of only 11 pc or less. We suggested that high pressures and an asymmetric collapse from the associated spiral density wave might have produced the right conditions for such concentrated star formation, but why this region should differ from so many other sites of star formation at spiral arm ends is unknown. Perhaps the pressure source that formed the outer arc of clusters, if that came first, was also involved with the formation of the globular cluster.
The possible formation of the globular cluster by the coalescence of smaller clusters was also discussed. For the assumed parameters, the cluster accretion rate is high enough to explain the globular cluster mass, and the gradual decrease in virial density by the addition of moving sub-clusters can explain the current low globular cluster density. Also explained is the apparent globular cluster age, which seems to equal the average age of all the clusters in the region, and the unusual cluster mass function, which has one massive cluster at least 10 times bigger than the rest in the midst of a swarm of what appear to be more normal clusters. Whether other globular clusters or even open clusters have had similar accretion histories is an intriguing possibility. Better observations at higher angular resolution are necessary to check this model.
Acknowledgements:
Yu.E. appreciates partial support from the Russian Foundation for Basic Research and the Council for Support of Scientific Schools. The research of SSL was supported by the Danish Natural Science Research Council through its Centre for Ground-Based Observational Astronomy. Many helpful suggestions by Tom Richtler are gratefully acknowledged, as are useful comments by the referee. |
warning/0001/quant-ph0001100.html | ar5iv | text | # Algebra, Logic and Qubits: Quantum 𝛼┬´┴’𝛽𝛼𝜉.
## 1 Introduction
In his first article about quantum computers David Deutsch wrote in relation with Feynman’s work :
> Feynman ($`1982`$) went one step closer to a true quantum computer with his ‘universal quantum simulator’. This consists of a lattice of spin systems with nearest-neighbour interactions that are freely specifiable. Although it can surely simulate any system with a finite-dimensional state space (I do not understand why Feynman doubts that it can simulate fermion systems), it is not a computing machine in the sense of this article. ‘Programming’ the simulator consists of endowing it by fiat with the desired dynamical laws, and then placing it in a desired initial state. But the mechanism that allows one to select arbitrary dynamical laws is not modelled. The dynamics of a true ‘computer’ in my sense must be given once and for all, and programming it must consist entirely of preparing it in a suitable state (or mixed case).
The cited article (together with ) forms some standard for many works about quantum information science and here is necessary to say few words for explanation of some difference in methods and purposes of present paper.
First, it maybe even more interesting to understand simulation of boson systems, because due to infinite-dimensional state space (see Sec. 2) it is impossible to apply directly approach mentioned by David Deutsch. Is it possible to extend the ideas from systems with finite number of states to countable number (discrete spectrum)? Quantization of harmonic oscillator can be considered as simple example — quantum model of natural numbers, ‘quantum abacus ($`\stackrel{\mathrm{}}{\stackrel{´}{\alpha }}\beta \alpha \xi `$)’.
Second, the cited methods of description of quantum computations are similar with attempts to learn classical information science only by ‘wiring schemes’ of processors with huge number of gates. But usually it is efficient to consider general methods of mathematical logic, algebra etc. together or even before binary numbers and logical gates.
The algebraic methods currently are common for quantum theory and if we are interesting in boson and fermion systems, we should consider canonical commutation and anticommutation relations.
## 2 (Anti)commutation relations
The using annihilation and creation operators $`a`$, $`a^{}`$ for definition of quantum gates was introduced by Feynman . The operators meet anticommutation relation $`\{a,a^{}\}aa^{}+a^{}a=1`$ for one qubit. For system with $`n`$ qubits we may use CAR for $`n`$ fermions, i.e. introduce operators $`a_i`$, $`a_i^{}`$ for each qubit with properties:
$$\{a_i,a_j\}=\{a_i^{},a_j^{}\}=0,\{a_i,a_j^{}\}=\delta _{ij}$$
(1)
The algebra generated by $`2n`$ operators Eq. (1) is isomorphic with algebra of all complex $`2^n\times 2^n`$ matrices and so can be used for representation of any quantum gate .
In basis $`|0=\left(\begin{array}{cccc}\hfill 1& & & \\ \hfill 0& & & \end{array}\right)`$, $`|1=\left(\begin{array}{cccc}\hfill 0& & & \\ \hfill 1& & & \end{array}\right)`$ the operators can be expressed as:
$`a`$ $`=`$ $`{\displaystyle \frac{\sigma _x+i\sigma _y}{2}}=\left(\begin{array}{cc}\hfill 0& \hfill 1\\ \hfill 0& \hfill 0\end{array}\right)`$
$`a^{}`$ $`=`$ $`{\displaystyle \frac{\sigma _xi\sigma _y}{2}}=\left(\begin{array}{cc}\hfill 0& \hfill 0\\ \hfill 1& \hfill 0\end{array}\right)`$
$`a_i`$ $`=`$ $`\underset{ni1}{\underset{}{1\mathrm{}1}}a\underset{i}{\underset{}{\sigma _z\mathrm{}\sigma _z}}`$
$`a_i^{}`$ $`=`$ $`\underset{ni1}{\underset{}{1\mathrm{}1}}a^{}\underset{i}{\underset{}{\sigma _z\mathrm{}\sigma _z}}`$
The representation of CAR also can be expressed in more invariant and clear way with using of Clifford algebras. Let us recall it briefly .
We can start with complex Clifford algebra $`\mathrm{l}\mathrm{l}(n,)`$ with $`n`$ generators $`e_i`$ with relations $`e_ie_j+e_je_i=2\delta _{ij}`$. For even $`n=2k`$ there are two useful properties: first, there are $`k`$ pairs $`a_l=(e_{2l}+ie_{2l+1})/2`$, $`a_l^{}=(e_{2l}ie_{2l+1})/2`$ those satisfy Eq. (1), i.e. CAR. Second, there is recursive construction of $`\mathrm{l}\mathrm{l}(n+2,)=\mathrm{l}\mathrm{l}(n,)\mathrm{l}\mathrm{l}(2,)`$: if $`e_i^{(n)}`$ ($`i0`$) are $`n`$ generators of $`\mathrm{l}\mathrm{l}(n,)`$ and $`e_0^{(2)}`$, $`e_1^{(2)}`$ are generators of $`\mathrm{l}\mathrm{l}(2,)`$ with $`e_{01}^{(2)}ie_0^{(2)}e_1^{(2)}`$, then $`e_i^{(n+2)}=e_i^{(n)}e_{01}^{(2)}`$, $`e_n^{(n+2)}=\mathrm{𝟏}_{}^{(n)}e_0^{(2)}`$, $`e_{n+1}^{(n+2)}=\mathrm{𝟏}_{}^{(n)}e_1^{(2)}`$ are $`n+2`$ generators of $`\mathrm{l}\mathrm{l}(n+2,)`$. Because $`\mathrm{l}\mathrm{l}(2,)`$ is isomorphic with algebra of Pauli matrices, the constructions correspond to equations above if $`e_0^{(2)}\sigma _x`$, $`e_1^{(2)}\sigma _y`$, $`e_{01}^{(2)}\sigma _z`$.
Let us now consider similar possibility for Bose particles with annihilation and creation operators $`c`$, $`c^{}`$ with commutation relation $`[c,c^{}]cc^{}c^{}c=1`$.
$$[c_i,c_j]=[c_i^{},c_j^{}]=0,[c_i,c_j^{}]=\delta _{ij}$$
(4)
The relation $`[A,B]=1`$ cannot be true for finite-dimensional matrices $`A,B`$ because $`\text{trace}[A,B]=\text{trace}(AB)\text{trace}(BA)=0\text{trace}(1)`$. On the other hand infinite-dimensional algebra with the property can be simple found. For example for algebra of functions on line it may be linear operators<sup>1</sup><sup>1</sup>1Cf. with well known representation of $`ip`$ and $`q`$ operators in quantum mechanics.:
$$\begin{array}{c}𝐃:\psi (x)\psi ^{}(x),𝐗:\psi (x)x\psi (x);\\ [𝐃,𝐗]\psi (x)𝐃\left(𝐗\left(\psi (x)^{}\right)^{}\right)𝐗\left(𝐃\left(\psi (x)^{}\right)^{}\right)=\\ =\left(x\psi (x)^{}\right)^{}x\psi ^{}(x)=\left(\psi (x)+x\psi (x)_{}^{}{}_{}{}^{}\right)x\psi ^{}(x)=\mathrm{𝟏}\psi (x)\end{array}$$
(5)
It is possible to define commutation relations on algebra of square-integrable functions with scalar product:
$$\psi \phi =\underset{\mathrm{}}{\overset{+\mathrm{}}{}}\overline{\psi (x)}\phi (x)dx$$
Then $`𝐗^{}=𝐗`$, $`𝐃^{}=𝐃`$ and so $`c=(𝐗+𝐃)/\sqrt{2}`$ and $`c^{}=(𝐗𝐃)/\sqrt{2}`$ satisfy commutation relation for simplest case of one variable $`[c,c^{}]=1`$. The operators is also related with quantization of harmonic oscillator — there is recurrence relation for eigenfunctions $`\psi _n(x)`$ of the 1D oscillator: $`\psi _n(x)=c^{}\psi _{n1}(x)/\sqrt{n}`$ and so $`c^{}`$ really corresponds to ‘rudimentary’ creation operator i.e. excitation of the system on next level.
For more general case of Eq. (4) it is possible to define $`n`$ pairs of operators $`c_i`$, $`c_i^{}`$ on algebra of functions with $`n`$ variables $`\psi (x_1,\mathrm{},x_2)`$:
$$c_i=\frac{1}{\sqrt{2}}\left(x_i+\frac{}{x_i}\right),c_i^{}=\frac{1}{\sqrt{2}}\left(x_i\frac{}{x_i}\right)$$
(6)
Now let us consider more formal representation of CCR used in secondary quantization for case of Bose statistics . Here is used occupation numbers representation $`|\Psi =|n_1,n_2,\mathrm{}`$ with $`n_i`$ is number of particles in $`i`$-th state. Then operators $`c_i`$ and $`c_i^{}`$ are formally defined as:
$$\begin{array}{ccc}\hfill c_i|n_1,n_2,\mathrm{},n_i,\mathrm{}& =& \sqrt{n_i}|n_1,n_2,\mathrm{},n_i1,\mathrm{}\hfill \\ \hfill c_i^{}|n_1,n_2,\mathrm{},n_i,\mathrm{}& =& \sqrt{n_i+1}|n_1,n_2,\mathrm{},n_i+1,\mathrm{}\hfill \end{array}$$
(7)
where definition of conjugated operator meets the standard condition i.e$`n_i|c_i^{}|n_i1=n_i1|c_i|n_i^{}`$.
In computational terminology it can be considered as quantum version of ancient<sup>2</sup><sup>2</sup>2Very ancient, because in Abriss der geschichte der mathematik (A Brief Review of the History of Mathematics) by von Dirk J. Struik (Berlin 1963) is mentioned that The Rhind Mathematical Papyrus written about 3650 years ago already used more or less directly both decimal and binary number systems. counting device with ‘unary’<sup>3</sup><sup>3</sup>3I ‘borrowed’ the term unary from Seth Lloyd (exchange about \[†\]), to avoid ambiguity of word ‘unitary’ in given context. number system (vs. binary or decimal) — number $`n`$ is represented as ‘$`\underset{n}{\underset{}{\mathrm{}}}`$’. Let us call it $`\stackrel{\mathrm{}}{\stackrel{´}{\alpha }}\beta \alpha \xi `$ (abacus).
## 3 Quantum infinite Turing machine
In classical theory of recursion the Turing machine is supplied with infinite tape. Usually it is not considered as practical linitation due to argument: let us start with finite tape and we always may add new sections (cells) if head of Turing machine going to reach the end of the tape.
In quantum computation instead of section of the tape with two states we have two-dimensional Hilbert space $`_2`$, instead of finite tape with $`n`$ sections here is tensor power of $``$:
$$𝖳_n()^n=\underset{n}{\underset{}{\mathrm{}}}$$
Then an analogue of discussed operation may be construction of space<sup>4</sup><sup>4</sup>4It is tensor algebra of $``$. :
$$𝖳_{}()\underset{k=1}{\overset{\mathrm{}}{}}𝖳_k()=()()\mathrm{}$$
(8)
sometime it is convenient to extend summation for $`k=0`$ with $`𝖳_0()`$.
Here subspace $`𝖳_k()`$ in direct sum corresponds to Turing machine with $`k`$-qubits tape. We considered only tape of quantum Turing machine, because it is enough for future description of CCR. The bounded quantum Turing machine also uses other elements and more than one tape and in addition to the property it should be mentioned that $`𝖳_{}()`$ maybe more appropriate for description of semi-infinite Turing tape with necessity of extension of the summation in Eq. (8) for negative $`k`$ (cf. also $`K`$-theory) or using two half-tapes for each infinite tape.
Because it is quantum system we could have superposition of two states with different number of qubits<sup>5</sup><sup>5</sup>5In physics may exist superselection laws those prohibit some superpositions., and here is important that states with different number of qubits belong to orthogonal subspaces.
It is possible to introduce Hermitian scalar product by summation of each component, i.e. for $`\mathrm{\Psi }=_k\psi __k`$, $`\mathrm{\Phi }=_k\phi __k`$ where $`\mathrm{\Psi },\mathrm{\Phi }𝖳_{}()`$ and $`\psi __k,\phi __k𝖳_k()`$:
$$\mathrm{\Psi }\mathrm{\Phi }=\underset{k=0}{\overset{\mathrm{}}{}}\psi __k\phi __k$$
(9)
## 4 Symmetric qubits
Let us now consider symmetric spaces $`𝖲_k()`$ . For two-dimensional $``$ with basis $`e_0`$, $`e_1`$ the space is produced by $`k+1`$ elements $`e_{00\mathrm{}00}`$, $`e_{00\mathrm{}01}`$, …, $`e_{01\mathrm{}11}`$, $`e_{11\mathrm{}11}`$. Let us use notation $`e_{\{i,ki\}}`$ for the elements: $`e_{\{k,0\}}`$, $`e_{\{k1,1\}}`$, …, $`e_{\{0,k\}}`$. It could be enough to use only one index $`e_{\{i\}}`$, $`i=0,\mathrm{},k`$ if we would work with space $`𝖲_k()`$, but it is not enough for $`𝖲_{}()`$ defined as:
$$𝖲_{}()\underset{k=1}{\overset{\mathrm{}}{}}𝖲_k()=()()\mathrm{}$$
(10)
where ‘$``$’ is used for symmetric product.
Due to isomorphism of symmetric space $`𝖲_k`$ with space of homogeneous $`k`$-polynomials elements of $`𝖲_k()`$ can be represented as homogeneous polynomials with two variables $`\xi `$, $`\eta `$: $`p_k(\xi ,\eta )=a_k\xi ^k+a_{k1}\xi ^{k1}\eta +\mathrm{}+a_0\eta ^k`$ and then element $`𝖲_{}()`$ corresponds to arbitrary (non-homogeneous) polynomial with two variables $`p(\xi ,\eta )`$; $`e_{\{i,j\}}\xi ^i\eta ^j`$.
It should be mentioned also, that $`𝖲_k()`$ can be considered as space of polynomials with one variable and with degree less or equal than $`k`$ i.e$`p_k(\zeta )=a_k\zeta ^k+a_{k1}\zeta ^{k1}+\mathrm{}+a_0`$ and there is special case with one or more higher coefficients are zeros $`a_i=0`$, $`ki>l`$. If $`\zeta _1,\mathrm{},\zeta _l`$, $`lk`$ are roots of the polynomial, then factorization of the $`p_k(\zeta )=a_l\mathrm{\Pi }_{i=1}^l(\zeta \zeta _i)`$ corresponds to factorization of $`p_k(\xi ,\eta )`$ on $`k`$ terms:
$$a_k\xi ^k+a_{k1}\xi ^{k1}\eta +\mathrm{}+a_0\eta ^k=(\alpha _1\xi \beta _1\eta )\times \mathrm{}\times (\alpha _k\xi \beta _k\eta )$$
where $`\alpha _i\zeta _i=\beta _i`$ and maybe $`\alpha _i=0`$ if $`l<k`$ and $`i>l`$ i.e. formally $`\zeta _i=\mathrm{}`$, but here is more rigorously to use projective spaces<sup>6</sup><sup>6</sup>6Projective spaces make possible to get rid of special cases like $`\zeta _i=\mathrm{}`$ if $`a_k\mathrm{}=0`$. — spaces of rays in terminology more usual for quantum mechanics. Projective coordinate $`\zeta `$ corresponds to ray $`(\xi ,\eta )(\lambda \xi ,\lambda \eta )`$ in $``$ or point on Riemann (Bloch) sphere and roots $`\zeta _i`$ correspond to pairs $`(\alpha _i,\beta _i)(\lambda \alpha _i,\lambda \beta _i)`$.
Due to the property element of $`𝖲_k()`$ up to multiplier is defined by $`k`$ points $`(\alpha _i,\beta _i)`$ from $``$: $`\mathrm{\Pi }_{i=1}^k(\alpha _i\xi \beta _i\eta )=\lambda p_k(\xi ,\eta )`$ and because multiplication of each pair $`(\alpha _i,\beta _i)(\lambda _i\alpha _i,\lambda _i\beta _i)`$ changes only common multiplier $`\lambda `$ for same $`p_k`$, it is correct map from $`k`$ rays in $``$ to ray in $`𝖲_k()`$.
So, symmetric qubits are never entangled, each element of $`𝖲_k()`$ can be represented as symmetrical product of $`k`$ qubits, elements of $``$. Because $`𝖲_k()`$ is $`k+1`$-dimensional linear space and can be used as space of states for particle with spin $`k/2`$, the factorization described before explains why state of the particle always can be described as $`k`$ points on Riemann sphere<sup>7</sup><sup>7</sup>7There is popular introduction and two references in \[10, §6; Objects with large spin\]..
Note: the model of symmetric qubits may looks like some violation of Pauli’s spin-statistics principle. Really, one qubit can be considered as fermion, especially in the context of the paper with anticommutation relation and Pauli’s matrices. It should be said for justification, that it is very convenient mathematical model of spin-$`\frac{k}{2}`$ system originated by Weyl, Majorana, Penrose and the symmetric spin-$`\frac{1}{2}`$ subsystems can be considered as formal math ‘ghosts’ (or ‘colored’ like quarks). Yet another reason — the qubit is an abstract two-states system and it is not quite correct to talk about spins and statistics, it maybe electron with spin half or photon with spin one, or some model described by Schrödinger equation with potential well, i.e. by one-component, scalar wave function. And next, the coordinate dependence is not considered in usual models of qubit and so Pauli’s exclusion principle sometime can be formally avoided by suggestion about different locations for each qubit. Cf. also misc. 2 at end of Sec. 5.
Now let us define scalar product on $`𝖲_k`$ and $`𝖲_{}`$. It is convenient together with basis $`e_{\{i,j\}}`$, $`i+j=k`$ to consider:
$$\stackrel{~}{e}_{\{i,j\}}\frac{e_{\{i,j\}}}{\sqrt{i!j!}}=\sqrt{\frac{C_k^i}{k!}}e_{\{i,j\}}$$
(11)
The basis is convenient by following reasons: First, in ‘more physical’ definition the $`𝖲_k()`$ is subspace ($`𝖲_k𝖳_k`$) of symmetrical tensors with operation of symmetrization by summation of $`k!`$ transpositions $`\sigma (T)`$ of indexes for given $`T𝖳_k`$: $`𝖲(T)=\frac{1}{k!}_\sigma \sigma (T)`$, and if to consider $`e_{\{i,ki\}}`$ as element $`e_{00\mathrm{}11}`$ of $`𝖳_k`$, then $`|𝖲(e_{\{i,ki\}})|=1/\sqrt{C_k^i}`$ (it is sum of all $`C_k^i`$ possible transpositions with coefficient $`1/C_k^i`$) and needs for normalizing multiplier proportional to $`\sqrt{C_k^i}`$ (in is used second definition of symmetric space as quotient space $`𝖲=𝖳/𝔖`$, where $`𝔖`$ is equivalence relation: $`T\sigma (T)`$).
Second, $`\stackrel{~}{e}_{\{i,j\}}`$ form representation of $`SU(2)`$ group in $`SU(k+1)`$ in such a way, that if we use other basis $`U:(e_0,e_1)(e_0^{},e_1^{})`$, $`USU(2)`$ then $`\stackrel{~}{e}_{\{i,j\}}^{}`$ are also connected with $`\stackrel{~}{e}_{\{i,j\}}`$ by some unitary transformation from $`SU(k+1)`$.
Let us use for $`𝖲_k`$ basis $`\stackrel{~}{e}_{\{i,j\}}`$ with standard $``$ (or $`/k!`$, see Eq. (12) below) together with $`e_{\{i,j\}}`$ considered as transformation to other basis with Hermitian scalar product defined by diagonal matrix $`h_{ii}=C_k^i`$ (or $`h_{ii}=C_k^i/k!=\frac{1}{i!(ki)!}`$) — the basis is convenient for representation of $`𝖲_k`$ as space of polynomials.
The scalar product on $`𝖲_{}()`$ may be defined as in Eq. (9), but it is more convenient to use also:
$$\mathrm{\Psi }^𝖲\mathrm{\Phi }^𝖲_{\mathrm{exp}}=\underset{k=0}{\overset{\mathrm{}}{}}\frac{1}{k!}\psi __k^𝖲\phi __k^𝖲$$
(12)
## 5 Quantum $`\stackrel{\mathbf{}}{\stackrel{\mathbf{´}}{𝜶}}𝜷𝜶𝝃`$
It is possible to use $`𝖲_{}()`$ as an example of quantum $`\stackrel{\mathrm{}}{\stackrel{´}{\alpha }}\beta \alpha \xi `$ (let us denote it as $`|\stackrel{\mathrm{}}{\stackrel{´}{\alpha }}\beta \alpha \xi `$ or $`|\stackrel{\mathrm{}}{\stackrel{´}{\alpha }}`$) introduced in Sec. 2. The approach makes possible to link it with $`𝖳_{}()`$ and quantum Turing machine, $`𝖲(|\text{TM tape})|\stackrel{\mathrm{}}{\stackrel{´}{\alpha }}\beta \alpha \xi `$
The elements $`e_{\{i,j\}}`$ or $`\stackrel{~}{e}_{\{i,j\}}`$ of basis $`𝖲_{}()`$ can be used as basis $`|i,j`$ of $`|\stackrel{\mathrm{}}{\stackrel{´}{\alpha }}`$ with two different kinds of states $`n_0=i`$, $`n_1=j`$ that can be considered also as composite system of two $`|\stackrel{\mathrm{}}{\stackrel{´}{\alpha }}`$ with infinite series of states for each one: $`𝖲_{}()|\stackrel{\mathrm{}}{\stackrel{´}{\alpha }}_0|\stackrel{\mathrm{}}{\stackrel{´}{\alpha }}_1`$.
It is possible to introduce operators $`c_i`$, $`c_i^{}`$; $`i=0,1`$ by Eq. (7). Let us $`\stackrel{~}{e}_{\{n_0,n_1\}}|n_0,n_1`$. Then:
$$\begin{array}{cccccc}\hfill c_0|n_0,n_1& =& \sqrt{n_0}|n_01,n_1,\hfill & \hfill c_0^{}|n_0,n_1& =& \sqrt{n_0+1}|n_0+1,n_1\hfill \\ \hfill c_1|n_0,n_1& =& \sqrt{n_1}|n_0,n_11,\hfill & \hfill c_1^{}|n_0,n_1& =& \sqrt{n_1+1}|n_0,n_1+1\hfill \end{array}$$
(13)
Here numbers of zeros and units $`n_0,n_1`$ are used instead of $`n_1,n_2,\mathrm{}`$
As understanding example of such system it is possible to use two-dimensional oscillator:
$$i\mathrm{}\dot{\psi }(x,y,t)=\left(\frac{m\omega ^2}{2}(x^2+y^2)\frac{\mathrm{}^2}{2m}\mathrm{\Delta }_{x,y}\right)\psi (x,y,t)$$
If $`\varphi _k(x)`$, $`k0`$ is stationary solution of one-dimensional oscillator for energy $`E=(k+\frac{1}{2})\mathrm{}\omega `$, then for 2D oscillator function $`\varphi _k(x)\varphi _j(y)`$ is solution for $`E=(k+j+1)\mathrm{}\omega `$ and so for any natural $`n0`$ there is $`n+1`$ dimensional space of solutions for given energy $`E=(n+1)\mathrm{}\omega `$:
$$\varphi _n(x,y)=\underset{k=0}{\overset{n}{}}\alpha _k\varphi _k(x)\varphi _{nk}(y)$$
and nonstationary solution has form:
$$\varphi (x,y,t)=\underset{n=0}{\overset{\mathrm{}}{}}A_n\varphi _n(x,y)e^{i(n+1)\mathrm{}\omega t}$$
The example shows, how tensor product of two infinite-dimensional spaces $`_{\mathrm{}}_{\mathrm{}}`$ is decomposed on direct sum of linear spaces with dimensions $`1,2,3,\mathrm{}`$ for the simple case in good agreement with formal mathematical constructions discussed above.
It is useful to consider Eq. (13) in basis $`e_{\{n_0,n_1\}}`$:
$$\begin{array}{ccc}\hfill c_0e_{\{n_0,n_1\}}& =& c_0|n_0,n_1\sqrt{n_0!n_1!}=\sqrt{n_0}|n_01,n_1\sqrt{n_0!}\sqrt{n_1!}\hfill \\ & =& n_0|n_01,n_1\sqrt{n_01!}\sqrt{n_1!}=n_0e_{\{n_01,n_1\}}\hfill \\ \hfill c_0^{}e_{\{n_0,n_1\}}& =& c_0^{}|n_0,n_1\sqrt{n_0!n_1!}=\sqrt{n_0+1}|n_0+1,n_1\sqrt{n_0!}\sqrt{n_1!}\hfill \\ & =& |n_0+1,n_1\sqrt{n_0+1!}\sqrt{n_1!}=e_{\{n_0+1,n_1\}}\hfill \end{array}$$
(14)
and similarly with $`c_1,c_1^{}`$ and $`n_1`$.
The equations Eq. (14) demonstrate relation between CCR in secondary quantization Eq. (7) with CCR in differential algebra Eq. (5). Really, let us consider polynomials with two variables $`\chi _0,\chi _1`$, then $`e_{\{i,j\}}\chi _0^i\chi _1^j`$, $`c_ip(\chi _0,\chi _1)=\frac{}{\chi _i}p(\chi _0,\chi _1)`$, $`c_i^{}p(\chi _0,\chi _1)=\chi _ip(\chi _0,\chi _1)`$.
### Miscellany
1. Let us return to initial notation ($`\xi =\chi _0`$, $`\eta =\chi _1`$) for polynomial basis $`e_{\{i,j\}}\xi ^i\eta ^j`$ defined earlier in Sec. 4. The linear operators $`c_0`$ and $`c_1`$ are isomorphic with two partial derivatives $`_\eta `$ and $`_\xi `$ and have some interesting property, ‘linear merging’ (‘anti-cloning’).
It was described in Sec. 4 that state of symmetric qubits up to multiplier, i.e. ray in $`𝖲_n()`$, can be described by $`n`$ points on Riemann sphere. The operators $`_\xi `$ and $`_\eta `$ can be considered as maps $`𝖲_n()𝖲_{n1}()`$ and also between spaces of rays — from sphere with $`n`$ marked points $`(\zeta _1,\mathrm{},\zeta _n)`$ and without one pole (the pole maps to zero) to sphere without same pole and with $`n1`$ marked points $`(\zeta _1^{},\mathrm{},\zeta _{n1}^{})`$. For $`n=2`$ they map $`𝖲_2()`$ and $`(\zeta _1,\zeta _2)(\zeta _1^{})`$. If two points on Riemann sphere coincide $`\zeta _1=\zeta _2`$, then due to standard property of differential we have $`\zeta _1^{}=\zeta _1=\zeta _2`$ and so we have maps $`(\zeta _1,\zeta _1)(\zeta _1)`$ described by linear operators $`_\xi `$ or $`_\eta `$.
Cloning is suggested may not to be linear, but for symmetric qubits ‘opposite’ operation may be defined as linear one almost everywhere except of one point and it may be arbitrary point of Riemann sphere because we can use operator $`\stackrel{}{}_vv_1_\xi +v_2_\eta `$.
2. Let us consider even subspace of $`𝖲_{}()`$ defined as: $`𝖲_{}^2()_{k=1}^{\mathrm{}}𝖲_{2k}()`$ and $`𝖳_{}^2()_{k=1}^{\mathrm{}}𝖳_{2k}()`$. Such spaces may be more appropriate for taking into account the Pauli’s exclusion principle, but here is discussed only simplest illustrative example.
Let us introduce operators $`\stackrel{~}{𝐃}c_0c_1`$ and $`\stackrel{~}{𝐗}\xi \eta `$ i.e. :
$$\stackrel{~}{𝐃}|n_0,n_1=\sqrt{n_0n_1}|n_01,n_11,\stackrel{~}{𝐗}|n_0,n_1=|n_0+1,n_1+1$$
The operators act on $`𝖲_{}^2()`$ and there is isomorphism of subspace with basis $`|n,n`$, i.e$`n_0=n_1`$ and space of functions with operators Eq. (5), $`|n,nx^n`$.
The operators are described here not only because of trivial identity $`\sqrt{n_0n_1}=n`$ for $`n_0=n_1=n`$, but as an illustrative introduction to more difficult 4D case where momentum operator $`p_i`$ and differentials $`\frac{}{x_i}`$ have two indexes as spinoral objects. But it is already away from theme of the paper<sup>8</sup><sup>8</sup>8It may be suggested, that in the words about simulation of bosons by spin lattice are also related with same area of research i.e. ‘Feynman checkerboard model’, ‘Penrose spin network’, ‘Finkelstein-Selesnick quantum net’ etc\[‡\]..
## 6 Conclusions and discussion
In the paper is discussed algebraic approach to quantum computational models with unlimited number of discrete states. Similarly with classical recursion theory instead of ‘actual infinity’ here is considered quantum analogue of Turing machine as sequence of systems with increasing number of states.
Here is also used simplest mathematical model with linear spaces, isomorphic to space of polynomials. It maybe defined also by canonical commutation relations, CCR. It is still far from more rigorous models like relativistic quantum theories of interacting fields, but it is some step in the direction.
Usual quantum networks correspond to physical approach with S-matrix. It is analogue of ideas cited in introduction. There are initial state, scattering process described by ‘quantum black box’ and final state. In such picture we have only two ‘points’ (in $``$ out) instead of 4D spacetime.
Because CCR are also related with algebra of smooth functions on spacetime (see Eq. (5)) the models discussed in the paper are useful possibility to take into account some properties of temporo-spatial systems in quantum networks approach. It is only necessary to accept infinite or dynamically changing number of qubits.
Does the mathematical models like $`𝖳_{}(_2)`$ with infinite sequence of linear spaces with increasing dimensions devote an attention instead of ‘actual’ infinite-dimensional space $`_{\mathrm{}}`$? The paper is an attempt to make positive answer. For example operators of creation and annihilation can be simply expressed via transition between spaces $`𝖲_k()`$ of such sequence.
A nontrivial property of such sequences is orthogonality of any states in different terms. For example a state of tape of quantum Turing machine is orthogonal with state that differs only on extra one empty section, but in classical case they are considered as the same. The quantum $`\stackrel{\mathrm{}}{\stackrel{´}{\alpha }}\beta \alpha \xi `$, $`𝖲_{}()`$ has same property, but it is more clear from physical model of secondary quantization, because $`|n_0,n_1`$ and $`|n_0+1,n_1`$ are obviously orthogonal. But $`𝖲_{}()`$ is simply symmetrization $`𝖲(𝖳_{}())`$ and so model of quantum Turing machine as infinite sequence of orthogonal linear spaces is not much more unusual than 2D harmonic oscillator discussed as a model of $`𝖲_{}()`$ in Sec. 5.
## Acknowledgements
I am grateful to David Deutsch for some useful exchange and for inspiration my interest to quantum information science in relation with reading \[\*\] few years ago. Many thanks to Seth Lloyd for interesting communication. Certainly, understanding of the particular area of quantum mechanics would be much less effective for me without big help of David Finkelstein with explanation and discussion about some general principles of philosophy of the quantum World. |
warning/0001/nucl-ex0001005.html | ar5iv | text | # Experimental Search for Molecular-Nuclear Transitions in WaterLANL e-print nucl-ex/0001005
## I Introduction
Phenomenon of a long-range effective interaction in some quantum systems including constituents which are only interacting at short distances is known already for a long time. Here, by a short-range interaction we understand any interaction that falls off at large distances between particles, i. e. as $`r\mathrm{}`$, not slower than exponentially. In the context below, longe-range forces decrease at infinity as some inverse powers of $`r`$.
The most pronounced manifestation of this phenomenon appears in the so-called “Efimov effect” . The Efimov effect consists, in particular, in arising an effective interaction between a particle and a coupled pair of particles which behaves as $`1/r^2`$. A sufficient condition for this effect to take place in a three-bosonic system is a closeness of the binding energies at least for two of pair subsystems to zero. Just the very small binding energies force the two-body wave function to be extremely extended generating an effective long-range interaction of the coupled pair with a complementary particle.
Another example of a long-range effect of short-range interaction was found by Ya. B. Zeldovich . He showed that if two particles interact via two potentials, a short-range one and an arbitrary long-range (say, Coulomb) one, then the spectrum of the two-body system can be drastically changed as compared to the case of a long-range potential only. This change takes place if the short-range potential generates a bound or resonant state with an energy sufficiently close to zero.
Thus, in both the above cases, long-range effects take place due to the fact that a short-range interaction causes a two-body system to be, nevertheless, very extended. Having this in mind, one can expect some enhancement for the probability of transition of the system from a molecular state to an extended (resonance) nuclear state as compared to a similar transition into a localized nuclear state. Indeed, model calculations of the overlap integral between wave functions of the H<sub>2</sub>O ($`1^{}`$) molecule and a resonant state ($`4.522,\mathrm{\hspace{0.17em}1}^{}`$) of the nucleus <sup>18</sup>Ne show an enhancement of this sort.
The purpose of the present experiment is to estimate the life time of water molecules with respect to the following decay chain
$$[\mathrm{H}_2\mathrm{O}]_1^{}{}_{}{}^{18}\mathrm{Ne}_{}^{}(4.522,1^{}){}_{}{}^{18}\mathrm{Ne}(\mathrm{g}.\mathrm{s}.){}_{}{}^{18}\mathrm{F}{}_{}{}^{18}\mathrm{O}+\beta ^+.$$
## II Experimental approach and results
A number of examples of nuclear systems with near-threshold resonances were analyzed from this point of view, among them are (p, p ,<sup>16</sup>O) and (p, <sup>17</sup>O) , i. e. the nuclear constituents of the usual water molecule H<sub>2</sub>O and hydroxyl ions OH based on the rare oxygen isotope <sup>17</sup>O (see Table I).
For the fist experimental study in this direction we choose the system (p, p ,<sup>16</sup>O), i. e. the H<sub>2</sub>O molecule. Its properties, in addition to obvious availability, make this case the most favorable for the experiment. From Fig. 1, displaying the energy-lever diagram for <sup>18</sup>Ne-nucleus, it is seen that rotation states 1<sup>-</sup> of a water molecule and a highly excited state (1<sup>-</sup>, 4.522 MeV) of this nucleus can be considered as energy degenerate. Thus, a real physical state with these quantum numbers appears to be a superposition of molecular and nuclear states. Experimental approach was designed with taking into account of the properties of both the components of the superposition pair.
As a rule, rotational states are excited only in free molecules while their population under the conditions of condensed-phase water is prohibited due to powerful hydrogen bonds . This makes carrying out the experiments more difficult, but at the same time it gives additional opportunities to manipulate with the hypothetical process of molecular-nuclear transitions in a water system. So, it becomes possible to use an accumulation-measuring cycling during searching for resulting radioactive products. The layout of measurements is shown in Fig. 2. The accumulation cycle represented heating the water portion within a sealed out measuring chamber to a critical point about 647 K, at which a total amount of water was for sure in a vapor phase regardless of pressure within the volume. To withstand this rather a high pressure (22.5 MPa ), the stainless-steel or titanium chambers were strengthened for the heating period by two thick steel removable plates at the top and bottom of the chambers. The latter were thick-wall shortened metallic cylinders, sealed at the faces by thin membranes almost fully transparent for the expected annihilation radiation $`E_\gamma =511`$ keV.
For the measuring period, the plates were removed, and the chamber was placed after cooling between two NaI(Tl)-scintillators, operating in the $`\gamma \gamma `$-coincidence mode. First, the measurements were performed in surface-laboratory conditions at the Institute of Physics and Technology Problems (Dubna), then the main part of the experiments was undertaken, this time at a deep underground laboratory of the Baksan Neutrino Observatory (Republic of Kabardino-Balkaria, the North Caucasus) of the Institute for Nuclear Research of the Russian Academy of Sciences. At the surface-laboratory stage, necessary measurement procedures were optimized, and a yield of molecular-nuclear transition in water in the condensed state was estimated. Analysis of these data showed main sources of the background counting: true coincidences due to the cosmic muons, decay of the natural <sup>40</sup>K and daughter products of <sup>222</sup>Rn (<sup>214</sup>Bi and <sup>214</sup>Po). To take into account the background in the region of “interest”, i. e., near $`E_\gamma =511`$ keV, some control background measurements were carried out under the identical conditions (geometry, amount of water, etc.), in which heavy water (D<sub>2</sub>O, 99.0%-enrichment) was used instead of the natural one. Within the limits of statistical fluctuations, both the spectra were identical. For the water half-life in the condensed-phase state with respect to decay by the chain H<sub>2</sub>O$``$<sup>18</sup>Ne($`\beta ^+`$; 1.7 s)$`^{18}`$F($`\beta ^+`$; 109 min), a lower limit was estimated as T$`{}_{1/2}{}^{}`$4.10<sup>21</sup> years (within the 99%-confidence level).
At the Baksan Observatory, all the efforts were made to subdue the background radiation as much as possible. The cosmic muons were subdued due to power shielding, as the measuring premises were situated inside a gallery created within the mountain Andyrchi (one of mountains of the Elbrus environment, where the screening thickness of a rock achieves about 600 m of water equivalent. Walls (the ceiling and floor, as well) were built of a special uranium-free concrete. The scintillators were prepared on the basis of materials with a low containment of potassium and radium. They were, for a long time, located in a deep underground. Thus all short-lived cosmic-ray-induced activities should be extinct. The detection unit was provided with an additional shielding composed of interchanged layers of pure tungsten, lead, and copper.
For every $`\gamma `$-quantum event, the time was registered with a channel width 10 s. This was more than sufficient for a detailed analysis of the time dynamics in the range of interest for this experiment. The accumulation interval was chosen to be about three half-lifes of <sup>18</sup>F, i. e. 5.5 hours. The measuring time-schedule for these runs is inserted in Fig. 3. Intervals of about four half-lives of <sup>18</sup>F, i. e. $`440`$ min each, were considered as time-periods when the decay of hypothetically accumulated <sup>18</sup>F could yet give some contribution to the counting in the region of $`E_\gamma =511`$ keV. The remaining part of measuring cycles was used to calculate the background counting in the same energy range.
Preliminary analysis of the results was carried out by means of comparing counting rates in two above measurement periods, each of a total duration $`4000`$ min for ten cycles. A certain excess (at a level about 1-2 RMS deviation) of the counting rate in the region of $`E_\gamma =511`$ keV was observed in the first period as compared to the second one, i. e. to the background. We emphasize that the above excess was non-stationary. Time dependence of the non-stationary process approximately corresponded to the half-life of <sup>18</sup>F. Under the assumption that the observed non-stationary component of the effect is associated with the accumulation and decay of the nuclei <sup>18</sup>F, our estimate for the efficient half-life of the water molecule with respect to the nuclear channel under consideration is $`T_{1/2}10^{18}`$ years.
Although the total statistics was not sufficient to make a decisive conclusion, the results obtained are rather encouraging for further experiments in this intriguing direction with new more sensitive approaches. More attention should be paid to analysis of population rates of suitable molecular states under the experimental conditions. Also some further theoretical consideration are desirable regarding the fusion probability estimates for this and other molecular systems.
###### Acknowledgements.
This work was supported by the Russian Foundation for Basic Research (grant # 98-02-16884). |
warning/0001/hep-th0001049.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The spectrum of supersymmetric string theories usually contains a special class of states known as BPS states, which are characterized by the property that their mass is completely determined by their charge under some gauge field. They form short (or ultra short) supersymmetric multiplets and, because of this fact, are stable and protected from quantum radiative corrections. A well-known example of such BPS states is provided by the supersymmetric D$`p`$-branes of the type II theories (with $`p`$ even in type IIA and $`p`$ odd in type IIB) . However, supersymmetric string theories quite often contain states that are stable without being BPS. These are in general the lighest states which carry some conserved quantum numbers. For them there is no particular relation between their mass and their charge; they form long multiplets of the supersymmetry algebra and receive quantum radiative corrections. However, being the lightest states with a given set of conserved quantum numbers, they are stable since they cannot decay into anything else. Usually, it is not difficult to find such non-BPS states with the standard string perturbative methods and analyze their properties at weak coupling; but, since they cannot decay, they should be present also in the strong coupling regime, or equivalently they should appear as non-perturbative (D-brane type) configurations in the weakly coupled dual theory. To verify the existence of these non-BPS states is therefore a very strong test on the duality relations between two string theories which does not rely on supersymmetry arguments. The study of the stable non-BPS D-branes in string theory, pioneered by A. Sen in a remarkable series of papers , has attracted a lot of interest during the last year (for reviews see Refs. ) also for several other reasons; among them we recall the fact the non-BPS D-branes might be useful for analyzing the non-perturbative properties of the non-supersymmetric field theories that live on their world-volumes, or the fact that they may lead to novel types of relations among string theories .
One of the most notable examples of stable non-BPS configurations is provided by the perturbative states at the first excited level of the $`SO(32)`$ heterotic string which carry the spinor representation of the gauge group and whose mass is given by
$$M_\mathrm{h}=\frac{2}{\sqrt{\alpha ^{}}}=\frac{g_{\mathrm{YM}}}{\kappa _{10}}.$$
(1.1)
In the last equality we have introduced the low-energy gauge and gravitational couplings $`g_{\mathrm{YM}}`$ and $`\kappa _{10}`$ of the heterotic string following the conventions of Ref. which are also reviewed in Appendix B. Being at the first massive level, these states are non-BPS, but being the lightest ones carrying the spinor representation of $`SO(32)`$, they are stable and should be present also when one increases the heterotic string coupling constant $`g_\mathrm{h}`$. In this process however, the mass $`M_\mathrm{h}`$ gets renormalized since there are no constraints on it coming from supersymmetry. Thus we can write
$$M_\mathrm{h}=\frac{g_{\mathrm{YM}}}{\kappa _{10}}f,$$
(1.2)
where the renormalization function $`f`$ can in principle be computed perturbatively in the heterotic string and is such that $`f1`$ for $`g_\mathrm{h}0`$.
If the heterotic/type I duality is correct, also in the type I theory there should exist stable non-BPS configurations that are spinors of $`SO(32)`$. Such states do indeed exist and were identified by A. Sen as the non-BPS D-particles of type I ; then, an explicit boundary state description for them was provided in Ref. <sup>1</sup><sup>1</sup>1For the description of other non-BPS D-branes using the boundary state formalism see Refs. .. The mass of these D-particles turns out to be
$$M_{\stackrel{~}{0}}=\frac{1}{\sqrt{\alpha ^{}}g_\mathrm{I}}=\frac{g_{\mathrm{YM}}}{\kappa _{10}}2^{3/4}g_{\mathrm{I}}^{}{}_{}{}^{1/2},$$
(1.3)
where $`g_\mathrm{I}`$, $`g_{\mathrm{YM}}`$ and $`\kappa _{10}`$ are, respectively, the string, the gauge and the gravitational coupling constants of the type I theory in the conventions of Ref. (see also Appendix B). Comparing Eqs. (1.2) and (1.3), and remembering that under the duality map the heterotic gauge and gravitational couplings turn into the corresponding ones of type I, we can deduce that the renormalization function $`f`$ must be such that $`f2^{3/4}g_{\mathrm{I}}^{}{}_{}{}^{1/2}`$ for $`g_\mathrm{I}0`$ in order for the masses to agree on both sides. Clearly, this result cannot be obtained using perturbative methods, but is a prediction of the heterotic/type I duality <sup>2</sup><sup>2</sup>2The T-dual heterotic/type I’ correspondence has been analyzed at the non-BPS level in Ref. .
In this paper we elaborate further on these stable non-BPS particles and study in detail their gravitational and gauge interactions. On the heterotic side, these can be easily obtained using standard perturbative techniques from correlation functions of vertex operators. In this way one can show, for example, that, at the lowest order in the heterotic string coupling constant, the gravitational and gauge potential energies of two such particles at large distance are given, respectively, by the Newton’s law and the Coulomb’s law for massive and charged point-like objects in ten dimensions. On the type I side, instead, the interactions of the non-BPS D-particles must be obtained using less standard methods and have not been fully investigated so far; indeed, only the general rules for computing string amplitudes with these D-particles have been given in the literature . It is the purpose of this paper to fill this gap.
In particular, we will concentrate on processes involving massless string modes that are responsible for the long range interactions among D-particles. To study the gravitational interactions we adopt the boundary state formalism and obtain the energy due to the exchange of closed string states between two D-particles by simply computing the diffusion amplitude between the two corresponding boundary states in relative motion (for recent reviews on the boundary state formalism and its applications see Ref. ). Then, by taking the large distance limit to which only graviton and dilaton exchanges contribute, we find that the gravitational potential energy of two D-particles exactly agrees with the one of their heterotic duals, provided that the duality relations and the mass renormalization previously discussed are taken into account.
For the gauge interactions, instead, the situation is a bit more involved. In fact, we cannot use any more the boundary state formalism since this accounts only for the couplings of the D-particles with the closed strings that live in the bulk, but is completely blind to the other bulk sector of the type I theory consisting of open strings with Neumann boundary conditions in all directions to which the $`SO(32)`$ gauge fields belong. On the other hand, the open strings attached to the D-particles have Neumann boundary conditions only along the time direction. Therefore, to study the gauge interactions of our D-particles we should consider scattering amplitudes involving open strings with mixed boundary conditions in an odd number of dimensions. Calculations of open string amplitudes with mixed boundary conditions have already appeared in the analyses of systems of several D-branes with different dimensionality (see for instance Ref. ), and require the use of twist operators to produce mixed boundary conditions in certain directions. These twist operators were used in the past to study strings on orbifolds , and have been recently reconsidered from an abstract conformal field theory point of view . Using such twist operators and applying the rules of Refs. , we will describe how to compute scattering amplitudes involving non-BPS D-particles and bulk open strings of type I. Special care is required in these calculations because the twist operators that we use change the boundary conditions in an odd number of directions. In particular, we will explicitly determine the gauge coupling of the D-particles by evaluating a correlation function on a disk with two boundary components produced by the insertion of two twist operators. The result of this calculation is extremely simple, namely the non-BPS D-particles couple minimally to the gauge field. Exploiting this fact, we then determine the gauge potential energy of a pair of D-particles at large distance and see that after taking into account the duality map, this exactly agrees with the corresponding energy computed in the heterotic theory.
This paper is organized as follows: in Section 2 we compute the gauge coupling of a non-BPS D-particle of type I by evaluating a disk diagram with two twist insertions, and then determine the gauge potential energy between two D-particles. In Section 3 we use the boundary state formalism to compute the gravitational contribution to the potential energy of two (moving) D-particles. In Section 4 we study the gauge and gravitational interactions of the non-BPS heterotic states that are dual to the D-particles. In Section 5 we compare the results for the non-BPS D-particles and for their dual heterotic states, and discuss their relations. In Appendix A we show how to compute the gauge interactions between two BPS D-strings of type I by extending the method of Section 2 and verify the no-force condition. Finally, Appendices B and C contain our conventions and a list of more technical formulas.
## 2 Type I D-particle interactions: the gauge amplitude
As we mentioned in the introduction, an important check of the heterotic/type I duality has been the discovery by A. Sen that the stable non-BPS heterotic states carrying the spinor representation of $`SO(32)`$ at the first massive level are dual to the non-BPS D-particles of type I. Specific rules for computing amplitudes involving such D-particles have been given by A. Sen and E. Witten in two different ways which we briefly recall here. Sen’s approach heavily relies of the use of Chan-Paton factors to distinguish the various kinds of open strings. The 0-0 strings, whose end-points lie on the non-BPS D-particle, contain both states that are even and states that are odd under $`(1)^F`$; the former carry a Chan-Paton factor 1l, the latter a Chan-Paton factor $`\sigma _1`$. The 9-0 strings stretching between one of the 32 D9 branes of the type I background and a D-particle contain only $`(1)^F`$ even states but, due to the existence of an odd number of fermionic zero modes, their vertex operators comprise the standard GSO-even part as well as the corresponding GSO-odd part, weighted by Chan-Paton factors $`\left(\begin{array}{c}1\\ 0\end{array}\right)`$ and $`\left(\begin{array}{c}0\\ 1\end{array}\right)`$ respectively. Besides these factors, the 9-0 strings also carry a Chan-Paton factor $`\lambda ^A`$ ($`A=1,\mathrm{},32`$) labeling the fundamental representation of the $`SO(32)`$ gauge group. Finally, the 9-9 strings are the usual open strings of the type I theory which are GSO projected and carry only the standard Chan-Paton factors of the gauge group.
The presence of the unusual Chan-Paton factors 1l, $`\sigma _1`$, $`\left(\begin{array}{c}1\\ 0\end{array}\right)`$ or $`\left(\begin{array}{c}0\\ 1\end{array}\right)`$ shows that the states of the 0-0 and 9-0 sectors have a non trivial structure which is really due to the presence of an odd number of fermionic zero modes. In order to remedy to this oddity, in Ref. Witten has proposed to introduce an extra one-dimensional fermion $`\eta `$ on each boundary of the string world-sheet lying on a D-particle. In this way, in the 9-0 sector one recovers an even number of fermionic zero modes and can perform the usual GSO projection. Also in the 0-0 sector one performs a (generalized) GSO projection to obtain physical states, but since the extra fermion $`\eta `$ is odd under this GSO parity, one obtains two types of 0-0 states, similarly to what found by Sen.
Let us now give some details on how to construct the massless states in the various open string sectors using Witten’s rules. We start with the NS sector of the 0-0 strings where at the massless level there are nine scalars $`x^i`$ ($`i=1,\mathrm{},9`$) corresponding to the freedom of moving the D-particle in its nine tranverse directions. These modes, which are present also on the BPS D0 brane of the type IIA theory, correspond to vertex operators $`𝒱_{x^i}`$ that do not depend on the boundary fermion $`\eta `$. In the $`(1)`$ superghost picture, these vertex operators are simply
$$𝒱_{x^i}^{(1)}=\psi ^i\mathrm{e}^\varphi .$$
(2.1)
Notice that there is no factor of $`\mathrm{e}^{\mathrm{i}kX}`$ in (2.1) because massless states of $`0`$-$`p`$ strings have no momentum. Let us now consider the R sector of the 0-0 strings. Here both the ten world-sheet fermions $`\psi ^\mu `$ and the boundary fermion $`\eta `$ possess zero modes so that the massless R states form a GSO-even spinor $`\xi ^\alpha `$ of $`SO(1,10)`$. Note that in this case the GSO projection is simply the ten-dimensional chirality projection which is natural when one extends $`SO(1,9)`$ to $`SO(1,10)`$ by adding $`\eta `$. Thus, in the $`(1/2)`$ superghost picture the vertex operator for the massless R states reads
$$𝒱_{\xi ^\alpha }^{(1/2)}=\frac{1+\eta }{2}S^\alpha \mathrm{e}^{\varphi /2}$$
(2.2)
where $`S^\alpha `$ is the spin field of conformal dimension $`10/16`$ associated to the ten world-sheet fermions. Upon quantization, the $`16`$ massless fermionic modes described by (2.2) account for the $`2^{16/2}=256`$ degeneracy of the non-BPS D-particle.
We now turn to the 9-0 strings which are more relevant for our purposes. Since the NS sector does not contain massless states, we just consider the R sector. In this case, the only world sheet fermion to have a zero mode is $`\psi ^0`$ so that the ground state is a GSO-even (chiral) spinor of the algebra $`SO(1,1)`$ generated by $`\psi _0^0`$ and $`\eta `$. Hence, the vertex operator describing the massless modes of the 9-0 sector should contain
* a spin field $`S`$ associated to the fermion $`\psi ^0`$, of conformal dimension 1/16;
* a boundary changing operator $`\mathrm{\Delta }`$ for the nine space directions transverse to the D-particle, of conformal dimension 9/16;
* a GSO (or chirality) projector for the Clifford algebra of $`SO(1,1)`$ $`\frac{1+\psi _0^0\eta }{2}`$, of conformal dimension zero;
* a superghost contribution in the $`1/2`$ picture $`\mathrm{e}^{\varphi /2}`$, of conformal dimension 3/8;
* a gauge Chan-Paton factor $`\lambda ^A`$ to specify which of the 32 D9 branes one is considering.
Thus, we have
$$𝒱_{90}^{(1/2)}=\lambda ^A\frac{1+\psi _0^0\eta }{2}S\mathrm{\Delta }\mathrm{e}^{\varphi /2}.$$
(2.3)
It is easy to check that the operator (2.3) has indeed conformal dimension 1 as it should be for a physical vertex operator. Notice that the GSO projection in (2.3) keeps only one fermionic degree of freedom for each value of the index $`A`$ of the fundamental representation of $`SO(32)`$. Upon quantization, the states described by $`𝒱_{90}`$ form a spinorial representation of $`SO(32)`$, and hence we can conclude that the marginal operator (2.3) accounts for the $`SO(32)`$ degeneracy of the non-BPS D-particle. Since in type I the strings are unoriented, we should consider also the 0-9 sector. This is merely related to the 9-0 sector through the action of the world-sheet parity $`\mathrm{\Omega }`$. Recalling that $`\mathrm{\Omega }`$ simply acts by transposition on the Chan-Paton factors without changing the physical content of the vertex operators, we have
$$𝒱_{09}^{(1/2)}=\mathrm{\Omega }𝒱_{90}^{(1/2)}=\lambda _{}^{t}{}_{}{}^{A}\frac{1+\psi _0^0\eta }{2}S\mathrm{\Delta }\mathrm{e}^{\varphi /2}.$$
(2.4)
Notice in particular that the $`SO(1,1)`$ GSO projection is the same in both vertices (2.3) and (2.4).
Finally, there are the 9-9 strings which, as we mentioned above, are the usual open strings of the type I theory; in particular in the NS sector at the massless level we find the $`SO(32)`$ gauge bosons which are described by the following vertex operators in the $`(1)`$ superghost picture
$$𝒱_{\mathrm{gauge}}^{(1)}=\mathrm{\Lambda }^{AB}A_\mu \psi ^\mu \mathrm{e}^{\mathrm{i}kX}\mathrm{e}^\varphi ,$$
(2.5)
where $`\mathrm{\Lambda }^{AB}`$ are the generators of $`SO(32)`$ in the fundamental representation (see Appendix B for our conventions) and $`A_\mu `$ is the polarization vector.
We now face the problem of finding the coupling between the non-BPS D-particle and the gauge field. Since the latter belongs to the 9-9 massless sector, the diagram we have to compute corresponds to a disk with a part of its boundary on the D-particle and a part on the D9-branes from which the gauge boson is emitted. This is represented in the Figure 1.
We thus have to insert two vertices containing the boundary changing operator $`\mathrm{\Delta }`$ which turns a boundary of type 0 into one of type 9 (or viceversa). The obvious choice is then to make insertions of the vertices $`𝒱_{90}`$ and $`𝒱_{09}`$ given in (2.3) and (2.4) which correspond to the $`SO(32)`$ degeneracy of the D-particle. Note that, although no momentum is carried by these vertices, the emitted gauge boson may have non-zero space momentum. Indeed, the twist operators $`\mathrm{\Delta }`$ are reservoirs of transverse momentum which allow emissions with non-zero momentum in the transverse directions; on the other hand, this is to be expected because the presence of a D-brane breaks the translational invariance in transverse space. Note also that, due to the insertion of the two vertices $`𝒱_{90}`$ and $`𝒱_{09}`$, the boundary component associated with the D-particle carries indices in the bi-fundamental representation of $`SO(32)`$. This is consistent with the fact that a D-particle should emit all, both massless and massive, perturbative open string states which group in the adjoint or in the symmetric representation of $`SO(32)`$.
The gauge coupling of a (static) D-particle is then given by the expectation value of the gauge boson emission vertex (2.5) in the “vacuum” representing the D-particle. Thus, the diagram of Figure 1 corresponds to
$$^{\mathrm{gauge}}=\mathrm{\Omega }_{09}|𝒱_{\mathrm{gauge}}^{(1)}|\mathrm{\Omega }_{90}$$
(2.6)
where
$$|\mathrm{\Omega }_{90}=\underset{z0}{lim}𝒱_{09}^{(1/2)}(z)|0\mathrm{and}\mathrm{\Omega }_{09}|=\underset{z\mathrm{}}{lim}0|𝒱_{09}^{(1/2)}(z).$$
(2.7)
Note that due to the presence of the twist operators in $`𝒱_{09}`$ and $`𝒱_{90}`$, the expectation value of the gauge emission vertex is not vanishing, as we will explicitly see in the following.
After including the normalization factor $`𝒞_{\mathrm{disk}}`$ appropriate of any disk amplitude, the normalization factors $`𝒩_\mathrm{R}`$ for the R vertices (2.3) and (2.4), and $`𝒩_{\mathrm{NS}}`$ for the NS vertex (2.5<sup>3</sup><sup>3</sup>3We refer to Appendix B for the explicit expression of these normalization factors and to Refs. for their derivation., $`^{\mathrm{gauge}}`$ may be reexpressed as a 3-point function on the world sheet and reads
$$^{\mathrm{gauge}}=𝒞_{\mathrm{disk}}𝒩_\mathrm{R}^2𝒩_{\mathrm{NS}}𝑑\mu (z_i)c𝒱_{09}^{(1/2)}(z_1)c𝒱_{\mathrm{gauge}}^{(1)}(z_2)c𝒱_{90}^{(1/2)}(z_3)_\eta .$$
(2.8)
where we have also added a ghost $`c`$ in each vertex operator. The notation $`_\eta `$ means that the correlator must be evaluated by including the action for the boundary fermion $`\eta `$ as explained in Ref. .
The correlation function in (2.8) may be decomposed into a longitudinal and a transverse piece. The latter vanishes because
$$\mathrm{\Delta }(z_1)\psi ^i(z_2)\mathrm{e}^{\mathrm{i}kX(z_2)}\mathrm{\Delta }(z_3)=0$$
(2.9)
for $`i=1,\mathrm{},9.`$ Thus, there is no emission of gauge bosons with polarization $`A_i`$ along the transverse directions, as it should be for a minimally coupled particle at rest. We then consider the longitudinal part for which the basic correlators are
$`c(z_1)c(z_2)c(z_3)`$ $`=`$ $`z_{12}z_{13}z_{23},`$ (2.10)
$`\mathrm{e}^{\varphi (z_1)/2}\mathrm{e}^{\varphi (z_2)}\mathrm{e}^{\varphi (z_3)/2}`$ $`=`$ $`z_{12}^{1/2}z_{13}^{1/4}z_{23}^{1/2},`$ (2.11)
$`\mathrm{\Delta }(z_1)\mathrm{e}^{\mathrm{i}kX(z_2)}\mathrm{\Delta }(z_3)`$ $`=`$ $`z_{12}^{\alpha ^{}k^2}z_{13}^{\alpha ^{}k^29/8}z_{23}^{\alpha ^{}k^2}.`$ (2.12)
Notice that in (2.12) the transverse momentum $`k^i`$ of the emitted gauge boson is not subject to any constraint, as we have anticipated. The remaining correlator to be considered is
$$\left(\frac{1+\psi _0^0\eta }{2}S(z_1)\right)\psi ^0(z_2)\left(\frac{1+\psi _0^0\eta }{2}S(z_3)\right)_\eta .$$
(2.13)
This splits into four pieces, two of which vanish. Indeed, according to Ref. the only non-vanishing correlation functions are those containing one factor of $`\eta `$. In particular one has
$$\eta _\eta =\sqrt{2},1_\eta =0.$$
(2.14)
Finally, we have
$$\psi _0^0S(z_1)\psi ^0(z_2)S(z_3)=S(z_1)\psi ^0(z_2)\psi _0^0S(z_3)=z_{12}^{1/2}z_{13}^{3/8}z_{23}^{1/2}.$$
(2.15)
Notice that a correlation function similar to (2.15) appears in the 2D Ising model. Indeed, the spin field $`S`$ may be identified with the order parameter $`\sigma `$ (i.e. the magnetization) while the other spin field $`\psi _0^0S`$ plays the role of the disorder parameter $`\mu `$ .
Inserting Eqs. (2.10)-(2.15) into (2.8) and exploiting the projective invariance to fix the position of the three punctures at arbitrary values, we easily get
$`^{\mathrm{gauge}}`$ $`=`$ $`{\displaystyle \frac{𝒞_{\mathrm{disk}}𝒩_\mathrm{R}^2𝒩_{\mathrm{NS}}}{\sqrt{2}}}\mathrm{Tr}(\lambda _{}^{t}{}_{}{}^{A}\mathrm{\Lambda }^{BC}\lambda ^D)A_0{\displaystyle 𝑑\mu (z_i)\left(\frac{z_{13}}{z_{12}z_{23}}\right)^{\alpha ^{}k^2}}`$ (2.16)
$`=`$ $`{\displaystyle \frac{𝒞_{\mathrm{disk}}𝒩_\mathrm{R}^2𝒩_{\mathrm{NS}}}{\sqrt{2}}}\mathrm{Tr}(\lambda _{}^{t}{}_{}{}^{A}\mathrm{\Lambda }^{BC}\lambda ^D)A_0.`$
Then, using the explicit expressions of the normalization coefficients and Chan-Paton factors given in Appendix B, we can rewrite $`^{\mathrm{gauge}}`$ as follows
$$^{\mathrm{gauge}}=\mathrm{i}\frac{g_{\mathrm{YM}}}{\sqrt{2}}\left(\delta ^{AB}\delta ^{CD}\delta ^{AC}\delta ^{BD}\right)A_0,$$
(2.17)
where $`g_{\mathrm{YM}}`$ is the gauge coupling constant of the type I theory.
Eq. (2.17) represents the amplitude for the emission of a gauge boson with longitudinal polarization $`\xi _0`$ and color index $`(BC)`$ from a 0-boundary in the bi-fundamental of $`SO(32)`$. The appearance of this representation is a direct consequence of our construction in which the D-particle is represented by the 0-component of a disk boundary produced by the insertion of the vertex operators $`𝒱_{90}`$ and $`𝒱_{09}`$. On the other hand, the $`SO(32)`$ spinor degeneracy of the non-BPS D-particle of type I arises from the (second) quantization of the 32 massless fermionic zero-modes of the 0-9 open strings, and thus it is clear that such a degeneracy cannot be seen in our operator formalism. This fact should not be surprising because a completely analogous situation occurs in the familiar description of D$`p`$-branes using boundary states. Indeed, a boundary state is a single state that correctly represents a D$`p`$-brane and its couplings to the bulk closed strings, even if it does not account for the degeneracy of the D$`p`$-brane under the supersymmetry algebra. Similarly, in our case the 0-component of the disk boundary produced by the insertion of $`𝒱_{90}`$ and $`𝒱_{09}`$ correctly describes a D-particle and allows to obtain its coupling with the bulk 9-9 open strings, even if does not account for its degeneracy under the gauge group. In fact, as we will see later and in the following sections, using this construction we are able to obtain non trivial information about the gauge interactions between two D-particles at large distance. Moreover, after taking into account the known duality relations, we will show that the results obtained in this way exactly agree with those in the heteroric theory, as required by the heterotic/type I duality, thus confirming the validity of our construction.
Using the result (2.17) we can now easily compute the gauge potential energy $`V_\mathrm{I}^{\mathrm{gauge}}`$ due to the exchange of the $`SO(32)`$ gauge bosons between two D-particles. As indicated in Figure 2,
this can be obtained simply by sewing two emission amplitudes $`^{\mathrm{gauge}}`$ with the gauge boson propagator
$$𝒫=\frac{\eta ^{\alpha \beta }}{q^2}\left(\delta ^{BB^{}}\delta ^{CC^{}}\delta ^{BC^{}}\delta ^{CB^{}}\right),$$
(2.18)
yielding
$$V_\mathrm{I}^{\mathrm{gauge}}=\frac{g_{\mathrm{YM}}^2}{2}\left(\delta ^{AE}\delta ^{FD}\delta ^{AF}\delta ^{DE}\right)\frac{1}{q^2}.$$
(2.19)
Performing a Fourier transform, we get the following (static) gauge potential in configuration space
$$V_\mathrm{I}^{\mathrm{gauge}}(r)=\frac{g_{\mathrm{YM}}^2}{2}\left(\delta ^{AE}\delta ^{FD}\delta ^{AF}\delta ^{DE}\right)\frac{1}{7\mathrm{\Omega }_8r^7},$$
(2.20)
where $`\mathrm{\Omega }_q=2\pi ^{(q+1)/2}/\mathrm{\Gamma }((q+1)/2)`$ is the area of a unit $`q`$-dimensional sphere. Eq. (2.20) clearly represents a “Coulomb-like” potential energy for point particles at a distance $`r`$ in ten dimensions.
We conclude this section by mentioning that the same results (2.17) and (2.20) can be obtained also using the rules given by A. Sen in Ref. for computing amplitudes with non-BPS D-particles.
## 3 Type I D-particle interactions: the gravitatio- nal amplitude
The gravitational contribution to the scattering of two non-BPS D-particles of type I can be calculated, at the leading order in the string coupling constant, from the diffusion amplitude between two corresponding boundary states. The boundary state description of the non-BPS D-particles has already been given in Refs. from which we recall the results that are relevant in the forthcoming analysis. For details and conventions on boundary states, we refer the reader to Refs. .
In the closed string operator formalism, one describes a D$`p`$-brane by means of a boundary state $`|Dp`$ . This is a closed string state which inserts a boundary on the world-sheet, enforces on it the appropriate boundary conditions and represents the source for the closed strings emitted by the brane. As an example, the boundary state for a BPS D-particle of type IIA may formally be written as <sup>4</sup><sup>4</sup>4In order to avoid clutter, we shall denote the NS-NS (resp. R-R) component of a boundary state with the simplified subscript NS (resp. R)
$$|D0_{\mathrm{IIA}}=|D0_{\mathrm{NS}}+|D0_\mathrm{R},$$
(3.1)
where the NS-NS and the R-R components are both proportional to $`T_0`$ which is the tension of the D-particle in units of the gravitational coupling constant, namely
$$T_0=8\pi ^{7/2}\alpha ^{\mathrm{\hspace{0.17em}3}/2}.$$
(3.2)
The presence of both the NS-NS and the R-R components implies that the spectrum of the open strings living on the D-particle is GSO-projected. The partition function of such open strings may be obtained by evaluating the cylinder/annulus amplitude in the closed string channel which is given by
$${}_{\mathrm{IIA}}{}^{}D0|P|D0_{\mathrm{IIA}}^{},$$
(3.3)
and then performing a modular transformation. In Eq. (3.3), $`P`$ denotes the closed string propagator
$`P={\displaystyle \frac{\alpha ^{}}{2}}{\displaystyle \frac{1}{L_0+\stackrel{~}{L}_0a\stackrel{~}{a}}}.`$ (3.4)
where $`a`$ ($`\stackrel{~}{a}`$) is the left (right) intercept ($`a_{\mathrm{NS}}=1/2`$, $`a_\mathrm{R}=0`$).
The boundary state for the non-BPS D-particle of the IIB theory has instead only a component along the NS-NS sector and a tension $`T_{\stackrel{~}{0}}`$ greater by a factor of $`\sqrt{2}`$ than $`T_0`$. Thus, we can write
$$|D\stackrel{~}{0}_{\mathrm{IIB}}=\sqrt{2}|D0_{\mathrm{NS}}.$$
(3.5)
As a consequence, there is no GSO-projection in the spectrum of the open strings lying on the non-BPS D-particle and the presence of a tachyon in the NS sector renders it unstable. However, if we consider the type I theory , the tachyon is removed by the projection onto states invariant under the world-sheet parity $`\mathrm{\Omega }`$ . In the boundary state formalism, the $`\mathrm{\Omega }`$ projection is implemented by adding the so-called crosscap state $`|C`$ , which corresponds to inserting on the closed string world-sheet a boundary with opposite points identified. The negative $`(32)`$ charge for the non-propagating R-R 10-form that the crosscap generates in the background, must be canceled by the introduction of 32 D9-branes. Hence, the type I theory possesses a background “boundary state” given by
$$\frac{1}{\sqrt{2}}\left(|C+32|D9\right)$$
(3.6)
where the factor of $`1/\sqrt{2}`$ has been introduced to obtain the right normalization of the various spectra. Then, the partition function for unoriented open 9-9 strings, given by the sum of the annulus and the Möbius strip contributions, is
$$\frac{1}{2}\left(2^{10}D9|P|D9+2^5D9|P|C+2^5C|P|D9\right),$$
(3.7)
while the contribution of the Klein bottle
$$\frac{1}{2}C|P|C$$
(3.8)
added to the torus contribution gives the partition function for unoriented closed strings.
The boundary state of the non BPS D-particle of type I reads
$$|D\stackrel{~}{0}_I=\frac{1}{\sqrt{2}}\left(\sqrt{2}|D0_{\mathrm{NS}}\right)=|D0_{\mathrm{NS}}$$
(3.9)
where we have added the same factor of $`1/\sqrt{2}`$ for consistency with (3.6). The mass of the D-particle is then given by
$`M_{\stackrel{~}{0}}={\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \frac{T_{\stackrel{~}{0}}}{\kappa _{10}}}={\displaystyle \frac{T_0}{\kappa _{10}}}={\displaystyle \frac{1}{\sqrt{\alpha ^{}}g_I}}`$ (3.10)
where $`\kappa _{10}`$ is the ten dimensional gravitational coupling constant of the type I theory (see Appendix B). The partition function for open 0-0 strings living on the D-particle, obtained by summing the contributions from the annulus and the Möbius strip, is
$${}_{\mathrm{NS}}{}^{}D0|P|D0_{\mathrm{NS}}^{}+\frac{1}{2}(\sqrt{2}C|P|D0_{\mathrm{NS}}+\sqrt{2}{}_{\mathrm{NS}}{}^{}D0|P|C).$$
(3.11)
In this theory, there are also 0-9 and 9-0 open strings with one end on the D-particle and the other on one of the 32 D9 branes of the type I background. The world-sheet parity $`\mathrm{\Omega }`$ exchanges the two sectors 0-9 and 9-0 so that we only retain symmetric combinations corresponding to the partition function
$$\frac{32\sqrt{2}}{2}({}_{\mathrm{NS}}{}^{}D0|P|D9+D9|P|D0_{\mathrm{NS}}).$$
(3.12)
The spectrum of open strings stretching between two different (distant) D-particles at rest, one labeled with a prime, has a partition function given by
$$\frac{1}{2}((\sqrt{2})^2{}_{\mathrm{NS}}{}^{}D0|P|D0^{}_{\mathrm{NS}}^{}+(\sqrt{2})^2{}_{\mathrm{NS}}{}^{}D0^{}|P|D0_{\mathrm{NS}}^{})$$
(3.13)
where the factor of one-half indicates that, compared to the IIB case, only the $`\mathrm{\Omega }`$ symmetric combinations are retained. Notice that, at sufficiently small distance, a tachyon develops in this open string spectrum signaling the instability of the configuration which decays into the vacuum .
Our aim is to study the diffusion of a moving D-particle with a velocity $`v`$ along one space direction, say $`X^1`$, on another D-particle at rest at the origin. Such an interaction may be evaluated analyzing the spectrum of the open strings stretching between the two objects with modified boundary conditions in the 0 and 1 directions. This can be done generalizing the treatment for the BPS D-branes presented in Ref. , but we find it simpler to use the method of the boosted boundary state . Indeed, the interaction amplitude just reads
$$𝒜_\mathrm{I}(v)={}_{\mathrm{NS}}{}^{}D0^{}|P\mathrm{\Lambda }|D0_{\mathrm{NS}}^{}+{}_{\mathrm{NS}}{}^{}D0|\mathrm{\Lambda }^{}P|D0^{}_{\mathrm{NS}}^{}$$
(3.14)
where $`\mathrm{\Lambda }`$ is the boost operator
$$\mathrm{\Lambda }=\mathrm{e}^{\mathrm{i}\pi \nu J^{01}}$$
(3.15)
acting on the boundary state of a particle at rest. Here we have $`v=\mathrm{th}(\pi \nu )`$ and $`J^{\mu \nu }`$ is the generator of the Lorentz transformations. Notice that the amplitude (3.14) reduces to the static one (3.13) in the limit of vanishing velocity. The boosted boundary state <sup>5</sup><sup>5</sup>5The signs $`\pm `$ correspond to the two possible implementations of boundary conditions for world-sheet fermions. In a physical (GSO projected) boundary state, only a suitable linear combination of them is retained. reads
$`\mathrm{\Lambda }|D0,\pm _{\mathrm{NS}}`$ $`=`$ $`{\displaystyle \frac{T_0}{2}}{\displaystyle \frac{1}{\gamma }}\delta ^{(8)}(x)\delta (x^0v+x^1)\mathrm{exp}\left[{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}a_n^{}S\stackrel{~}{a}_n^{}\right]`$ (3.16)
$`\times \mathrm{exp}\left[\pm \mathrm{i}{\displaystyle \underset{r=\frac{1}{2}}{\overset{\mathrm{}}{}}}\psi _r^{}S\stackrel{~}{\psi }_r^{}\right]{\displaystyle \underset{\mu =0}{\overset{9}{}}}|k^\mu =0`$
where the boundary conditions are encoded in the matrix $`S=(𝒱_{01},\text{ 1l}_8)`$ with
$`𝒱_{01}=\left(\begin{array}{cc}\mathrm{cosh}(2\pi \nu )& \mathrm{sinh}(2\pi \nu )\\ \mathrm{sinh}(2\pi \nu )& \mathrm{cosh}(2\pi \nu )\end{array}\right).`$ (3.19)
Note that $`\mathrm{cosh}(\pi \nu )\gamma `$ is the Lorentz factor.
The interaction amplitude $`𝒜_\mathrm{I}(v)`$ can be evaluated using standard techniques and explicitly reads <sup>6</sup><sup>6</sup>6See for instance Ref. for definitions and conventions about the modular functions $`f_k`$ and $`\theta _k`$.
$`𝒜_\mathrm{I}(v)`$ $`=`$ $`(8\pi ^2\alpha ^{})^{\frac{1}{2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\tau {\displaystyle _0^{\mathrm{}}}𝑑ss^{\frac{9}{2}}\mathrm{e}^{\frac{b^2+\tau ^2v^2\gamma ^2}{2\pi \alpha ^{}s}}`$ (3.20)
$`\mathrm{\hspace{0.17em}2}\mathrm{sinh}(\pi \nu )\left[{\displaystyle \frac{f_3^6(q)\theta _3(\mathrm{i}\nu |\mathrm{i}s)f_4^6(q)\theta _4(\mathrm{i}\nu |\mathrm{i}s)}{f_1^6(q)\theta _1(\mathrm{i}\nu |\mathrm{i}s)}}\right]`$
in which $`q=\mathrm{e}^{\pi s}`$, $`\tau `$ is the proper time of the moving particle and $`b`$ is the impact parameter. We are now in a position to extract the long range interaction potential energy $`V_\mathrm{I}^{\mathrm{grav}}`$ due to gravitational exchange between the two particles. To do so we have to perform the limit $`s\mathrm{}`$ in the integrand of Eq. (3.20), then integrate on the variable $`s`$ and finally identify the potential energy according to
$$𝒜_\mathrm{I}(v)=_{\mathrm{}}^{\mathrm{}}𝑑\tau V_\mathrm{I}^{\mathrm{grav}}.$$
(3.21)
In the non relativistic limit, we obtain the Newton’s law with its first correction
$$V_\mathrm{I}^{\mathrm{grav}}(r)=(2\kappa _{10})^2\frac{M_{\stackrel{~}{0}}^2}{7\mathrm{\Omega }_8r^7}\left(1+\frac{1}{2}v^2\right)+o(v^2)v0$$
(3.22)
where we have introduced the radial coordinate $`r^2=b^2+v^2\gamma ^2\tau ^2`$. Thus, in the non relativistic limit the boundary state calculation reproduces correctly the gravitational potential energy that we expect for a pair of D-particles in relative motion.
## 4 Interactions of the heterotic non-BPS states
The non-BPS D-particles described in the previous sections account for the presence in the spectrum of the type I theory of long super-multiplets of states carrying the spinorial representation of $`SO(32)`$. These non-perturbative states are dual to those appearing at the first massive level in the heterotic theory. Carrying the same quantum numbers, one naturally expects that these heterotic states have the same kind of interactions as the D-particles of type I. In this section we will check this idea and investigate the gauge and gravitational interactions of the non-BPS heterotic states using standard tools of perturbative string theory. In doing so, we will adopt the bosonized formulation of the heterotic string in which the gauge degrees of freedom are described by sixteen chiral bosons $`\stackrel{~}{X}^I`$ ($`I=1,\mathrm{},16`$) appropriately compactified .
The long super-multiplet of the stable heterotic states appears at the first massive level $`\left(M_\mathrm{h}^2=4/\alpha ^{}\right)`$, and contains the following bosonic states
$`\psi _{3/2}^\mu |k|K^I`$ , $`\alpha _1^{(\mu }\psi _{1/2}^{\nu )}|k|K^I,`$ (4.1)
$`\alpha _1^{[\mu }\psi _{1/2}^{\nu ]}|k|K^I`$ , $`\psi _{1/2}^\mu \psi _{1/2}^\nu \psi _{1/2}^\rho |k|K^I.`$ (4.2)
with $`\mu ,\nu ,\mathrm{}=0,\mathrm{},9`$. In these formulas $`k`$ denotes the space-time momentum ($`k^2=M_\mathrm{h}^2`$) while $`K^I`$ is the adimensional momentum associated to the sixteen internal coordinates $`\stackrel{~}{X}^I`$. The states of Eq. (4.1) describe massive degrees of freedom which transform in the 44 representation of the Lorentz group, whereas those of Eq. (4.2) transform in the 84 <sup>7</sup><sup>7</sup>7 The fermionic states that complete this long multiplet transform in the 128 representation of the Lorentz group.. The level matching condition requires that $`K^2=4`$. This may be realized for example by taking $`K^I`$ to be of the form $`(\pm \frac{1}{2},\pm \frac{1}{2},,\mathrm{},\pm \frac{1}{2})`$ with an even number of $`+`$ signs, thus obtaining the spinorial representation of $`SO(32)`$ with positive chirality. The vertex operators for the states (4.1) and (4.2) will be denoted by $`𝒱`$ and can be found in Appendix C both in the $`(1)`$ and in the $`(0)`$ superghost pictures.
We now study the interactions of these states with the massless gauge bosons of $`SO(32)`$. In the bosonized formulation of the heterotic string we must distinguish between the states associated to the $`16`$ Cartan generators that are given by
$$A_\mu \psi _{1/2}^\mu |q\stackrel{~}{\alpha }_1^I|Q=0\text{with}q^2=0\text{and}I=1,\mathrm{},16,$$
(4.3)
and those associated to the remaining $`480`$ generators which are instead given by
$$A_\mu \psi _{1/2}^\mu |q|Q\text{with}q^2=0$$
(4.4)
and the internal momentum $`Q`$ of the form $`(0,\mathrm{},0,\pm 1,0,\mathrm{},0,\pm 1,0,\mathrm{},0)`$. Also the vertex operators for the states (4.3) and (4.4), which we denote collectively by $`𝒱_{\mathrm{gauge}}`$, can be found in Appendix C in the $`(1)`$ and $`(0)`$ superghost pictures.
The gauge coupling of the states (4.1) and (4.2) is obtained by simply computing the 3-point function on the sphere among two vertex operators $`𝒱`$ and one vertex operator $`𝒱_{\mathrm{gauge}}`$ (see Figure 3).
Including the normalization factor $`𝒞_0`$ appropriate of any tree-level closed string amplitude and a normalization factor $`\widehat{𝒩}`$ for each vertex operator, we have
$$\left(_3^{\mathrm{gauge}}\right)_{\alpha _1\alpha _2}=𝒞_0\widehat{𝒩}^3d^2\mu (z_i,\overline{z}_i)c\overline{c}𝒱_1^{(1)}(z_1,\overline{z}_1)c\overline{c}𝒱_2^{(1)}(z_2,\overline{z}_2)c\overline{c}𝒱_{\mathrm{gauge}}^{(0)}(z_3,\overline{z}_3)$$
(4.5)
where $`\alpha _1`$ and $`\alpha _2`$ label the spinor representation of $`SO(32)`$ carried by the non-BPS states and a ghost factor $`c\overline{c}`$ has been added in each puncture. Actually, we are not interested in the complete expression of this correlation function but only in the scalar part of it, namely in the terms where the polarizations $`\zeta _1`$ and $`\zeta _2`$ of the two spinor states are contracted between themselves <sup>8</sup><sup>8</sup>8The polarization $`\zeta _i`$ can be either a vector, a symmetric or antisymmetric two-index tensor or an antisymmetric three-index tensor depending on which particular states (4.1) and (4.2) are considered.. This is because we want to compare our results with those of the non-BPS D-particles of type I obtained in the previous sections in which the Lorentz group structure was not manifest.
Using the explicit expression of the vertex operators reported in Appendix C, we find that the terms of (4.5) proportional to $`\zeta _1\zeta _2`$ are given by
$$\left(_3^{\mathrm{gauge}}\right)_{\alpha _1\alpha _2}^I=\frac{𝒞_0\widehat{𝒩}^3\sqrt{2\alpha ^{}}}{4}(\zeta _1\zeta _2)A_\mu (k_1^\mu k_2^\mu )(K_{\alpha _1}K_{\alpha _2})^I\delta _{K_{\alpha _1}+K_{\alpha _2},0}$$
(4.6)
when $`𝒱_{\mathrm{gauge}}`$ corresponds to the gauge bosons associated to the Cartan generators, and by
$$\left(_3^{\mathrm{gauge}}\right)_{\alpha _1\alpha _2}^Q=\frac{𝒞_0\widehat{𝒩}^3\sqrt{2\alpha ^{}}}{2}(\zeta _1\zeta _2)A_\mu (k_1^\mu k_2^\mu )C_Q(K_{\alpha _1})\delta _{K_{\alpha _1}+K_{\alpha _2}+Q,0}$$
(4.7)
when $`𝒱_{\mathrm{gauge}}`$ corresponds to the gauge bosons associated to the remaining generators (here $`K_\alpha `$ denotes the lattice vector corresponding to the spinorial state $`\alpha `$ and $`C_Q(K_{\alpha _1})`$ is a cocycle factor; see e.g. Refs. for details). The dependence on the internal momenta looks different in the two expressions (4.6) and (4.7), but a moment thought reveals that it is actually of the same form as required by gauge invariance. Indeed, we can rewrite both equations in the following form
$$\left(_3^{\mathrm{gauge}}\right)_{\alpha _1\alpha _2}^{AB}=\frac{g_{\mathrm{YM}}}{\sqrt{2}}(\zeta _1\zeta _2)A_\mu (k_1^\mu k_2^\mu )\mathrm{\Gamma }_{\alpha _1\alpha _2}^{AB}$$
(4.8)
where we have used the definitions of the normalization coefficients to express the prefactor in terms of the Yang-Mills coupling constant of the heterotic theory. Here $`\mathrm{\Gamma }_{\alpha _1\alpha _2}^{AB}`$ are the matrix elements of the antisymmetrized product of two $`\mathrm{\Gamma }`$-matrices of $`SO(32)`$ which represent the fusion coefficients among the adjoint and two spinor representations of $`SO(32)`$<sup>9</sup><sup>9</sup>9See for instance Ref. for the expression of these $`\mathrm{\Gamma }`$-matrices in terms of the internal momenta and cocycle factors.
We are now in the position of evaluating the contribution to the diffusion amplitude among four particles due to the exchange of gauge bosons at tree level. In fact, this can be simply obtained by sewing two 3-point functions $`_3^{\mathrm{gauge}}`$ with the massless propagator (2.18); in this way we obtain
$$\left(_4^{\mathrm{gauge}}\right)_{\alpha _1\alpha _4;\alpha _2\alpha _3}=\frac{g_{\mathrm{YM}}^2}{2}(\zeta _1\zeta _4)(\zeta _2\zeta _3)\frac{us}{t}\left(\frac{1}{2}\mathrm{\Gamma }_{\alpha _1\alpha _4}^{AB}\mathrm{\Gamma }_{\alpha _2\alpha _3}^{AB}\right)$$
(4.9)
where $`s`$, $`t`$ and $`u`$ are the usual Mandelstam variables which satisfy $`s+t+u=16/\alpha ^{}`$. For later convenience, we introduce the adimensional variable $`\sigma =u\alpha ^{}/4`$, which in the limit $`t0`$ is related to the Lorentz parameter $`\gamma `$ according to $`\sigma =2(\gamma 1)`$. Then, for $`t0`$ Eq. (4.9) becomes
$$\left(_4^{\mathrm{gauge}}\right)_{\alpha _1\alpha _4;\alpha _2\alpha _3}=g_{\mathrm{YM}}^2(\zeta _1\zeta _4)(\zeta _2\zeta _3)\frac{M_\mathrm{h}^2(2+\sigma )}{t}\left(\frac{1}{2}\mathrm{\Gamma }_{\alpha _1\alpha _4}^{AB}\mathrm{\Gamma }_{\alpha _2\alpha _3}^{AB}\right).$$
(4.10)
Reverting to the standard field theory normalization by multiplying each external leg by $`1/\sqrt{2E}`$, and removing for simplicity the polarization factors, we finally obtain the gauge potential energy
$$\left(V_\mathrm{h}^{\mathrm{gauge}}\right)_{\alpha _1\alpha _4;\alpha _2\alpha _3}=\frac{g_{\mathrm{YM}}^2}{2}\frac{1}{t}\left(\frac{1}{2}\mathrm{\Gamma }_{\alpha _1\alpha _4}^{AB}\mathrm{\Gamma }_{\alpha _2\alpha _3}^{AB}\right),$$
(4.11)
which in configuration space becomes
$$\left(V_\mathrm{h}^{\mathrm{gauge}}\right)_{\alpha _1\alpha _4;\alpha _2\alpha _3}(r)=\frac{g_{\mathrm{YM}}^2}{2}\left(\frac{1}{2}\mathrm{\Gamma }_{\alpha _1\alpha _4}^{AB}\mathrm{\Gamma }_{\alpha _2\alpha _3}^{AB}\right)\frac{1}{7\mathrm{\Omega }_8r^7}.$$
(4.12)
Notice that this potential does not depend on the relative velocity of the particles involved in the interaction; moreover, as expected, it is a “Coulomb-like” potential for point particles in ten dimensions carrying the spinor representation of the gauge group.
We now turn to the gravitational interactions of the non-BPS heterotic particles (4.1) and (4.2) following the same steps we have described for the gauge interactions. Let us recall that the massless bosonic states of the graviton multiplet of the heterotic theory are
$$ϵ_{\mu \nu }\psi _{1/2}^\mu \stackrel{~}{\alpha }_1^\nu |q|Q=0\text{with}q^2=0,$$
(4.13)
where the polarization is
$$ϵ_{\mu \nu }=h_{\mu \nu }=h_{\nu \mu },q^\mu h_{\mu \nu }=0$$
(4.14)
for the graviton,
$$ϵ_{\mu \nu }=\frac{\varphi }{\sqrt{8}}(\eta _{\mu \nu }q_\mu \mathrm{}_\nu q_\nu \mathrm{}_\mu ),q\mathrm{}=1,\mathrm{}^2=0$$
(4.15)
for the dilaton, and
$$ϵ_{\mu \nu }=\frac{1}{\sqrt{2}}B_{\mu \nu }=\frac{1}{\sqrt{2}}B_{\nu \mu },q^\mu B_{\mu \nu }=0$$
(4.16)
for the antisymmetric Kalb-Ramond field. The vertex operators corresponding to these states are written in Appendix C in the $`(1)`$ and $`(0)`$ superghost pictures, and will be denoted generically by $`𝒱_{\mathrm{grav}}`$.
The gravitational coupling of the non-BPS particles can be determined by evaluating the correlation function among two vertex operators $`𝒱`$ and one vertex operator $`𝒱_{\mathrm{grav}}`$, namely
$$\left(_3^{\mathrm{grav}}\right)_{\alpha _1\alpha _2}=𝒞_0\widehat{𝒩}^3d^2\mu (z_i,\overline{z}_i)c\overline{c}𝒱_1^{(1)}(z_1,\overline{z}_1)c\overline{c}𝒱_2^{(1)}(z_2,\overline{z}_2)c\overline{c}𝒱_{\mathrm{grav}}^{(0)}(z_3,\overline{z}_3).$$
(4.17)
As before, also now we are interested only in the scalar part of this expression which is proportional to $`\zeta _1\zeta _2`$, since we want to compare it with the boundary state calculation of Section 3. Using the expression of the vertex operators reported in Appendix C, it is not difficult to find that
$$\left(_3^{\mathrm{grav}}\right)_{\alpha _1\alpha _2}=𝒞_0\widehat{𝒩}^3\alpha ^{}\left(\zeta _1\zeta _2\right)ϵ_{\mu \nu }k_2^\mu k_2^\nu ,$$
(4.18)
from which we read that the couplings of the non-BPS states with the graviton, the dilaton and the antisymmetric tensor field are
$`\left(_3^{\mathrm{grav}}\right)_{\alpha _1\alpha _2}^{(h)}`$ $`=`$ $`4\kappa _{10}(\zeta _1\zeta _2)h_{\mu \nu }k_2^\mu k_2^\nu ,`$
$`\left(_3^{\mathrm{grav}}\right)_{\alpha _1\alpha _2}^{(\varphi )}`$ $`=`$ $`\sqrt{2}\kappa _{10}M_\mathrm{h}^2(\zeta _1\zeta _2)\varphi ,`$ (4.19)
$`\left(_3^{\mathrm{grav}}\right)_{\alpha _1\alpha _2}^{(B)}`$ $`=`$ $`0.`$
Note that the vanishing of the heterotic non BPS states coupling with the antisymmetric Kalb-Ramond field is consistent with the fact that the type I D-particle does not couple to the R-R 2-form.
Now we can evaluate the diffusion amplitude of the non-BPS particles due to gravitational exchanges by gluing two 3-point functions (4.19) with the appropriate massless propagators. Summing over graviton and dilaton exchanges and introducing the same notation adopted for the gauge interactions, we obtain
$$\left(_4^{\mathrm{grav}}\right)_{\alpha _1\alpha _4;\alpha _2\alpha _3}=(2\kappa _{10})^2(\zeta _1\zeta _4)(\zeta _2\zeta _3)\frac{M_\mathrm{h}^4(2+\sigma )^2}{t}.$$
(4.20)
Notice that, as expected, neither the 3-point functions (4.19) nor the 4-point function (4.20) depend on the indices $`\alpha _i`$ that span the $`SO(32)`$ spinor representations carried by the non-BPS particles; therefore they can be suppressed. Normalizing each external leg by a factor of $`1/\sqrt{2E}`$ and removing for simplicity the polarization terms, we can obtain from Eq. (4.20) the following gravitational potential energy
$$V_\mathrm{h}^{\mathrm{grav}}(r)=(2\kappa _{10})^2\frac{M_\mathrm{h}^2\gamma }{7\mathrm{\Omega }_8r^7},$$
(4.21)
which in the small velocity limit becomes
$$V_\mathrm{h}^{\mathrm{grav}}(r)=(2\kappa _{10})^2\frac{M_\mathrm{h}^2}{7\mathrm{\Omega }_8r^7}\left(1+\frac{1}{2}v^2\right)+o(v^2).$$
(4.22)
In Eq. (4.21) we recognize Newton’s law for point particles of mass $`M_\mathrm{h}`$ separated by a distance $`r`$ in ten dimensions with the appropriate relativistic correction.
We conclude this section by mentioning that the same results (4.10) and (4.20) can be directly obtained by evaluating a 4-point function of non BPS states on the sphere, or more precisely its “universal” part in the $`t`$-channel which is proportional to $`(\zeta _1\zeta _4)(\zeta _2\zeta _3)`$. In fact, using standard techniques, one can show that this part of the 4-point amplitude is
$$𝒜_4=\frac{4}{\pi ^2}𝒞_0\widehat{𝒩}^4(\zeta _1\zeta _4)(\zeta _2\zeta _3)A(s,t,u;S,T,U),$$
(4.23)
where
$`A(s,t,u;S,T,U)`$ $`=`$ $`ϵ(K)\mathrm{sin}\left[\pi \left({\displaystyle \frac{s\alpha ^{}}{4}}\right)\right]\mathrm{sin}\left[\pi \left({\displaystyle \frac{t\alpha ^{}}{4}}\right)\right]\mathrm{sin}\left[\pi \left({\displaystyle \frac{u\alpha ^{}}{4}}\right)\right]`$
$`\times `$ $`\mathrm{\Gamma }\left(3{\displaystyle \frac{s\alpha ^{}}{4}}\right)\mathrm{\Gamma }\left({\displaystyle \frac{t\alpha ^{}}{4}}\right)\mathrm{\Gamma }\left(3{\displaystyle \frac{u\alpha ^{}}{4}}\right)`$
$`\times `$ $`\mathrm{\Gamma }\left(1{\displaystyle \frac{s\alpha ^{}}{4}}{\displaystyle \frac{S}{2}}\right)\mathrm{\Gamma }\left(1{\displaystyle \frac{t\alpha ^{}}{4}}{\displaystyle \frac{T}{2}}\right)\mathrm{\Gamma }\left(1{\displaystyle \frac{u\alpha ^{}}{4}}{\displaystyle \frac{U}{2}}\right).`$
In this expression $`S`$, $`T`$ and $`U`$ are the Mandelstam variables for the internal momenta which obey $`S+T+U=16`$, and $`ϵ(K)c_{K_3}(K_1+K_2)c_{K_2}(K_1)()^{U/2}`$ is a cocycle factor whose values are $`\pm 1`$ (see for instance Ref. ). Since we are interested only in the contributions due to exchanges of massless states in the $`t`$ channel, we must look for the poles of $`𝒜_4`$ with respect to $`t`$. Inspection of Eq. (4) shows that these occur only for $`T=0`$ and $`T=2`$. The poles for $`T=0`$ correspond to exchanges of gravitons, dilatons and gauge bosons associated to the Cartan generators of $`SO(32)`$, while those for $`T=2`$ correpond to exchanges of the remaining 480 gauge bosons. In the limit $`t0`$, we have
$`𝒜_4|_{T=0}`$ $``$ $`{\displaystyle \frac{16\pi 𝒞_0\widehat{𝒩}^4}{\alpha ^{}}}\left[(4+S/2)(2+\sigma )+(2+\sigma )^2\right]{\displaystyle \frac{1}{t}},`$ (4.25)
$`𝒜_4|_{T=2}`$ $``$ $`{\displaystyle \frac{16\pi 𝒞_0\widehat{𝒩}^4}{\alpha ^{}}}(2+\sigma ){\displaystyle \frac{1}{t}}.`$ (4.26)
We can disentangle the gauge and gravity pieces of Eq. (4.25) by observing that, because of gauge invariance, the gauge part at $`T=0`$ should have the same dynamical dependence as the amplitude (4.26) for $`T=2`$. Hence, we can conclude that the term of Eq. (4.25) linear in $`(2+\sigma )`$ is due to gauge interactions, while the term quadratic in $`(2+\sigma )`$ comes from gravity. Inspection of the coefficients and a little algebra show that these expressions indeed match with Eqs. (4.10) and (4.20), thus providing a strong check on our previous calculations and on their interpretation.
## 5 Conclusions
We now compare the results obtained in the previous sections and discuss their relation in the light of the heterotic/type I duality. For the gravitational interactions, the comparison is quite simple since in both theories we have found a potential energy of the form
$$V^{\mathrm{grav}}(r)=(2\kappa _{10})^2\frac{M^2}{7\mathrm{\Omega }_8r^7}\left(1+\frac{1}{2}v^2\right)+o(v^2)$$
(5.1)
in the non-relativistic limit (see Eqs. (3.22) and (4.22)). The only thing that one has to do to have complete agreement is to change the values of the gravitational coupling constant $`\kappa _{10}`$ and of the mass $`M`$ according to the duality map as we discussed in the introduction. What is nice to observe is that these changes make the two gravitational potential energies agree not only at the static level but also at the first non-trivial order in the velocity $`v`$.
For the gauge interactions the situation is a bit different. Both in the type I theory and in the heterotic string we have found that the gauge potential energy of the stable non-BPS states is in the form of Coulomb’s law (see Eqs. (2.20) and (4.12)). However, the detailed gauge group structure is not the same in the two cases. The reason for this is quite simple. In the heterotic theory one is able to describe the non-BPS particles in a complete way because they are perturbative configurations of the heterotic string, and in particular one can fully specify the polarizations of these states also with respect to the gauge group. This is why the gauge amplitudes involving these particles explicitly depend on the indices of the spinorial representation of $`SO(32)`$ (see Eqs. (4.8) and (4.10)). On the other hand, in the type I theory the non-BPS particles are non-perturbative configurations of the type I string, and thus the description one is able to provide for them using perturbative methods is necessarily incomplete. This fact should not be surprising, because also in the case of the supersymmetric BPS D-branes one is not able to account for their degeneracy (with respect to both the Lorentz group and the gauge group) using open strings with Dirichlet boundary conditions or equivalently boundary states. Indeed with these methods one can compute only the “universal” parts of the interactions involving D-branes.
In Section 2 we have introduced a method to describe the emission of a colored gauge boson from a non-BPS D-particle viewed as a source carrying not the spinorial indices of $`SO(32)`$, but rather those of the bi-fundamental representation formed with the Chan-Paton factors of the boundary changing vertex operators $`𝒱_{09}`$ and $`𝒱_{90}`$ (see Eqs. (2.4) and (2.3)). In this framework, using the various kinds of open strings of type I we have been able to account for the gauge interactions of the non-BPS D-particles, but then the comparison with the heterotic theory is not immediate. In order to do such a comparison, we must “reduce” the heterotic gauge potential energy by taking into account the contribution of all pairs of states compatible with the emission of a gauge boson of definite color. From the group theory point of view, this amounts to transform the spinorial indices of $`\left(V_\mathrm{h}^{\mathrm{gauge}}\right)_{\alpha _1\alpha _4;\alpha _2\alpha _3}`$ given in Eq. (4.12) into those of the bi-fundamental representation. This can be easily done by noting that
$$\left(\mathrm{\Gamma }^A\mathrm{\Gamma }^D\right)^{\alpha \beta }\left(\mathrm{\Gamma }^{BC}\right)_{\beta \alpha }=\mathrm{Tr}\left(\mathrm{\Gamma }^\mathrm{A}\mathrm{\Gamma }^\mathrm{D}\mathrm{\Gamma }^{\mathrm{BC}}\right)=\mathrm{Tr}(\text{ 1l})\left(\delta ^{\mathrm{BD}}\delta ^{\mathrm{AC}}\delta ^{\mathrm{CD}}\delta ^{\mathrm{AB}}\right).$$
(5.2)
Then, using this identity and Eq. (4.12), we obtain the following reduced gauge potential energy for the heterotic non-BPS particles
$`V_\mathrm{h}^{\mathrm{gauge}}(r)`$ $``$ $`\mathrm{Tr}(\text{ 1l})^2\left(\mathrm{\Gamma }^A\mathrm{\Gamma }^D\right)^{\alpha _4\alpha _1}\left(\mathrm{\Gamma }^E\mathrm{\Gamma }^F\right)^{\alpha _3\alpha _2}\left(V_\mathrm{h}^{\mathrm{gauge}}\right)_{\alpha _1\alpha _4;\alpha _2\alpha _3}(r)`$ (5.3)
$`=`$ $`{\displaystyle \frac{g_{\mathrm{YM}}^2}{2}}\left(\delta ^{AE}\delta ^{FD}\delta ^{AF}\delta ^{DE}\right){\displaystyle \frac{1}{7\mathrm{\Omega }_8r^7}}.`$
This expression exactly agrees with the corresponding one for the type I theory given in Eq. (2.20).
We remark that it is not meaningful to perform this reduction directly on the heterotic 3-point function (4.8) and then compare it with the type I amplitude (2.17) describing the emission of a gauge boson from a D-particle. In fact, in the perturbative type I theory the D-particle is an infinitely massive object which acts as a reservoir of momentum, and thus the space-time structure of its amplitudes cannot match with that of the heterotic scattering amplitudes. In other words, the diagram represented in Figure 1 describing the gauge emission from a D-particle of type I must not be considered as a vertex or a 3-point function in the field theory sense, but rather as a 1-point function in some definite background. A similar situation occurs also in the gravitational sector where the boundary state representing the D-particle generates all its 1-point functions, i.e. all its couplings with the closed string states of the bulk. In this sense, what we have done in Section 2 is to find the 1-point function of the non-BPS D-particle with the massless states of the other sector of the bulk, namely the open strings of type I. It would be nice to extend these results to all states of this open string sector.
In conclusion, in this paper we have described how to compute the gauge and gravitational potential energies of the non-BPS D-particles of type I and shown that these agree with the corresponding ones computed for the dual heterotic states provided that one uses the known duality and renormalization effects. Our results thus provide a dynamical test of the heterotic/type I duality at the non-BPS level.
Acknowledgments
We would like to thank C. Bachas, M. Billò, P. Di Vecchia, M. Frau, B. Pioline, C. Schweigert and R. Russo for several useful discussions. We especially thank B. Pioline for valuable discussions and remarks.
Appendix A
In this appendix we describe the gauge emission from a D-string of type I, following the scheme we proposed in Section 2.
As a consequence of the BPS condition, two parallel BPS D$`p`$-branes of type I or type II do not exert any force on each other. In the type II theories, the interaction between two branes is mediated only by the exchange of closed strings and the vanishing of the force is easily seen using boundary states. Indeed one finds that
$$Dp|P|Dp=0$$
(A.1)
at the leading order in the string coupling constant. In the limit of large distance between the branes, when only massless closed string states are exchanged, this means that the attraction due to gravitons and dilatons is compensated by the repulsion due to the $`p+1`$ R-R form under which the D$`p`$-branes are charged. In the case of two parallel D-strings the interaction (A.1) is globally invariant under the world-sheet parity $`\mathrm{\Omega }`$ so that it vanishes also in type I theory. However, in this theory one has to consider also the exchange of open 9-9 strings which are present in the bulk and whose first contribution – associated to a disk with two boundary components on the D-strings and two on the D9 branes – appears at the next-to-leading order in the string coupling constant. For the no force condition to be true, this disk amplitude has thus to vanish identically. At first sight, this seems striking since the D-string is charged under the gauge potential (in fact it carries the $`SO(32)`$ spinorial representation). However, one must recall that, being an extended object, the D-string cannot be minimally coupled to the gauge potential and thus the naive conclusion does not apply. The aim of this appendix is to evaluate the gauge emission from a D-string of type I using the same methods applied in the case of the D-particle (but without the technicalities due to the boundary fermion $`\eta `$), and then to compute its contribution to the diffusion process between two D-strings.
We first briefly discuss the spectrum of the 1-9 open strings stretching between a D-string and a D9 brane. As for the 0-9 strings, also here the NS sector is massive and does not represent any degeneracy of the D-string; thus we do not consider it. In the R sector, instead, there are massless states. Since the world-sheet fermions $`\psi ^0`$, $`\psi ^1`$ have zero modes, the massless R ground state is a GSO even (chiral) spinor of $`SO(1,1)`$. The corresponding vertex operator reads
$$𝒱_{91}^{(1/2)}=\lambda ^AS^+\mathrm{\Delta }^{}\mathrm{e}^{\varphi /2}$$
(A.2)
where $`S^+`$ is a positive chirality spin field of conformal dimension 1/8, and $`\mathrm{\Delta }^{}`$ is a boundary changing operator for the eight space directions transverse to the D-string of conformal dimension $`1/2`$. The vertex operator for the massless R states of the 1-9 sector is obtained by acting with $`\mathrm{\Omega }`$ on $`𝒱_{91}`$; thus
$$𝒱_{19}^{(1/2)}=\mathrm{\Omega }𝒱_{91}^{(1/2)}=\lambda _{}^{t}{}_{}{}^{A}S^+\mathrm{\Delta }^{}\mathrm{e}^{\varphi /2}.$$
(A.3)
We now evaluate the coupling of the D-string with the gauge bosons of type I. As in the case of the D-particle, this is determined by the amplitude on a disk with one boundary component on the D-string and one component on one of the 32 D9-branes from which a gauge boson is emitted, and is represented by the 1-point function of the gauge boson in the vacuum representing the D-string, i.e.
$$^{\mathrm{gauge}}=\mathrm{\Omega }_{19}|V_{\mathrm{gauge}}^{(1)}(1)|\mathrm{\Omega }_{91}$$
(A.4)
where
$$|\mathrm{\Omega }_{19}=\underset{z0}{lim}𝒱_{19}^{(1/2)}(z)|0\mathrm{and}\mathrm{\Omega }_{19}|=\underset{z\mathrm{}}{lim}0|𝒱_{19}^{(1/2)}(z).$$
(A.5)
After introducing the appropriate normalization factors, we can express $`^{\mathrm{gauge}}`$ as
$$^{\mathrm{gauge}}=𝒞_{\mathrm{disk}}𝒩_\mathrm{R}^2𝒩_{\mathrm{NS}}𝑑\mu (z_i)c𝒱_{19}^{(1/2)}(z_1)c𝒱_{\mathrm{gauge}}^{(1)}(z_2)c𝒱_{91}^{(1/2)}(z_3).$$
(A.6)
For the same reasons discussed in the case of the D-particle, also here there is no emission of gauge fields with polarizations along the directions transverse to the D-string; thus we have emissions only in the two longitudinal directions $`\mu =0,1`$, but these can occur with arbitrary transverse momentum. To evaluate (A.6) we need the following basic correlators
$`\mathrm{\Delta }^{}(z_1)\mathrm{e}^{\mathrm{i}kX(z_2)}\mathrm{\Delta }^{}(z_3)`$ $`=`$ $`z_{12}^{\alpha ^{}k^2}z_{13}^{\alpha ^{}k^21}z_{23}^{\alpha ^{}k^2},`$
$`S^+(z_1)\psi ^\mu (z_2)S^+(z_3)`$ $`=`$ $`(\mathrm{\Gamma }^\mu C^1)^{++}z_{12}^{1/2}z_{13}^{1/4}z_{23}^{1/2}`$ (A.7)
for $`\mu =0,1`$. Inserting them into Eq. (A.6), we obtain
$`^{\mathrm{gauge}}`$ $`=`$ $`𝒞_{\mathrm{disk}}𝒩_\mathrm{R}^2𝒩_{\mathrm{NS}}\mathrm{Tr}(\lambda _{}^{t}{}_{}{}^{A}\mathrm{\Lambda }^{BC}\lambda ^D)A_\mu (\mathrm{\Gamma }^\mu C^1)^{++}{\displaystyle 𝑑\mu (z_i)\left(\frac{z_{13}}{z_{12}z_{23}}\right)^{\alpha ^{}k^2}}`$ (A.8)
$`=`$ $`\mathrm{i}g_{\mathrm{YM}}(\delta ^{AB}\delta ^{CD}\delta ^{AC}\delta ^{BD})(A_0A_1)`$
where we have used that $`(\mathrm{\Gamma }^0C^1)^{++}=(\mathrm{\Gamma }^1C^1)^{++}=1`$. As anticipated, this is not a minimal gauge coupling because the D-string is an extended object.
Now, using this coupling and the propagator (2.18) for the massless gauge bosons, we can obtain the gauge potential energy $`V_\mathrm{I}^{\mathrm{gauge}}`$ between two D-strings given by $`V_\mathrm{I}^{\mathrm{gauge}}^{\mathrm{gauge}}𝒫^{\mathrm{gauge}}`$. Inserting the explicit values, we find
$`V_\mathrm{I}^{\mathrm{gauge}}=g_{\mathrm{YM}}^2\left(\delta _{AE}\delta _{DF}\delta _{AF}\delta _{DE}\right){\displaystyle \frac{(1+1)}{q^2}}=0.`$ (A.9)
The vanishing of this contribution confirms our previous statements about the no force condition.
Appendix B
In this appendix we present the definitions of the various normalization factors that are needed for the calculations presented in Sections 2, 3 and 4.
The gravitational coupling constant $`\kappa _{10}`$, the normalization $`𝒞_0`$ of the closed-string tree-level diagrams (sphere diagrams) and the normalization factor $`\widehat{𝒩}`$ of the closed string vertex operators have the same expressions for both the $`SO(32)`$ heterotic string and the type I theory, and are given by
$`\kappa _{10}`$ $`=`$ $`8\pi ^{7/2}\alpha ^{\mathrm{\hspace{0.17em}2}}g,`$ (B.1)
$`𝒞_0`$ $`=`$ $`(2\pi )^4\alpha ^5g^2,`$ (B.2)
$`\widehat{𝒩}`$ $`=`$ $`8\pi ^{5/2}\alpha ^{\mathrm{\hspace{0.17em}2}}g,`$ (B.3)
where $`g`$ is the string coupling constant ($`g_\mathrm{h}`$ for the heterotic string and $`g_\mathrm{I}`$ for the type I theory).
In the heterotic string, the gauge coupling constant $`g_{\mathrm{YM}}`$ is related to $`\kappa _{10}`$ as follows
$$g_{\mathrm{YM}}^2=\frac{4}{\alpha ^{}}\kappa _{10}^2=2^8\pi ^7\alpha ^{\mathrm{\hspace{0.17em}3}}g_\mathrm{h}^2,$$
(B.4)
while in the type I theory the relation between $`g_{\mathrm{YM}}`$ and $`\kappa _{10}`$ is
$$g_{\mathrm{YM}}^2=\frac{2^{3/2}}{\alpha ^{}g_\mathrm{I}}\kappa _{10}^2=2^{15/2}\pi ^7\alpha ^{\mathrm{\hspace{0.17em}3}}g_\mathrm{I}.$$
(B.5)
In type I string theory one must consider also diagrams involving open strings. The normalization of the disk diagrams is
$$𝒞_{\mathrm{disk}}=g_{\mathrm{YM}}^2(2\alpha ^{})^2,$$
(B.6)
while the normalization factors of the open string vertex operators in the NS and R sectors are
$`𝒩_{\mathrm{NS}}`$ $`=`$ $`g_{\mathrm{YM}}\sqrt{2\alpha ^{}}`$
$`𝒩_\mathrm{R}`$ $`=`$ $`g_{\mathrm{YM}}(2\alpha ^{})^{3/4}.`$ (B.7)
We now list the expression of the various Chan-Paton factors that were used in Section 2. The factor $`\mathrm{\Lambda }^{AB}`$ carried by the gauge boson has indices in the adjoint representation of $`SO(32)`$. Its matrix elements explicitly read
$$(\mathrm{\Lambda }^{AB})_{CD}=\mathrm{i}\left(\delta _C^A\delta _D^B\delta _D^A\delta _C^B\right).$$
(B.8)
The Chan-Paton factor $`\lambda ^A`$ associated to 0-9 strings is a column vector with an index in the fundamental representation of $`SO(32)`$. It reads
$$(\lambda ^A)_B=\delta _B^A.$$
(B.9)
With these expressions it is easy to see that
$`\mathrm{Tr}\left(\mathrm{\Lambda }^{AB}\mathrm{\Lambda }^{CD}\right)`$ $`=`$ $`2\left(\delta ^{AC}\delta ^{BD}\delta ^{AD}\delta ^{BC}\right),`$
$`\mathrm{Tr}\left(\lambda _{}^{t}{}_{}{}^{A}\mathrm{\Lambda }^{BC}\lambda ^D\right)`$ $`=`$ $`\mathrm{i}\left(\delta ^{AB}\delta ^{CD}\delta ^{AC}\delta ^{BD}\right).`$ (B.10)
Appendix C
In this appendix we write the vertex operators of the non-BPS heterotic states (4.1) and (4.2), of the gauge bosons of $`SO(32)`$ given in Eqs. (4.3) and (4.4), and of the bosonic states (4.13) of the graviton multiplet.
In the $`(1)`$ superghost picture, the vertices associated to the non-BPS states (4.1) and (4.2) are
$`𝒱_𝒜^{(1)}(z,\overline{z})`$ $`=`$ $`\zeta _\mu 𝒜^\mu (z)\mathrm{e}^{\mathrm{i}kX(z,\overline{z})}C_K\mathrm{e}^{\mathrm{i}\frac{K}{\sqrt{2\alpha ^{}}}\stackrel{~}{X}(\overline{z})},`$ (C.1)
$`𝒱_{}^{(1)}(z,\overline{z})`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\alpha ^{}}}}\zeta _{\mu \nu }^{\mu \nu }(z)\mathrm{e}^{\mathrm{i}kX(z,\overline{z})}C_K\mathrm{e}^{\mathrm{i}\frac{K}{\sqrt{2\alpha ^{}}}\stackrel{~}{X}(\overline{z})},`$ (C.2)
$`𝒱_𝒞^{(1)}(z,\overline{z})`$ $`=`$ $`{\displaystyle \frac{1}{3!}}\zeta _{\mu \nu \rho }𝒞^{\mu \nu \rho }(z)\mathrm{e}^{\mathrm{i}kX(z,\overline{z})}C_K\mathrm{e}^{\mathrm{i}\frac{K}{\sqrt{2\alpha ^{}}}\stackrel{~}{X}(\overline{z})},`$ (C.3)
where
$`𝒜^\mu (z)`$ $`=`$ $`\psi ^\mu (z)\mathrm{e}^{\varphi (z)},`$
$`^{\mu \nu }(z)`$ $`=`$ $`\psi ^\mu (z)x^\nu (z)\mathrm{e}^{\varphi (z)},`$ (C.4)
$`𝒞^{\mu \nu \rho }(z)`$ $`=`$ $`\psi ^\mu (z)\psi ^\nu (z)\psi ^\rho (z)\mathrm{e}^{\varphi (z)}.`$
In Eq. (C.2) the polarization tensor $`\zeta ^{\mu \nu }`$ is symmetric or antisymmetric depending on whether one considers a state in the 44 or in the 84 representation of the Lorentz group. Moreover, in all vertex operators we have introduced suitable cocycle factors $`C_K`$, which depend only on the internal momenta and satisfy
$$C_K(P)C_K^{}(P)=C_{K+K^{}}(P).$$
(C.5)
Applying the picture-changing operator to $`𝒱^{(1)}`$ we can obtain the vertices $`𝒱^{(0)}`$ in the $`(0)`$ superghost picture. They are given by the same expressions (C.1) – (C.3) with $`𝒜^\mu `$, $`^{\mu \nu }`$ and $`𝒞^{\mu \nu \rho }`$ replaced respectively by
$`\widehat{𝒜}^\mu (z)`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}}{\sqrt{2\alpha ^{}}}}\left[^2X^\mu (z)\mathrm{i}\alpha ^{}(k\psi )\psi ^\mu (z)\right],`$
$`\widehat{}^{\mu \nu }(z)`$ $`=`$ $`\mathrm{i}\sqrt{2\alpha ^{}}\left[\psi ^\mu (z)\psi ^\nu (z)(\mathrm{i}/2)(k\psi )\psi ^\mu (z)X^\nu (z)+{\displaystyle \frac{1}{2\alpha ^{}}}X^\mu (z)X^\nu (z)\right],`$
$`\widehat{𝒞}^{\mu \nu \rho }(z)`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}}{\sqrt{2\alpha ^{}}}}[\mathrm{i}\alpha ^{}(k\psi )\psi ^\mu (z)\psi ^\nu (z)\psi ^\rho (z)+X^\mu (z)\psi ^\nu (z)\psi ^\rho (z)`$ (C.6)
$`\psi ^\mu (z)X^\nu (z)\psi ^\rho (z)+\psi ^\mu (z)\psi ^\nu (z)X^\mu (z)].`$
The vertex operators for the gauge bosons (4.3) associated to the 16 Cartan generators of $`SO(32)`$ in the $`(1)`$ superghost picture are
$$𝒱_{\mathrm{gauge}}^{(1)}(z,\overline{z})=\frac{\mathrm{i}}{\sqrt{2\alpha ^{}}}A_\mu \psi ^\mu (z)\mathrm{e}^{\varphi (z)}\mathrm{e}^{\mathrm{i}qX(z,\overline{z})}\overline{}\stackrel{~}{X}^I(\overline{z}),$$
(C.7)
while those for the gauge bosons (4.4) associated to the 480 remaining generators are
$$𝒱_{\mathrm{gauge}}^{(1)}(z,\overline{z})=A_\mu \psi ^\mu (z)\mathrm{e}^{\varphi (z)}\mathrm{e}^{\mathrm{i}qX(z,\overline{z})}C_Q\mathrm{e}^{\mathrm{i}\frac{Q}{\sqrt{2\alpha ^{}}}\stackrel{~}{X}(\overline{z})}.$$
(C.8)
In the $`(0)`$ superghost picture these vertices become respectively
$$𝒱_{\mathrm{gauge}}^{(0)}(z,\overline{z})=\frac{1}{2\alpha ^{}}A_\mu \left[X^\mu (z)\mathrm{i}\alpha ^{}(q\psi )\psi ^\mu (z)\right]\mathrm{e}^{\mathrm{i}qX(z,\overline{z})}\overline{}\stackrel{~}{X}^I(\overline{z}),$$
(C.9)
and
$$𝒱_{\mathrm{gauge}}^{(0)}(z,\overline{z})=\frac{\mathrm{i}}{\sqrt{2\alpha ^{}}}A_\mu \left[X^\mu (z)\mathrm{i}\alpha ^{}(q\psi )\psi ^\mu (z)\right]\mathrm{e}^{\mathrm{i}qX(z,\overline{z})}C_Q\mathrm{e}^{\mathrm{i}\frac{Q}{\sqrt{2\alpha ^{}}}\stackrel{~}{X}(\overline{z})}.$$
(C.10)
Finally, the vertices for the bosonic states (4.13) of the graviton multiplet are
$$𝒱_{\mathrm{grav}}^{(1)}(z,\overline{z})=\frac{\mathrm{i}}{\sqrt{2\alpha ^{}}}ϵ_{\mu \nu }\psi ^\mu (z)\mathrm{e}^{\varphi (z)}\overline{}\stackrel{~}{X}^\nu (\overline{z})\mathrm{e}^{\mathrm{i}qX(z,\overline{z})}$$
(C.11)
in the $`(1)`$ superghost picture, and
$$𝒱_{\mathrm{grav}}^{(0)}(z,\overline{z})=\frac{1}{2\alpha ^{}}ϵ_{\mu \nu }\left[X^\mu (z)\mathrm{i}\alpha ^{}(q\psi )\psi ^\mu (z)\right]\overline{}\stackrel{~}{X}^\nu (\overline{z})\mathrm{e}^{\mathrm{i}qX(z,\overline{z})}$$
(C.12)
in the $`(0)`$ superghost picture. |
warning/0001/astro-ph0001291.html | ar5iv | text | # CADIS deep star counts: Galactic structure and the stellar luminosity function
## 1 Introduction
The stellar structure of the Galaxy has been studied by many groups using a variety of methods, such as metallicity-age determinations or studies of the kinematics of nearby halo and disk stars (for a review see Norris 1998). Other discussions can be found in Fuchs & Jahreiß (1998), or Fuhrmann (1998). The method predominantly used to study Galactic strucure is the method of starcounts. This provides a measurement of the density distribution of the stellar component of the Galaxy Gilmore & Reid (1983); Reid et al. (1996, 1997); Gould et al. (1998).
In the so called standard model (Bahcall & Soneira 1980) the vertical structure of the disk follows an exponential law ($`\rho e^{z/h_z}`$) with scaleheight $`h_z`$. As Gilmore & Reid (1983) showed, the data can be fitted much better by a superposition of two exponentials. Whether or not this deviation from the single exponential is due to a distinct population of stars is not clear, although this is suggested by the different kinematics and lower metallicities Freeman (1992). We will call the deviation ”thick disk”, regardless of its unknown physical origin.
Although the existence of the “thick disk” component seems well established now Norris (1998), the proper values for the scaleheights and lengths are not exactly determined yet, since in most studies the distance of observable stars has been limited to a few kpc above the Galactic plane.
Another topic of debate has been whether the stellar luminosity function (hereafter SLF) of main sequence stars in the disk declines beyond $`M_V=12^{\mathrm{mag}}`$. Most photometric studies find a down-turn of the slope of the SLF at $`M_V=12^{\mathrm{mag}}`$ Stobie et al. (1989); Kroupa (1995); Gould et al. (1998), whereas local observations of stars in the solar neighbourhood indicate a flat continuation to fainter magnitudes Wielen et al. (1983); Jahreiß & Wielen 1997a . Therefore, it is worthwhile to derive the SLF from the CADIS starcounts.
The paper is structured as follows: the Calar Alto Deep Imaging Survey is outlined in Sect. 2. Sect. 3 describes the preparation of the stellar data, Sect. 4 deals with the density distribution of the stars, Sect. 5 with the SLF and its implications for the mass function. Sect. 6 gives a brief summary and an outlook on future prospects.
## 2 The Calar Alto Deep Imaging Survey
The Calar Alto Deep Imaging Survey combines an emission line survey carried out with an imaging Fabry-Perot interferometer with a deep multicolour survey using three broad-band optical to NIR filters and up to eighteen medium-band filters when fully completed. The combination of different observing strategies facilitates not only the detection of emission line objects but also to derive photometric spectra of all objects in the fields without performing time consuming slit spectroscopy. Details of the survey and its calibration will be given in Meisenheimer et al. (in preparation).
All observations were performed on Calar Alto, Spain, in the optical wavelength region with the focal reducers CAFOS (Calar Alto Faint Object Spectrograph) at the 2.2 m telescope and MOSCA (Multi Object Spectrograph for Calar Alto) at the 3.5 m telescope, and with the Omega Prime camera for the NIR observations.
As a byproduct of the survey we obtain a lot of multi-color data about faint stars in the Galaxy. Although for the object classification (see below) exposures in two broadband filters and seven medium-band filters (in the case of the 9 h field eight medium-band filters) are used, the present analysis of the stellar component of CADIS is based only on exposures in three filters, $`R_C`$ (central wavelength/width $`\lambda _c/\mathrm{\Delta }\lambda =649\mathrm{nm}/170`$ nm), $`B_C`$ ($`\lambda _c=461\mathrm{nm}/100`$ nm) and $`I_{815}`$ ($`\lambda _c=815\mathrm{nm}/32`$ nm). Exposure times converted to the 2.2 m telescope are given in Table 1.
The nine CADIS fields measure $`1/30\mathit{}\mathrm{°}`$ each and are located at high Galactic latitude to avoid dust absorption and reddening. In all fields the total flux on the IRAS 100 $`\mu `$m maps is less than 2 MJy/sr which corresponds to $`E_{BV}<0.07`$, so we do not have to apply any color corrections. A second selection criterium for the fields was that there should be no star brighter than $`16^{mag}`$ in the CADIS $`R`$ band. In fact the brightest star in the two fields under consideration has an $`R`$ magnitude of $`15.42^{mag}`$.
### 2.1 Object detection and classification
Objects are identified on each of the deep images (superposition of 5 to 15 individual exposures) using the Source Extractor software SExtractor Bertin & Arnouts (1996), and the resulting lists merged into a master catalogue.
Photometry is done using the program Evaluate, which has been developed by Meisenheimer & Röser (1986). Variations in seeing in between individual exposures are taken into account, in order to get accurate colors. For photometric calibration we use a system of ”tertiary” standard stars in the CADIS fields, which are calibrated with secondary standard stars Oke (1990); Walsh (1995) in photometric nights.
From the locus of the stars in the 8-dimensional color space we conclude that the relative calibration between each pair of wavebands is better than 3 % for all objects with $`R=22`$.
Since one of the major goals of the survey is the classification of every object found in all CADIS fields ($`\mathrm{80\hspace{0.17em}000}`$ to $`\mathrm{100\hspace{0.17em}000}`$ in total), a classification scheme was developed which is based on template spectral energy distributions (see Wolf 1998). The observed colors of every object are compared with a color library of known objects, whose colors are obtained from synthetic photometry performed on our CADIS filterset. The input library for stellar spectra was the Gunn & Stryker (1983) catalogue. For each object the probability to belong to a certain object class (stars – quasars – galaxies) is computed.
Objects classified as stars have stellar colors with a likelihood of more than 75%, and images the profile of which does not deviate significantely from that of well defined stars.
Details about the performance and reliability of the classification are given in Wolf et al. 1999, and Wolf (in preparation).
With the current filter set and exposure times the classification is reliable down to a limit of $`R23^{mag}`$. This was checked by spectroscopic follow-up observations of 245 arbitrarily chosen objects (55 stars, 153 galaxies and 20 quasars) with $`R<23^{mag}`$. One galaxy has been classified as a star by it’s colors, two quasars as galaxies and two galaxies as quasars. The star counts are not significantly affected by the misclassifications down to $`R=23^{mag}`$. Thus we restrict our present analysis to stars with $`R23^{mag}`$.
## 3 Preparation of the stellar data
### 3.1 Photometric parallaxes
To derive distances of the stars from the distance modulus $`m_RM_R`$, it is essential to know the absolute magnitudes of the stars, $`M_R`$.
In principle absolute magnitudes of main sequence stars can be obtained from a color-magnitude diagramm. This main sequence approximation is valid for all stars in our sample since we can be sure that it is free from contamination of any non-main sequence stars, as the faint magnitude intervall we observe ($`16R23`$) does not allow the detection of a giant star.
We took a mean $`M_{V_J}`$ versus $`(BV)_J`$ relation from Lang (1992). A complication arises, because there is no $`V`$ filter included in our filter set. Thus we have to convert $`M_V`$ and $`(BV)_J`$ into our filter system, in order to derive the absolute magnitudes from mean main sequence fit in $`M_R`$ versus $`(br)`$, where the CADIS color is defined by:
$`(br)=2.5\mathrm{log}{\displaystyle \frac{F_R^\gamma }{F_B^\gamma }}.`$ (1)
Here $`F_k^\gamma `$ is the flux outside the atmosphere in units of Photons $`\mathrm{m}^2\mathrm{s}^1\mathrm{nm}^1`$ in the CADIS filter $`k`$.
In order to convert $`(br)`$ into the Johnson $`(BV)_J`$ we performed synthetic photometry on the Gunn-Stryker stars (Wolf, priv. comm.).
$`br`$ can be calibrated to the Johnson-Cousins system (the CADIS $`R_C`$ is very close to the Cousins $`R`$) by using Vega as a zero point:
$`(BR)_C`$ $`=`$ $`(br)+2.5\mathrm{log}{\displaystyle \frac{F_{\mathrm{Vega}}^\gamma (\lambda =440\mathrm{n}\mathrm{m})}{F_{\mathrm{Vega}}^\gamma (\lambda =648\mathrm{n}\mathrm{m})}}`$ (2)
$`=`$ $`(br)+0.725.`$
The absolute magnitude $`M_{R_C}`$ is then given by
$`M_{R_C}`$ $`=`$ $`M_{V_J}(V_JR_C),\mathrm{with}`$ (3)
$`(V_JR_C)`$ $`=`$ $`(B_CR_C)(BV)_J,`$
where we assume $`B_J=B_C`$ (see Huang et al., in preparation).
With the above conversions, the main sequence ($`M_R`$ vs $`(br)`$) can be approximated by a fourth order polynomial in the range $`1(br)1.8`$:
$`M_R`$ $`=`$ $`c_0+c_1(br)+c_2(br)^2`$ (4)
$`+`$ $`c_3(br)^3+c_4(br)^4,`$
the parameters of which are:
$`c_0=4.01236`$ $`c_1=4.12575`$ $`c_2=1.89076`$ $`c_3=0.762053`$ $`c_4=0.341384`$ ,
as shown in Fig. 1.
One further complication arises due to the fact that the mean main sequence relation is valid strictly only for stars with solar metallicities, whereas our sample may contain stars spread over a wide range of different metalicities. For the blue stars ($`br<0.7`$), which are almost all halo stars (see Fig. 4), and therefore supposed to be metal-poor, we use the same main sequence relation, but shifted towards fainter magnitudes by $`M_R=0.75^{mag}`$ (see Fig. 1). This value is the mean deviation from the mean main sequence defined by the CNS 4 stars Jahreiß & Wielen 1997b of a subsample of 10 halo stars, for which absolute R magnitudes were available (Jahreiß, priv. com.).<sup>1</sup><sup>1</sup>1Note that this artificial separation may well lead to wrong absolute magnitudes for individual stars (e.g. a disk star with $`r<1`$ kpc but $`br<0.7`$), but should be correct on average.
The spread in a two color diagram $`(br)`$ versus $`(ri)`$ (that is the CADIS color between $`R_C`$ and $`I_{815}`$ analog to Eq. (1)) becomes significant at $`(br)1.0`$, see Fig. 2. The maximal photometric error for the very faint stars is $`0.15^{mag}`$. Here metallicity effects will distort the relation between the measured $`(br)`$ colors and the spectral type (temperature) and thus lead to wrong absolute magnitudes, so we have to correct for metallicity in order to avoid errors in the photometric parallaxes.
The $`R_C`$ filter is strongly affected by metallicity effects like absorption bands of TiO<sub>2</sub> and VO molecules in the stars’ atmosphere, whereas the $`B_C`$ and the medium-band filter $`I_{815}`$ (the wavelength of which was chosen in order to avoid absorption bands in cool stars) are not. So in a first approximation we can assume the ”isophotes” of varying metallicity in a $`(br)`$ versus $`(ri)`$ two color diagram to be straight lines with a slope of $`1`$, along of which we project the measured colors with $`(br)1.0`$ onto the mean main sequence track which in the interval $`1.0(br)1.8`$ is defined by
$`(ri)`$ $`=`$ $`0.39(br)_{\mathrm{corr}}^40.36(br)_{\mathrm{corr}}^3`$ (5)
$`+`$ $`0.09(br)+0.06.`$
This projection implies that stars with $`(br)_{\mathrm{corr}}1.8`$ cannot exist in Fig. 3, which shows the spatial distribution of metallicity corrected $`(br)_{\mathrm{corr}}`$ colors (the limit is indicated by the dashed–dotted line).
Both the upper and lower magnitude limits lead to selection effects which have to be taken into account.
As one can see in Fig. 3 there is an bimodality of the observed color distribution. The accumulation of red stars in the upper left consists mainly of disk stars, and is separated by a void from the blue stars, which predominantly belong to the halo. Thus a crude disk-halo separation can be drawn by a color cut – we take stars with $`(br)_{\mathrm{corr}}<0.7`$ to be halo, stars with $`(br)_{\mathrm{corr}}>0.7`$ to be disk stars. In the following we will make use of this color cut to derive the distribution of the disk stars separately. In a second step, the distribution as a whole will be analysed. In Fig. 3 the cut is denoted by a dotted horizontal line.
In the two fields we have analysed so far we find 95 halo and 178 disk stars, that is a factor of two more disk stars than predicted by the standard model Bahcall & Soneira (1980). This surplus of disk stars was already noted by Reid & Majewski (1993). Fig. 4 shows the distribution of the $`(br)_{\mathrm{corr}}`$ in the two fields under consideration.
## 4 Density distribution of the stars
### 4.1 Completeness correction
As expected, the detection limit of absolute magnitudes (colors) is distance dependent (see Fig. 3).
We use the two-dimensional distribution of stars in the $`(br)_{\mathrm{corr}}`$ vs $`\mathrm{log}r`$ diagram (Fig. 3) to correct for this incompleteness in the following way: First, we divide the distance in logarithmic bins of 0.2 as indicated in Fig. 3 and count the stars up to the the upper color limit $`(br)^{lim}`$ (this is the distance dependent color (luminosity) limit, up to which stars can be detected. The metal-poor halo stars are intrinsically fainter (see paragraph 3)), thus the limits are shifted accordingly).
The nearest bins ($`0.8\mathrm{log}r0.2`$) are assumed to be complete. For the incomplete bins we multiply iteratively with a factor given by the ratio of complete to incomplete number counts in the previous bin, where the limit for the uncorrected counts is defined by the bin currently under examination ($`j`$):
$`N_j^{\mathrm{corr}}=N_j{\displaystyle \underset{i=1}{\overset{j1}{}}}(1+{\displaystyle \frac{N_i^{^{\prime \prime }}}{N_i^{^{}}}}),`$ (6)
where $`N_j`$ is the number of stars in bin $`j`$, $`N^{^{}}`$ is the number of stars in the previous bin ($`j1`$), up to the limit given by the bin $`j`$, and $`N^{^{\prime \prime }}`$ is the number of stars from that limit up to the limit given by bin $`j1`$, see also appendix.
With the poissonian errors $`\sigma _N=\sqrt{N}`$, $`\sigma _N^{^{}}=\sqrt{N^{^{}}}`$, and $`\sigma _{N^{^{\prime \prime }}}=\sqrt{N^{^{\prime \prime }}}`$ the error of the corrected number counts becomes:
$`\sigma _{N_j^{\mathrm{corr}}}^2`$ $`=`$ $`\sigma _{N_j}^2{\displaystyle \underset{i=1}{\overset{j1}{}}}\left(1+{\displaystyle \frac{N_i^{^{\prime \prime }}}{N_i^{^{}}}}\right)^2`$ (7)
$`+`$ $`{\displaystyle \underset{i=1}{\overset{j1}{}}}\left[\left({\displaystyle \frac{1}{N_i^{^{}}}}\right)^2\sigma _{N_i^{^{\prime \prime }}}^2+\left({\displaystyle \frac{N_i^{^{\prime \prime }}}{N_i^{^{}}}}\right)^2\sigma _{N_i^{^{}}}^2\right]`$
$``$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{m=1}{mj}}{\overset{j1}{}}}\left(1+{\displaystyle \frac{N_i^{^{\prime \prime }}}{N_i^{^{}}}}\right)`$
For a detailed deduction see appendix.
The completeness correction is done for each field separately.
With the corrected number counts the density in the logarithmic spaced volume bins ($`V_j=\frac{1}{3}\omega (r_{j+1}^3r_j^3)`$) can than be calculated according to
$`\rho _j={\displaystyle \frac{N_j^{\mathrm{corr}}}{V_j}}={\displaystyle \frac{cN_j}{V_j}},`$ (8)
For every logarithmic distance bin we use the mean height $`z`$ above the Galactic plane $`<z_j>=\mathrm{sin}b<r>`$, where $`<\mathrm{log}r>=\mathrm{log}r_j+(\mathrm{log}r_{j+1}\mathrm{log}r_i)/2=\mathrm{log}r_i+0.1<r>=1.259r_i`$.
### 4.2 Vertical density distribution
We first study the density distribution of the disk stars by taking only stars in the corresponding color interval into account. Although the color-cut at $`(br)_{\mathrm{corr}}=0.7`$ is a rather crude separation between disk and halo, we gain a clearer insight into the disk distribution since the contamination by halo stars is suppressed considerably.
As the nearest stars in our fields have still distances of about 200 pc the normalization at $`z=0`$ has to be established by other means: We take stars from the CNS4 Jahreiß & Wielen 1997b , which are located in a sphere with radius 20 pc around the sun.
The stars in our normalization sample are selected from the CNS4 by their absolute visual magnitudes, according to the distribution of absolute magnitudes of the CADIS disk stars ($`6.5M_v14.5`$).
Fig. 5 shows the resulting density distribution of the disk stars in the two CADIS fields. The solid line represents a fit with a superposition of two exponentials, the dotted line is the fit for the thin disk component (the first seven data points). Obviously a single exponential is not a good description.
It was suggested to fit the thin disk with a secans hyperbolicus – the exponential is unphysical in that sense, that it is not continuously differentiable at $`z=0`$. A squared secans hyperbolicus (which represents a self-gravitating isothermal disk) can be proved not to fit the data very well – indicating that the stellar disk is in no way isothermal (the velocity dispersion depends on the spectral type).
The fits for the three functions under consideration:
$`\rho _{exp}(z)`$ $`=`$ $`n_1\mathrm{exp}(z/h_1)+n_2\mathrm{exp}(z/h_2)`$ (9)
$`\rho _{sech}(z)`$ $`=`$ $`n_3\mathrm{sech}(z/z_0)+n_4\mathrm{exp}(z/h_3)`$ (10)
$`\rho _{sech^2}(z)`$ $`=`$ $`n_5\mathrm{sech}^2(z/\stackrel{~}{z_0})+n_6\mathrm{exp}(z/h_4).`$ (11)
are shown in Fig. 6, the corresponding parameters are given in Table 2.
As the secans hyperbolicus is a sum of two exponentials
$`\mathrm{sech}(z/z_0)={\displaystyle \frac{2}{\mathrm{exp}(z/z_0)+\mathrm{exp}(z/z_0)}},`$ (12)
$`z_0`$ is not really a scaleheight, but has to be compared to $`h_1`$ by multiplying it with arcsech$`\frac{1}{\mathrm{e}}1.65745`$: $`h_1^{}=z_01.65745`$ (cf. Tab. 2).
The errors of the scaleheights and corresponding parameters $`z_0`$ and $`\stackrel{~}{z_0}`$ are estimated by changing their values until $`\chi ^2`$ increases by 1. We find that at $`z=0`$ the contribution of the thick disk component to the entire disk is $`2\pm 4\%`$.
Our values for the scaleheights lie within the range given by different authors in the literature. The proposed values for the exponential scaleheights of the thin disk range between 200 pc Kent et al. 1991 and 325 pc Bahcall & Soneira (1980). The parameters $`\stackrel{~}{z_0}`$ and $`h_4`$ we found in the 16 h field are equal within the errors to the values found by Gould et al (1997).
### 4.3 Density distribution in the halo
Fig. 7 shows the density distribution of all stars in the two CADIS fields. The dotted line is the secans hyperbolicus + exponential fit for the thin and thick disk components, the dashed line is a deVaucouleurs law. Although the space density corresponding to the $`r^{1/4}`$ law in projection has no simple analytic form Young (1976) it is possible to form expansions about the origin and the point at infinity. The latter expansion is valid for a large range of distances. An analytic approximation is Bahcall (1986):
$`\rho _H(z,b,l)=\rho _0{\displaystyle \frac{\mathrm{exp}\left[10.093\left(\frac{R}{R_{}}\right)^{1/4}+10.093\right]}{\left(\frac{R}{R_{}}\right)^{7/8}}}`$
$``$ $`1.25{\displaystyle \frac{\mathrm{exp}\left[10.093\left(\frac{R}{r_{}}\right)^{1/4}+10.093\right]}{\left(\frac{R}{R_{}}\right)^{6/8}}},R<0.03R_{}`$
$``$ $`\left[10.08669/(R/R_{})^{1/4}\right],R0.03R_{}`$
where $`R=(R_{}^2+z^2+\frac{z^2}{\mathrm{tan}^2b}2R_{}\frac{z}{\mathrm{tan}b}\mathrm{cos}l)^{1/2}`$, $`b`$ and $`l`$ are Galactic latidude and longitude, the distance of the sun from the Galactic center $`R_{}=8`$ kpc, $`R=(x^2+z^2)^{1/2}`$, $`x=(R_{}^2+d^2\mathrm{cos}^2b2R_{}d\mathrm{cos}b\mathrm{cos}l)^{1/2}`$, $`z=d\mathrm{sin}b`$. The dot-dashed is the sum of both halo and disk distribution function.
The density of halo stars extrapolated to $`z=0`$ is compared with the data from the CNS4 with an estimated error of $`10`$%. The stars are discriminated against disk stars by their metallicities and kinematics Fuchs & Jahreiß (1998). They are further selected according to the color cut of the CADIS halo stars ($`(br0.7`$).
Fig. 8 shows the distribution of all stars in the range $`z=0`$ to $`z=10`$ kpc, the fit for a single exponential disk, the deVaucouleurs law, the sum of both, and the sum of the thin plus thick disk fit and the deVaucouleurs law. It is obvious that between $`1.5`$ kpc to 5.0 kpc the thick disk provides the predominant contribution to the overall distribution. Thus, the sum of thin disk and deVaucouleurs halo does not fit the distribution.
If we assume the sun to be 8 kpc away from the Galactic center and the normalisation $`\rho _0`$ to be fixed by the data from the CNS4, all parameters are completely determined. As can be seen in Fig. 7, the corresponding plots fit the distribution of the halo stars up to a distance of about 22 kpc above the Galactic plane.
The de Vaucouleurs law is an empirical description of the density distributions of stars in ellipticals, bulges of spirals and galactic halos de Vaucouleurs & Pence (1978) but it is equally justified to fit other distribution functions to the data: a smoothed power law fits the halo stars equally well Fuchs & Jahreiß (1998); Gould et al. (1998), see e.g. Fig. 9.
When applying the power-law fit, we add a further free parameter which describes the flattening of the halo:
$`\rho _H(z,b,l)=`$
$`\rho _{}\left({\displaystyle \frac{r_c^2+r_{}^2}{r_c^2+r_{}^2+\frac{z^2}{tan^2b}2r_{}\frac{z}{\mathrm{tan}b}\mathrm{cos}l+\frac{z^2}{(c/a)^2}}}\right)^{\alpha /2},`$ (14)
with an arbitrarily chosen core radius of $`r_c=1`$ kpc.
We find an axial ratio $`c/a`$ of $`(0.63\pm 0.07)`$, with an exponent $`\alpha =(2.99\pm 0.13)`$.
## 5 The stellar luminosity function
The knowledge of the distribution function enables us to calculate the SLF for the thin disk stars. We selected the stars with distances less than 1.5 kpc. Beyond this point the contribution by thick disk and halo stars becomes dominant.
Although CADIS does not include a filter which is close to the Johnson-V, we calculate absolute visual magnitudes from the $`R`$ magnitudes to make comparison with literature easier. In the relevant magnitude intervall ($`5M_V14`$, i.e $`1<(br)<1.8`$) there holds a linear relation (see Sect. 3.1):
$`M_{V_J}=1.058M_{R_C}.`$ (15)
We calculate an effective volume for every luminosity bin by integrating the distribution function along the line of sight, $`r`$, where the integration limits are given by the minimum between 1.5 kpc and the distance modulus derived for upper and lower limiting apparent magnitude:
$`V_{\mathrm{max}}^{\mathrm{eff}}=\omega {\displaystyle \underset{R_{\mathrm{min}}}{\overset{R_{\mathrm{max}}}{}}}\nu (r,b)r^2𝑑r,`$ (16)
where
$`R_{\mathrm{min}}`$ $`=`$ $`10^{0.2(16^{\mathrm{mag}}M_R)2.0},`$
$`R_{\mathrm{max}}`$ $`=`$ $`\mathrm{min}(1.5\mathrm{kpc},10^{0.2(23^{\mathrm{mag}}M_\mathrm{R})2.0}).`$
The distribution function
$`\nu (r,b)=\mathrm{exp}(r\mathrm{sin}b/h_1)`$ (17)
is normalised to unity at $`z=0`$.
The weighted mean of the SLFs of the two CADIS fields is shown in Fig. 10, in comparison with the SLF of the stars inside a distance of 20 pc Jahreiß & Wielen 1997a , which is based on HIPPARCOS parallaxes.
As can be seen from Fig. 10, the CADIS SLF is equal within the errors to the local SLF. Since the errors are dominated by Poissonian statistics, the error bars are small at the faint end, and larger at the bright end, complementary to the local SLF.
Thus the weighted mean of CADIS and local SLF which combines the large number of bright stars in the local sample with our superior sampling of faint stars (shown in Fig. 11) can be regarded as the most accurate determination of the SLF. The combined CADIS/local SLF keeps rising with constant slope to its limit at $`M_V=13`$. This is in pronounced contrast to previous determinations of the SLF based on faint star counts Stobie et al. (1989); Kroupa (1995). The discrepancy is demonstrated in Fig. 11 where we compare the combined SLF with the most recent photometric SLF based on HST observations Gould et al. (1998).
At this faint end both incompleteness and uncertainties of the exact location of the main sequence may introduce systematic errors which are hard to quantify.
### 5.1 Implications for the mass function
The stellar mass function (SMF) can in principle be inferred from the SLF by converting the $`M_V`$ magnitudes of the disk stars ($`r1.5`$ kpc) into masses; this requires a mass-luminosity relation, for which we adopted the analytical fit taken from Henry et al. (1993, 1999). Their relation is valid for stellar masses from $`0.08M/M_{}`$ to $`2.0M/M_{}`$. If we count the stars in equally spaced mass intervals of $`0.1`$ and divide the number counts in the $`j_{th}`$ interval by the maximum effective volume (see Eq. (16)), calculated for the corresponding luminosity,
the SMF can be represented by a power law:
$`\mathrm{\Psi }(M/M_{})\left[{\displaystyle \frac{M}{M_{}}}\right]^a.`$ (18)
From the data of the disk stars in the two CADIS fields we derived $`a=1.28\pm _{0.42}^{0.68}`$ for \[$`0.2<M/M_{}<1.1`$\]. Its slope is equal (within the errors) to the value proposed by Henry et al. (1993) ($`a=1`$).
## 6 Summary and future prospects
Although the Calar Alto Deep Imaging Survey is designed as an extragalactic survey, we obtain a substantial set of multicolor-data about faint stars in the Galaxy. With the current filter set and exposure times the classification is reliable down to $`M_R=23`$.
From the $`300`$ stars we identified in two CADIS fields covering a total area of 1/15 , we deduced the density distribution of the stars up to a distance of about 20 kpc above the Galactic plane. The density distribution shows unambiguously the contribution of a thick disk component of the Galaxy with a scaleheight of $`1.3`$ kpc.
There has been discussion whether this ”thick disk” is introduced artificially by the assumption that all stars are on the main sequence Bahcall & Soneira (1984). Our present sample of very faint field stars does rule out this ambiguity: if there were any giant stars included in our data, they would have distances of about 250 kpc, thus their contribution to our sample can be neglected.
The density distribution in the halo, which essentially is corroborated by comparing the local density of CNS4 stars with the density between 5.5 kpc and 15 kpc above the Galactic plane from CADIS, is perfectly fit by a de Vaucouleurs’ law.
Based on our present, limited data set of 72 stars beyond z=5.5 kpc, it is equally well justified to fit an inverse power law with exponent $`\alpha =2.99\pm 0.07`$ to the data, with the axial ratio of the halo as an additional free parameter. The best fit value for the axial ratio is ($`c/a=0.63\pm 0.07`$).
Based on this knowledge of the density distribution, we determined the stellar luminosity function (SLF) for the disk stars. To this end, we confine our sample to stars with distances less than 1.5 kpc, for at this point the contribution of the thick disk and the halo population becomes predominant. Our result is within the errors indistinguishable from the local SLF Jahreiß & Wielen 1997a . Thus we conclude that the weighted average between the local SLF (containing predominantly stars at bright magnitudes ($`M_V<10`$)) and the CADIS SLF (with superior statistics at the faint end ($`M_V>10`$)) can be regarded as the best estimate of the SLF of the disk stars. The SLF continues to rise at least up to $`M_V=13`$, in contrast to most other photometric SLFs, which show a down-turn at $`M_V=12`$.
We regard this as a hint that incompleteness corrections and the closely connected issue of the ”true” sampling volume of these faint star counts might need some critical revision.
The stellar mass function (SMF) of the disk stars was derived by converting the visual magnitudes into masses Henry et al. (1993, 1999). For a power-law SMF $`\mathrm{\Psi }(M/M_{})^a`$ we find a slope $`a=1.28\pm _{0.42}^{0.68}`$. This is (within the errors) consistent with the slope proposed by Henry et al. (1993) ($`a=1`$).
The main purpose of the present paper was to explore to what extent the faint star counts in CADIS can be used for determining the stellar density distribution and the stellar luminosity function of the disk. However, one should keep in mind that our present analysis is based only on about a quarter of of the entire Calar Alto Deep Imaging Survey data. In the entire sample which will contain eight fields of 1/30 each, we expect to identify a total of $`1200`$ stars with $`R23`$, of which 800 are supposed to be disk stars and 400 to be halo stars.
In addition with the complete multicolor data of the survey it should be possible to push the limit for a reliable star classification to $`R=23.5`$ or beyond.
As soon as better template spectra of stars with sub-solar metallicities become available, the currently rough correction for metallicity effect can be replaced by a more accurate metallicity dependent main sequence.
With this much more accurate determination of photometric distances and an at least 4-fold increased statistics it should be possible to adress the following issues concerning the density and luminosity function of the stars in the Galaxy:
* the scalelength $`h_r`$ of the density distribution of stars in the disk. This can be done if one takes the longitude dependent radial decrease of the density into account, e.g. the scaleheight measured in the CADIS 16h field ($`l90\mathrm{°}`$) should be only very slightly affected by the radial density gradient, whereas the 9h field is suffering from a maximal radial decrease in density, so one measures an effective scaleheight $`h_{\mathrm{eff}}`$. Likely the star counts in the 18 h and 23 h field (at Galactic coordinates $`b=30\mathrm{°}`$, $`l=95\mathrm{°}`$, and $`b=43\mathrm{°}`$, $`l=90\mathrm{°}`$, respectively, are only very slightly affected by the radial density gradient, so it should be possible to measure the ”true” scaleheights $`h_1`$ and $`h_2`$ for thin and thick disk component with high precision. From this, the counts in the 9h field and the 10 h field ($`b=53\mathrm{°}`$, $`l=150\mathrm{°}`$), we will deduce the scalelength $`r_i=[1/h_{\mathrm{eff}}1/h]^1`$.
* the axial ratio of a flattened halo and the exponent of the inverse power law with higher precision.
* a bulge – halo separation. By analysing the star counts in all eight fields it will be possible to separate the contribution of the bulge from the halo distribution.
* the slope of the SLF with higher precision. From the entire data set with limiting magnitude $`R23.5`$ we will be able to deduce the SLF down to $`M_V14`$ with good statistics and thus decide whether the features like the Wielen dip Wielen et al. (1983) at $`M_V7`$ and the apparent rise of the SLF beyond $`M_V=12`$ are real.
###### Acknowledgements.
The authors would like to thank H.Jahreiß who identified the subset of stars from the Forth Catalogue of Nearby Stars (CNS4) and thus provided the local normalisation for the disk component.
We are greatly indebted to the anonymous referee who pointed out several points which had not received sufficient attention in the original manuscript. This led to a substantial improvement of the paper.
We also thank the Calar Alto staff for their help and support during many observing runs at the observatory.
## Appendix A Completeness correction
Since the first distance bin is assumed to be complete, we can iteratively correct the number counts in the next bins.
$`N_j`$ is the number of stars in bin $`j`$, $`N^{^{}}`$ is the number of stars in the previous bin ($`j1`$), up to the limit given by the bin $`j`$, and $`N^{^{\prime \prime }}`$ is the number of stars from that limit up to the limit given by bin $`j1`$, as indicated in Fig. 12.
The division into $`N^{^{}}`$ and $`N^{^{\prime \prime }}`$ is done to avoid correlated errors.
The first bin ($`j=1`$) is assumed to be complete ($`N_1`$). The corrected number $`N_2^{\mathrm{corr}}`$ in the second bin is then
$`N_2^{\mathrm{corr}}=N_2\left({\displaystyle \frac{N_1}{N_1^{^{}}}}\right)=N_2\left({\displaystyle \frac{N_1^{^{}}+N_1^{^{\prime \prime }}}{N_1^{^{}}}}\right)=N_2\left(1+{\displaystyle \frac{N_1^{^{\prime \prime }}}{N_1^{^{}}}}\right),`$
and for the third bin
$`N_3^{\mathrm{corr}}`$ $`=`$ $`N_3\left({\displaystyle \frac{N_2^{\mathrm{corr}}}{N_2^{^{}}}}\right)=N_3\left({\displaystyle \frac{N_1}{N_1^{^{}}}}\right)\left({\displaystyle \frac{N_2}{N_2^{^{}}}}\right)`$
$`=`$ $`N_3\left({\displaystyle \frac{N_1^{^{}}+N_1^{^{\prime \prime }}}{N_1^{^{}}}}\right)\left({\displaystyle \frac{N_2^{^{}}+N_2^{^{\prime \prime }}}{N_2^{^{}}}}\right)`$
$`=`$ $`N_3\left(1+{\displaystyle \frac{N_1^{^{\prime \prime }}}{N_1^{^{}}}}\right)\left(1+{\displaystyle \frac{N_2^{^{\prime \prime }}}{N_2^{^{}}}}\right).`$
Thus this can generally be written as
$`N_j^{\mathrm{corr}}=N_j{\displaystyle \underset{i=1}{\overset{j1}{}}}\left(1+{\displaystyle \frac{N_i^{^{\prime \prime }}}{N_i^{^{}}}}\right)`$ (19)
The errors of the number counts $`N_j`$, $`N_{j1}^{^{}}`$ and $`N_{j1}^{^{\prime \prime }}`$ are independent of each other and thus Poissonian: $`\sigma _{N_j}=\sqrt{N_j}`$, $`\sigma _{N_{j1}^{^{}}}=\sqrt{N_{j1}^{^{}}}`$, $`\sigma _{N_{j1}^{^{\prime \prime }}}=\sqrt{N_{j1}^{^{\prime \prime }}}`$. Following Gaussian error propagation the error for the corrected number of stars in the second bin is
$`\sigma _{N_2^{\mathrm{corr}}}^2`$ $`=`$ $`\sigma _{N_2}^2+\left({\displaystyle \frac{N_1^{^{\prime \prime }}}{N_1^{^{}}}}\right)^2\sigma _{N_2}^2`$
$`+`$ $`\left({\displaystyle \frac{N_2}{N_1^{^{}}}}\right)^2\sigma _{N_1^{^{\prime \prime \prime \prime }}}^2+\left({\displaystyle \frac{N_2N_1^{^{\prime \prime }}}{N_1^{{}_{}{}^{\prime \prime }2}}}\right)^2\sigma _{N_1^{^{}}}^2,`$
and for the third bin
$`\sigma _{N_3^{\mathrm{corr}}}^2`$ $`=`$ $`\sigma _{N_3}^2\left(1+{\displaystyle \frac{N_2^{^{\prime \prime }}}{N_2^{^{}}}}\right)^2\left(1+{\displaystyle \frac{N_1^{^{\prime \prime }}}{N_2^{^{}}}}\right)^2+`$
$`+`$ $`N_3^2\left(1+{\displaystyle \frac{N_1^{^{\prime \prime }}}{N_1^{^{}}}}\right)\left[\left({\displaystyle \frac{1}{N_2^{^{}}}}\right)^2\sigma _{N_2^{^{\prime \prime }}}+\left({\displaystyle \frac{N_2^{^{\prime \prime }}}{N_2^{^{}}}}\right)^2\sigma _{N_2^{^{\prime \prime }}}^2\right]`$
$`+`$ $`N_3^2\left(1+{\displaystyle \frac{N_2^{^{\prime \prime }}}{N_2^{^{}}}}\right)\left[\left({\displaystyle \frac{1}{N_1^{^{}}}}\right)^2\sigma _{N_1^{^{\prime \prime }}}+\left({\displaystyle \frac{N_1^{^{\prime \prime }}}{N_1^{^{}}}}\right)^2\sigma _{N_1^{^{\prime \prime }}}^2\right]`$
Thus the error of the completeness corrected counts in the $`j^{th}`$ bin is
$`\sigma _{N_j^{\mathrm{corr}}}^2`$ $`=`$ $`\sigma _{N_j}^2{\displaystyle \underset{i=1}{\overset{j1}{}}}\left(1+{\displaystyle \frac{N_i^{^{\prime \prime }}}{N_i^{^{}}}}\right)^2`$ (20)
$`+`$ $`{\displaystyle \underset{i=1}{\overset{j1}{}}}\left[\left({\displaystyle \frac{1}{N_i^{^{}}}}\right)^2\sigma _{N_i^{^{\prime \prime }}}^2+\left({\displaystyle \frac{N_i^{^{\prime \prime }}}{N_i^{^{}}}}\right)^2\sigma _{N_i^{^{}}}^2\right]`$
$``$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{m=1}{mj}}{\overset{j1}{}}}\left(1+{\displaystyle \frac{N_i^{^{\prime \prime }}}{N_i^{^{}}}}\right).`$ |
warning/0001/cs0001015.html | ar5iv | text | # Multi-Agent Only Knowing
## 1 Introduction
Levesque introduced a notion of “only knowing”, with the goal of capturing certain types of nonmonotonic reasoning. In particular, he hoped to capture the type of reasoning that says “If all I know is that Tweety is a bird, and that birds typically fly, then I can conclude that Tweety flies”.<sup>1</sup><sup>1</sup>1The reader should feel free to substitute “believe” anywhere we say “know”. Indeed, the formal logic that we use, which is based on the modal logic K45, is more typically viewed as a logic of belief rather than knowledge. Levesque’s logic dealt only with the case of a single agent. It is clear that in many applications of such nonmonotonic reasoning, there are several agents in the picture. For example, it may be the case that all Jack knows about Jill is that Jill knows that Tweety is a bird and that birds typically fly. Jack may then want to conclude that Jill knows that Tweety flies.
Recently, each of us independently attempted to extend Levesque’s logic to the multi-agent case. Although there are a number of similarities in the approaches, there are some significant differences. In this paper, we reexamine the notion of only knowing, going back to first principles. In the process, we point out some problems with both of the earlier definitions. This leads us to consider what the properties of only knowing ought to be. We provide an axiom system that captures all our desiderata, and show that it has a semantics that corresponds to it. The axiom system has an added feature of interest: it involves enriching the language with a modal operator for satisfiability, and thus provides an axiomatization for satisfiability in K45. Unfortunately, the semantics corresponding to this axiomatization is not as natural as we might like. It remains an open question whether there is a natural semantics for only knowing that corresponds to this axiomatization.
The rest of this paper is organized as follows. In the next section, we review the basic ideas of Levesque’s logic and provide an alternative semantics. The use of the alternative semantics leads to a simplification of Levesque’s completeness proof. In Section 3, we review Lakemeyer’s approach, which we call the canonical-model approach, and discuss some of its strengths and weaknesses. In Section 4, we go through the same process for Halpern’s approach. In Section 5, we consider our new approach. Much of the discussion in Sections 3, 4, and 5 is carried out in terms of three critical properties of Levesque’s approach which we cull out in Section 2. In particular, we examine to what extent each of the approaches satisfies these properties. In Section 6, we show how the logic can be used, and discuss its relationship to Moore’s autoepistemic logic . Levesque showed that the single-agent version of his logic of only knowing was closely connected to autoepistemic logic. We extend his result to the multi-agent case. We conclude in Section 7 with some discussion of only knowing.
## 2 Levesque’s Logic of Only Knowing
We begin by reconsidering Levesque’s definition. Let $`\mathrm{\Phi }`$ be a set of primitive propositions. Let $`𝒪𝒩(\mathrm{\Phi })`$ be a propositional modal language formed by starting with the primitive propositions in $`\mathrm{\Phi }`$, and closing off under the classical operators $`\neg `$ and $``$ and two modalities, $`L`$ and $`N`$. We omit the $`\mathrm{\Phi }`$ whenever it is clear from context or not relevant to the discussion. We freely use other connectives like $``$, $``$, and $``$ as syntactic abbreviations of the usual kind. In addition, we take $`O\alpha `$ to be an abbreviation for $`L\alpha N\neg \alpha `$. Here $`L\alpha `$ should be read as “the agent knows or believes (at least) $`\alpha `$”, $`N\alpha `$ should be read as “the agent believes at most $`\neg \alpha `$” (so that $`N\neg \alpha `$ is “the agent believes at most $`\alpha `$”) and $`O\alpha `$ should be read as “the agent knows only $`\alpha `$”. We define an objectivex formula to be a propositional formula (i.e., a formula with no modal operators), a subjectivex formula to be a Boolean combination of formulas of the form $`L\phi `$ or $`N\phi `$, and a basicx formula to be a formula which does not mention $`N`$.
Levesque gave semantics to knowing and only knowing using the standard possible-worlds approach. In the single-agent case, we can identify a situationx with a pair $`(W,w)`$, where $`w`$ is a possible world (represented as a truth assignment to the primitive propositions) and $`W`$ consists of a set of possible worlds. Intuitively, $`W`$ is the set of worlds which the agent considers (epistemically) possible, and $`w`$ describes the real world. We do not require that $`wW`$ or that $`W\mathrm{}`$.<sup>2</sup><sup>2</sup>2By requiring that $`W`$ is nonempty, we get the modal logic KD45; by requiring that $`wW`$, we get S5. As usual, we say that the agent knows (at least) $`\alpha `$ if $`\alpha `$ is true in all the worlds that the agent considers possible. Formally, the semantics of the modality $`L`$ and the classical connectives is given as follows.
| $`(W,w)p`$ if $`wp`$ if $`p`$ is a primitive proposition. |
| --- |
| $`(W,w)\neg \alpha `$ if $`(W,w)|\ne \alpha `$. |
| $`(W,w)\alpha \beta `$ if $`(W,w)\alpha `$ or $`(W,w)\beta `$. |
| $`(W,w)L\alpha `$ if $`(W,w^{})\alpha `$ for all $`w^{}W`$. |
Notice that if $`L\alpha `$ holds, then the agent may know more than $`\alpha `$. For example, $`Lp`$ does not preclude $`L(pq)`$ from holding. This is why we should think of $`L\alpha `$ as saying that the agent knows at least $`\alpha `$.
It is well-known that this logic is characterized by the axiom system K45. For convenience, we describe K45 here:
| P. | All instances of axioms of propositional logic. |
| --- | --- |
| K. | $`(L\phi L(\phi \psi ))L\psi `$. |
| 4. | $`L\phi LL\phi `$. |
| 5. | $`\neg L\phi L\neg L\phi `$. |
| R1. | From $`\phi `$ and $`\phi \psi `$ infer $`\psi `$. |
| --- | --- |
| R2. | From $`\phi `$ infer $`L\phi `$. |
The axioms 4 and 5 are called the positive introspection axiom and negative introspection axiom, respectively. They are appropriate for agents that are sufficiently introspective so that they know what they know and do not know.
How do we give precise semantics to $`N`$? That is, when should we say that $`(W,w)N\beta `$? Intuitively, $`N\beta `$ is true if $`\beta `$ is true at all the worlds that the agent does not consider possible. It seems fairly clear from the intuition that we need to evaluate the truth of $`\beta `$ in worlds $`w^{}W`$,<sup>3</sup><sup>3</sup>3Note that, since we defined worlds extensionally as truth assignments, the set of impossible worlds is well-defined and fixed for a given $`W`$. since these are the worlds that the agent considers impossible in $`(W,w)`$. But if $`\beta `$ is a complicated formula involving nested $`L`$ operators, then we cannot simply evaluate the truth of $`\beta `$ at a world $`w^{}`$. We need to have a set of worlds too. In fact, the set of possible worlds we use is still $`W`$. That is, while evaluating the truth of $`\beta `$ in the impossible worlds, the agent keeps the set of worlds he considers possible fixed. Formally, we define
$`(W,w)N\alpha `$ if $`(W,w^{})\alpha `$ for all $`w^{}W`$.
Let us stress three important features of this definition.
* First, as we have already observed, the set of possibilities is kept fixed when we evaluate $`N\alpha `$.
* Second, the set of conceivable worldsx—the union of the set of “possible” worlds considered when evaluating $`L`$ and the set of “impossible” worlds considered when evaluating $`N`$—is fixed, independent of the situation $`(W,w)`$; it is always the set of all truth assignments.
* Finally, for every set of conceivable worlds, there is a model where that set is precisely the set of worlds that the agent considers possible.
Roughly speaking, the first property is what is required for A4, while the second property is what is required for A5. The intent of the third property is to make it possible that any objective formula can be “all you know.” As we shall see in Section 5, in a precise sense, the third property is somewhat stronger than we actually need. All we really need is that for any objective formulas the set of all worlds satisfying these formulas can be considered possible. We return to these properties for guidance when we discuss possible ways of extending Levesque’s semantics to the multi-agent case.
Since $`O\alpha `$ is an abbreviation for $`L\alpha N\neg \alpha `$, we have that
$`(W,w)O\alpha `$ if for all worlds $`w^{}`$, $`w^{}W`$ iff $`(W,w^{})\alpha `$.
As it stands, the semantics has the somewhat odd property that there are situations that agree on all basic beliefs yet disagree on what is only believed. As pointed out by Levesque , the problem is that there are far too many sets of worlds than there are basic belief sets. In order to find a perfect match between the sets of basic beliefs an agent may hold and sets of worlds, Levesque introduces what he calls maximalx sets of worlds. In essence, a maximal set is the largest set in the sense that adding any other world to it would change the agent’s basic beliefs. Furthermore, every set of worlds can be extended to a unique maximal set of worlds. It is well known that in the logic K45, an agent’s beliefs are completely determined by his beliefs about objective formulas (see, for example, for a proof). Thus, we define a maximal set as follows:
###### Definition 2.1
If $`W`$ is a set of worlds, let
$$W^+=\{w|\text{for all objective formulas }\phi \text{, if }(W,w)|=L\phi \text{ then }(W,w)|=\phi \}.$$
$`W`$ is called maximal iff $`W=W^+`$.
Levesque defines validity and satisfiability with respect to maximal sets only. In particular, a formula $`\alpha `$ is valid iff for every maximal set of worlds $`W`$ and every world $`wW`$, we have $`(W,w)|=\alpha `$.
We end this review of Levesque’s logic by presenting (a slight variant of) his proof theory.
| A1. | All instances of axioms of propositional logic. |
| --- | --- |
| A2. | $`L(\alpha \beta )(L\alpha L\beta )`$. |
| A3. | $`N(\alpha \beta )(N\alpha N\beta )`$. |
| A4. | $`\sigma L\sigma N\sigma `$ for every subjective formula $`\sigma `$. |
| A5. | $`N\alpha \neg L\alpha `$ if $`\neg \alpha `$ is a propositionally consistent objective formula. |
| MP. | From $`\alpha `$ and $`\alpha \beta `$ infer $`\beta `$. |
| --- | --- |
| Nec. | From $`\alpha `$ infer $`L\alpha `$ and $`N\alpha `$. |
Axioms A2A4 tell us that that $`L`$ and $`N`$ separately have all the properties of K45-operators. Actually, A4 tells us more; it says that $`L`$ and $`N`$ are mutually introspective, so that, for example, $`L\phi NL\phi `$ is valid. Perhaps the most interesting axiom is A5, which gives only-knowing its desired properties. Its soundness depends on the fact that the union of the set of worlds considered when evaluating $`L`$ and the set of worlds considered when evaluating $`N`$ is the set of all conceivable worlds.<sup>4</sup><sup>4</sup>4Note that, while unusual, the axiom schema A5 is recursive, since consistency of formulas in classical propositional logic is decidable. Hence the axioms themselves are recursive. As noted in , this is a problem in the first-order case, however. In fact, Levesque’s proof theory for the first-order version of his logic was recently shown to be incomplete .
###### Theorem 2.2
If $`\mathrm{\Phi }`$ is infinite, then Levesque’s axiomatization is sound and complete for the language $`𝒪𝒩(\mathrm{\Phi })`$ with respect to Levesque’s semantics.
As we shall see, the assumption that there are infinitely many primitive propositions in $`\mathrm{\Phi }`$ is crucial for Levesque’s completeness result. Extra axioms are required if $`\mathrm{\Phi }`$ is finite. In addition, it is interesting to note that the assumption that $`L`$ and $`N`$ are interpreted with respect to complementary sets of worlds is not forced by the axioms. In particular, for the soundness of Axiom A5, it suffices that the sets considered for $`L`$ and $`N`$ cover all conceivable worlds; they may overlap. The following semantics makes this precise.
Define an extended situationx to be a triple $`(W_L,W_N,w)`$, where $`W_L`$ and $`W_N`$ are sets of worlds (truth assignments) such that $`W_LW_N`$ consists of all truth assignments. Define a new satisfaction relation $`^x`$ that is exactly like Levesque’s except for $`L`$\- and $`N`$-formulas. For them, we have
| $`(W_L,W_N,w)^xL\alpha `$ if $`(W_L,W_N,w^{})^x\alpha `$ for all $`w^{}W_L`$ |
| --- |
| $`(W_L,W_N,w)^xN\alpha `$ if $`(W_L,W_N,w^{})^x\alpha `$ for all $`w^{}W_N`$. |
Note that $`L`$ and $`N`$ are now treated in a completely symmetric way.
###### Theorem 2.3
For all $`\mathrm{\Phi }`$, Levesque’s axiomatization is sound and complete for the language $`𝒪𝒩(\mathrm{\Phi })`$ with respect to $`^x`$.
###### Proof
We omit the soundness proof, which is straightforward. Note that for axiom A5 to be sound it suffices that $`W_L`$ and $`W_N`$ together cover all worlds. In particular, it does not matter whether or not the two sets overlap.
To prove completeness, we use the notion of a maximal consistent set. Given an arbitrary axiom system $`AX`$, we say that a formula $`\phi `$ is consistent with respect to AX if it is not the case that $`AX\neg \phi `$, where, as usual, we use $``$ to denote provability. A finite set of formulas $`\phi _1,\mathrm{},\phi _n`$ is consistent with respect to $`AX`$ if the conjunction $`\phi _1\mathrm{}\phi _n`$ is consistent with respect to $`AX`$. An infinite set of formulas is consistent with respect to $`AX`$ if every finite subset of its formulas is consistent with respect to $`AX`$. Finally, given a set $`F`$ of formulas, a maximal consistent subset of $`F`$ is a subset $`F^{}`$ of $`F`$ which is consistent with respect to $`AX`$ such that any superset of $`F^{}`$ is not consistent with respect to $`AX`$.
In the following, provability, consistency, and maximal consistency all refer to Levesque’s axiom system unless stated otherwise. To prove completeness we show that every consistent formula is satisfiable with respect to $`^x`$, using a standard canonical model construction .
Let $`\mathrm{\Gamma }_0`$ be the set of all maximal consistent sets of formulas in $`𝒪𝒩(\mathrm{\Phi })`$. For $`\theta \mathrm{\Gamma }_0`$, define $`\theta /L=\{\alpha |L\alpha \theta \}`$ and $`\theta /N=\{\alpha |N\alpha \theta \}`$. We then define
* $`\mathrm{\Gamma }_L^\theta =\{\theta ^{}\mathrm{\Gamma }_0|\theta /L\theta ^{}\}`$,
* $`\mathrm{\Gamma }_N^\theta =\{\theta ^{}\mathrm{\Gamma }_0|\theta /N\theta ^{}\}`$.
If we view maximal consistent sets as worlds, then $`\mathrm{\Gamma }_L^\theta `$ and $`\mathrm{\Gamma }_N^\theta `$ represent the worlds accessible from $`\theta `$ for $`L`$ and $`N`$, respectively. The following lemma reflects the fact that $`L`$ and $`N`$ are both fully and mutually introspective (axiom A4).
###### Lemma 2.4
If $`\theta ^{}\mathrm{\Gamma }_L^\theta \mathrm{\Gamma }_N^\theta `$, then $`\mathrm{\Gamma }_L^\theta ^{}=\mathrm{\Gamma }_L^\theta `$ and $`\mathrm{\Gamma }_N^\theta ^{}=\mathrm{\Gamma }_N^\theta `$.
###### Proof
We prove the lemma for $`\theta ^{}\mathrm{\Gamma }_L^\theta `$. The case $`\theta ^{}\mathrm{\Gamma }_N^\theta `$ is completely symmetric. To show that $`\mathrm{\Gamma }_L^\theta =\mathrm{\Gamma }_L^\theta ^{}`$, it clearly suffices to show that $`\theta /L=\theta ^{}/L`$. Let $`\alpha \theta /L`$. Then $`L\alpha \theta `$ and also $`LL\alpha \theta `$ by axiom A4. Thus $`L\alpha \theta ^{}`$ (since $`\theta ^{}\mathrm{\Gamma }_L^\theta `$ implies that $`\theta /L\theta ^{}`$) and, hence, $`\alpha \theta ^{}/L`$.
For the converse, let $`\alpha \theta ^{}/L`$. Thus, $`L\alpha \theta ^{}`$. Assume that $`\alpha \theta /L`$. Then $`\neg L\alpha \theta `$ (since $`\theta `$ is a maximal consistent set) and, therefore, $`L\neg L\alpha \theta `$, from which $`\neg L\alpha \theta ^{}`$ follows, a contradiction.
The proof that $`\mathrm{\Gamma }_N^\theta ^{}=\mathrm{\Gamma }_N^\theta `$ proceeds the same way, that is, we show that $`\theta /N=\theta ^{}/N`$. Let $`\alpha \theta /N`$. Then $`N\alpha \theta `$ and also $`LN\alpha \theta `$ by axiom A4. Hence $`N\alpha \theta ^{}`$, so $`\alpha \theta ^{}/N`$.
For the converse, let $`\alpha \theta ^{}/N`$. Thus, $`N\alpha \theta ^{}`$. Assume that $`\alpha \theta /N`$. Then $`\neg N\alpha \theta `$ and also $`L\neg N\alpha \theta `$, from which $`\neg N\alpha \theta ^{}`$ follows, a contradiction.
In traditional completeness proofs using maximal consistent sets (see, for example, ), a situation is constructed whose worlds consists of all maximal consistent sets. Here, we must be a little more careful.
We say that a maximal consistent set $`\theta `$ contains a truth assignment $`w`$ if for all atomic formulas $`p`$, we have $`wp`$ iff $`p\theta `$. Clearly a maximal consistent set $`\theta `$ contains exactly one world; we denote this world by $`w_\theta `$. For $`\theta \mathrm{\Gamma }_0`$, let $`W_L^\theta =\{w_\theta ^{}|\theta ^{}\mathrm{\Gamma }_L^\theta \}`$ and $`W_N^\theta =\{w_\theta ^{}|\theta ^{}\mathrm{\Gamma }_N^\theta \}`$.
###### Lemma 2.5
1. $`(W_L^\theta ,W_N^\theta ,w_\theta )`$ is an extended situation.
2. For all $`\alpha `$, we have $`\alpha \theta `$ iff $`(W_L^\theta ,W_N^\theta ,w_\theta )^x\alpha `$.
###### Proof
For part (a), to show that $`(W_L^\theta ,W_N^\theta ,w_\theta )`$ is an extended situation, we must show that $`W_L^\theta W_N^\theta `$ consists of all truth assignments. By way of contradiction, suppose there is a truth assignment $`w`$ not in $`W_L^\theta W_N^\theta `$. Let $`F_w=\{p\mathrm{\Phi }|wp\}\{\neg p|p\mathrm{\Phi },w\neg p\}`$. $`F_w\theta /L`$ cannot be consistent, for otherwise there would be some $`\theta ^{}\mathrm{\Gamma }_L^\theta `$ that contains $`F_w`$, which would mean that $`wW_L^\theta `$. Similarly $`F_w\theta /N`$ cannot be consistent. Thus, there must be formulas $`\phi _1,\phi _2,\phi _3,\phi _4`$ such that $`\phi _1`$ and $`\phi _2`$ are both conjunctions of a finite number of formulas in $`F_w`$, $`\phi _3`$ is the conjunction of a finite number of formulas in $`\theta /L`$, and $`\phi _4`$ is the conjunction of a finite number of formulas in $`\theta /N`$, and both $`\phi _1\phi _3`$ and $`\phi _2\phi _4`$ are inconsistent. Thus, we have $`\phi _3\neg \phi _1`$ and $`\phi _4\neg \phi _2`$. Using standard modal reasoning (A2, A3, and Nec), we have $`L\phi _3L\neg \phi _1`$ and $`N\phi _4N\neg \phi _2`$. Since $`L\psi \theta `$ for each conjunct $`\psi `$ of $`\phi _3`$, standard modal reasoning shows that $`L\phi _3\theta `$. Similarly, we have $`N\phi _4\theta `$. Since $`\theta `$ is a maximal consistent set, both $`L\neg \phi _1`$ and $`N\neg \phi _2`$ are in $`\theta `$. Since $`L\neg \phi _1L(\neg \phi _1\neg \phi _2)`$ and $`N\neg \phi _2N(\neg \phi _1\neg \phi _2)`$, it follows that both $`L(\neg \phi _1\neg \phi _2)`$ and $`N(\neg \phi _1\neg \phi _2)`$ are in $`\theta `$. But this contradicts A5, since $`\phi _1\phi _2`$ is a propositionally consistent objective formula.
For part (b), the proof proceeds by induction on the structure of $`\alpha `$. The statement holds trivially for atomic propositions, conjunctions, and negations. In the case of $`L\alpha `$, we proceed by the following chain of equivalences:
| | $`L\alpha \theta `$ |
| --- | --- |
| iff | for all $`\theta ^{}\mathrm{\Gamma }_L^\theta `$, we have $`\alpha \theta ^{}`$ |
| iff | for all $`\theta ^{}\mathrm{\Gamma }_L^\theta `$, we have $`(W_L^\theta ^{},W_N^\theta ^{},w_\theta ^{})^x\alpha `$ (using the induction hypothesis) |
| iff | for all $`w_\theta ^{}W_L^\theta `$, we have $`(W_L^\theta ,W_N^\theta ,w_\theta ^{})^x\alpha `$ (by Lemma 2.4) |
| iff | $`(W_L^\theta ,W_N^\theta ,w_\theta )^xL\alpha `$. |
The case $`N\alpha `$ is completely symmetric.
The completeness result now follows easily. Let $`\alpha `$ be a consistent formula and $`\theta `$ a maximal consistent set of formulas containing $`\alpha `$. By Lemma 2.5, $`(W_L^\theta ,W_N^\theta ,w_\theta )^x\alpha `$.
Levesque considered only maximal sets in his definition of validity. In fact, this restriction has no effect on the notion of validity.
###### Corollary 2.6
A formula $`\alpha 𝒪𝒩(\mathrm{\Phi })`$ is valid iff $`(W,w)\alpha `$ for all situations $`(W,w)`$ (including nonmaximal $`W`$).
###### Proof
If $`\mathrm{\Phi }`$ is finite, it is easy to check that $`W^+=W`$ for all sets $`W`$, so the result is trivially true if $`\mathrm{\Phi }`$ is finite. So suppose $`\mathrm{\Phi }`$ is infinite. Notice that each situation $`(W,w)`$ corresponds to an extended situation $`(W_L,W_N,w)`$, where $`W_L=W`$ and $`W_N`$ is the complement of $`W`$. Let us call such an $`(W_L,W_N,w)`$ an extended complementary situation. Theorems 2.2 and 2.3 together imply the valid formulas obtained when considering all extended situations remain the same when we restrict ourselves to complementary situations with maximal $`W_L`$. The corollary then follows from the fact that the set of all extended situations properly includes the set of all extended complementary situations, which in turn includes the set of all extended complementary situations with maximal $`W_L`$.
As Theorem 2.3 shows, for the $`^x`$ semantics, Levesque’s axioms are sound and complete for all sets $`\mathrm{\Phi }`$ of primitive propositions. On the other hand, as we said earlier, Levesque’s completeness proof (with respect to his semantics) depends crucially on the fact that $`\mathrm{\Phi }`$ is infinite. If $`\mathrm{\Phi }`$ is finite, Levesque’s axioms are still sound with respect to his semantics, but they are no longer complete. For example, if $`\mathrm{\Phi }=\{p\}`$, then $`\neg L\neg pN\neg p`$ would be valid under $``$; this does not follow from the axioms given above. In fact, for each finite set $`\mathrm{\Phi }`$ of primitive propositions, we can find a new axiom scheme that, taken together with the previous axioms, gives a complete axiomatization for $`𝒪𝒩(\mathrm{\Phi })`$ for Levesque’s semantics if $`\mathrm{\Phi }`$ is finite.<sup>5</sup><sup>5</sup>5This was also the situation for the logic considered in . In that paper, a simple axiomatization was provided for the case where $`\mathrm{\Phi }`$ was infinite; for each finite $`\mathrm{\Phi }`$, an extra axiom was needed (that depended on $`\mathrm{\Phi }`$). The new axiom, which subsumes axiom A5, allows us to reduce formulas involving $`N`$ formulas involving only $`L`$.
Note that worlds, which are truth assignments to the primitive propositions $`\mathrm{\Phi }`$, are themselves finite if $`\mathrm{\Phi }`$ is finite. Hence we can identify a world $`w`$ with the conjunction of all literals over $`\mathrm{\Phi }`$ that are true at $`w`$. For example, if $`\mathrm{\Phi }=\{p,q\}`$ and $`w`$ makes $`p`$ true and $`q`$ false, then we identify $`w`$ with $`p\neg q`$. For any objective formula $`\alpha `$, let $`W_{\alpha ,\mathrm{\Phi }}`$ be the set of all worlds (over the primitive propositions $`\mathrm{\Phi }`$) that satisfy $`\alpha `$. The axiom system AX<sub>Φ</sub> is then obtained from Levesque’s system by replacing A5 by the following axiom:
The axiom is easily seen to be sound since it merely expresses that $`N\alpha `$ holds at $`W`$ just in case $`W`$ contains all worlds that satisfy $`\neg \alpha `$. Note that this property depends only on the fact that $`L`$ and $`N`$ are defined with respect to complementary sets of worlds and, hence, also holds in the case of infinite $`\mathrm{\Phi }`$. However, it is only in the finite case that we can express this axiomatically. Completeness is also very easy to establish. Levesque showed that in his system, even without A5, every formula is provably equivalent to one without nested modalities. With A5<sub>Φ</sub>, we then obtain an equivalent formula that does not mention $`N`$. In other words, given a formula consistent with respect to AX<sub>Φ</sub>, a satisfying model can be constructed with the usual technique for $`K45`$ alone.
###### Theorem 2.7
AX<sub>Φ</sub> is sound and complete for the language $`𝒪𝒩(\mathrm{\Phi })`$ with respect to Levesque’s semantics, if $`\mathrm{\Phi }`$ is finite.
## 3 The Canonical-Model Approach
How do we extend our intuitions about only knowing to the multi-agent case? First we extend the language $`𝒪𝒩(\mathrm{\Phi })`$ to the case of many agents. That is, we now consider a language $`𝒪𝒩_n(\mathrm{\Phi })`$, which is just like $`𝒪𝒩`$ except that there are modalities $`L_i`$ and $`N_i`$ for each agent $`i`$, $`1in`$, for some fixed $`n`$. In the remainder of the paper, we omit the $`\mathrm{\Phi }`$, just writing $`𝒪𝒩`$ and $`𝒪𝒩_n`$, since the set of primitive propositions does not play a significant role. By analogy with the single-agent case, we call a formula basic if it does not mention any of the operators $`N_i`$ ($`i=1,\mathrm{},n`$) and $`i`$-subjectivex if it is a Boolean combination of formulas of the form $`L_i\phi `$ and $`N_i\phi `$. What should be the analogue of an objective formula? It clearly is more than just a propositional formula. From agent 1’s point of view, a formula like $`L_2p`$ or even $`L_2L_1p`$ is just as “objective” as a propositional formula. We define a formula to be $`i`$-objectivex if it is a Boolean combination of primitive propositions and formulas of the form $`L_j\phi `$ and $`N_j\phi `$, $`ji`$, where $`\phi `$ is arbitrary. Thus, $`qN_2L_1p`$ is 1-objective, but $`L_1p`$ and $`qL_1p`$ are not. The $`i`$-objective formulas true at a world can be thought of as characterizing what is true apart from agent $`i`$’s subjective knowledge of the world.
The standard model here is to have a Kripke structure with worlds and accessibility relations that describe what worlds the agents consider possible in each world. Formally, a (Kripke) structure or model is a tuple $`M=(W,\pi ,𝒦_1,\mathrm{},𝒦_n)`$, where $`W`$ is a set of worlds,<sup>6</sup><sup>6</sup>6Note that here $`W`$ denotes the set of all worlds of the particular model $`M`$, not just the (epistemically) possible ones as in Levesque’s logic. $`\pi `$ associates with each world a truth assignment to the primitive propositions, and $`𝒦_i`$ is agent $`i`$’s accessibility relation. Given such a Kripke structure $`M`$, let $`𝒦_i^M(w)=\{w^{}:(w,w^{})𝒦_i\}`$.<sup>7</sup><sup>7</sup>7We use the superscript $`M`$ since we shall later need to talk about the $`𝒦_i`$ relations in more than one model at the same time. $`𝒦_i^M(w)`$ is the set of worlds that agent $`i`$ considers possible at $`w`$ in structure $`M`$. As usual, we define
$`(M,w)L_i\alpha `$ if $`(M,w^{})\alpha `$ for all $`w^{}𝒦_i^M(w)`$.
We focus on structures where the accessibility relations are Euclidean and transitive, where a relation $`R`$ on $`W`$ is Euclidean if $`(u,v)R`$ and $`(u,w)R`$ implies that $`(v,w)R`$, and $`R`$ is transitive if $`(u,v)R`$ and $`(v,w)R`$ implies that $`(u,w)R`$. We call such structures K45<sub>n</sub>-structures. It is well known that these assumptions are precisely what is required to get belief to obey the K45 axioms (generalized to $`n`$ agents). We say that a formula consistent with these axioms is K45<sub>n</sub>-consistent. An infinite set of formulas is said to be K45<sub>n</sub>-consistent if the conjunction of the formulas in every one of its finite subsets is K45<sub>n</sub>-consistent.
Now the question is how to define the modal operator $`N_i`$. The problem in the multi-agent case is that we can no longer identify a possible world with a truth assignment. In the single-agent case, knowing the set of truth assignments that the agent considers possible completely determines his knowledge. This is no longer true in the multi-agent case. Somehow we must take the accessibility relations into account. A general semantics for an $`N`$-like operator was first given by Humberstone and later by Ben-David and Gafni . In this approach, the semantics of $`N_i`$ is given as follows:
$$(M,w)N_i\alpha \text{ if }(M,w^{})\alpha \text{ for all }w^{}W𝒦_i^M(w).$$
The problem with this definition is that it misses out on the intuition that when evaluating $`N_i\alpha `$, we keep the set of worlds that agent $`i`$ considers possible fixed. If $`w^{}W𝒦_i^M(w)`$, there is certainly no reason to believe that $`𝒦_i^M(w)=𝒦_i^M(w^{})`$.
One approach to solving this problem is as follows: If $`w`$ and $`w^{}`$ are two worlds in $`M`$, we write $`w_iw^{}`$x if $`𝒦_i^M(w)=𝒦_i^M(w^{})`$, i.e., if $`w`$ and $`w^{}`$ agree on the possible worlds according to agent $`i`$. We then define
$`(M,w)N_i\alpha `$ if $`(M,w^{})\alpha `$ for all $`w^{}`$ such that $`w^{}W𝒦_i^M(w)`$ and $`w_iw^{}`$.
While this definition does capture the first of Levesque’s properties, it does not capture the second. To see the problem, suppose we have only one agent and a structure $`M`$ with only one possible world $`w`$. Suppose that $`(w,w)𝒦_1^M`$ and $`p`$ is true at $`w`$. Then it is easy to see that $`(M,w)L_1pN_1p`$, contradicting axiom A5. The problem is that since the structure has only one world and it is in $`𝒦_1^M(w)`$, there are no worlds in $`W𝒦_1^M(w)`$. Thus, $`N_1p`$ is vacuously true. Intuitively, there just aren’t enough “impossible” worlds in this case; the set of conceivable worlds is not independent of the model. To deal with this problem, we focus attention on one particular model, the canonical model, which intuitively has “enough” worlds. Its worlds consist of all the maximal consistent subsets of basic formulas. (Recall that maximal consistent sets were defined in the proof of Theorem 2.3.) Thus, in some sense, the canonical model has as many worlds as possible.
###### Definition 3.1
The canonical model (for K45<sub>n</sub>) $`M^c=(W^c,\pi ^c,𝒦_1^c,\mathrm{},𝒦_n^c)`$ is defined as follows:
* $`W^c=\{w|w\text{ is a maximal consistent set of basic formulas with respect to K45}\text{n}\}`$,
* for all primitive propositions $`p`$ and $`wW^c`$, $`\pi (w)(p)=`$ true iff $`pw`$,
* $`(w,w^{})𝒦_i^c`$ iff $`w/L_iw^{}`$, where $`w/L_i=\{\alpha |L_i\alpha w\}`$.
Validity in the canonical-model approach is defined with respect to the canonical model only. More precisely, a formula $`\alpha `$ is said to be valid in the canonical-model approach, denoted $`^c\alpha `$, iff $`M^c|=\alpha `$, that is, if for all worlds $`w`$ in the canonical model we have $`(M^c,w)|=\alpha `$.
We clearly cannot use the canonical model in a practical way. It can be shown that it has uncountably many worlds. Each of its worlds is characterized by an infinite set of formulas, so cannot be described easily. Moreover, in general, both $`𝒦_i^c(w)`$ and $`W^c𝒦_i^c(w)`$ are infinite, so we cannot compute whether $`L_i\phi `$ or $`N_i\phi `$ holds at a given world. Thus, our interest in the canonical model is mainly to understand whether it gives reasonable semantics to the $`N_i`$ operator.
We start by arguing that, for an appropriate notion of “possibility” and “conceivability”, this semantics satisfies the first two of Levesque’s properties. What then is a conceivable world? Intuitively, it is an objective state of affairs from agent $`i`$’s point of view, which does not include $`i`$’s beliefs. In the single-agent case, this is simply a truth assignment. In the multi-agent case, things are more complicated, since beliefs of other agents are also part of $`i`$’s objective world. One way of characterizing a state of affairs from $`i`$’s point of view is by the set of $`i`$-objective formulas that are true at a particular world. For technical reasons, in this section we restrict even further to the $`i`$-objective basic formulas—that is, those formulas that do not mention any of the modal operators $`N_j`$, $`j=1,\mathrm{},n`$—that are true. If we assume that the basic formulas determine all the other formulas, which can be shown to be true in the single-agent case, and under this semantics for the multi-agent case, then it is arguably reasonable to restrict to basic formulas. However, as we shall see in Section 4, it is not clear that this restriction is appropriate, although we make it for now.
###### Definition 3.2
Given a situation $`(M,w)`$, let $`\text{obj}_i(M,w)`$x consist of all the $`i`$-objective basic formulas that are true at $`(M,w)`$. let $`\text{Obj}_i(M,w)`$x = $`\{\text{obj}_i(M,w^{})|w^{}𝒦_i^M(w)\}`$, and let $`\text{subj}_i(M,w)`$x = $`\{\text{basic}L_i\alpha |(M,w)|=L_i\alpha \}\{\text{basic}\neg L_i\alpha |(M,w)|\ne L_i\alpha \}`$.
We take $`\text{obj}_i(M,w)`$ to be $`i`$’s state at $`(M,w)`$. Notice that $`\text{obj}_i(M,w)`$ is a maximal consistent set of $`i`$-objective basic formulas. For ease of exposition, we say $`i`$-setx from now on rather than “maximal set of $`i`$-objective basic formulas”. Thus, the set of conceivable states for agent $`i`$ is the set of all $`i`$-sets. Notice that the set of conceivable states is independent of the model. It is easy to show that this is a generalization of the single-agent case, since in the single-agent case the $`i`$-objective basic formulas are just the propositional formulas, and an $`i`$-set can be identified with a truth assignment.
$`\text{Obj}_i(M,w)`$ is the set of $`i`$-sets that agent $`i`$ considers possible in situation $`(M,w)`$; $`\text{subj}_i(M,w)`$ characterizes $`i`$’s basic beliefs in $`(M,w)`$. Notice that if $`\alpha `$ is an $`i`$-objective basic formula, then $`L_i\alpha \text{subj}_i(M,w)`$ iff $`\alpha `$ is in every $`i`$-set in $`\text{Obj}_i(M,w)`$.
With these definitions, we can show that the first two of Levesque’s properties hold in the canonical model. The first property says that at all worlds $`w^{}`$ considered in evaluating a formula of the form $`N_i\phi `$ at a world $`w`$, the set of possible states—that is, the set $`\{\text{obj}_i(M^c,w^{\prime \prime })|w^{\prime \prime }𝒦_i^c(w^{})\}`$—is the same for all $`w^{}𝒦_i^c(w)`$. This is easy to see, since the only worlds $`w^{}`$ we consider are those such that $`𝒦_i^c(w^{})=𝒦_i^c(w)`$. The second property says that the union of the set of states associated with the worlds used in computing $`L_i\phi `$ at $`w`$ and the set of states associated with the worlds used in computing $`N_i\phi `$ at $`w`$ should consist of all conceivable states. To show this, we must show that for every world $`w`$ in the canonical model, the set $`\{\text{obj}_i(M^c,w^{})|w^{}_iw\}`$ consists of all $`i`$-sets.
To prove this, we need two preliminary lemmas.
###### Lemma 3.3
Let $`w`$ and $`w^{}`$ be worlds in $`M^c`$. Then $`w_iw^{}`$ iff agent $`i`$ has the same basic beliefs at $`w`$ and $`w^{}`$, that is, $`\text{subj}_i(M^c,w)=\text{subj}_i(M^c,w^{})`$.
###### Proof
The “only if” direction is immediate because $`w`$ and $`w^{}`$ are assumed to have the same $`𝒦_i`$-accessible worlds. To prove the “if” direction, suppose that $`\text{subj}_i(M^c,w)=\text{subj}_i(M^c,w^{})`$ but $`w_iw^{}`$. Without loss of generality, there is a world $`w^{}𝒦_i^c(w)𝒦_i^c(w^{})`$. By the definition of the canonical model, there must be a basic formula $`L_i\alpha w^{}`$ such that $`\alpha w^{}`$. By assumption, $`L_i\alpha w`$, contradicting the assumption that $`w^{}𝒦_i^c(w)`$.
###### Lemma 3.4
Suppose $`\mathrm{\Gamma }`$ consists only of $`i`$-objective basic formulas, $`\mathrm{\Sigma }`$ consists only of $`i`$-subjective basic formulas, and $`\mathrm{\Gamma }`$ and $`\mathrm{\Sigma }`$ are both K45<sub>n</sub>-consistent. Then $`\mathrm{\Gamma }\mathrm{\Sigma }`$ is K45<sub>n</sub>-consistent.
###### Proof
This follows immediately from part (c) of Proposition 4.2 below.
We can now prove that the set of conceivable states for agent $`i`$ is the same at all worlds of the canonical model. This follows from the following result.
###### Theorem 3.5
Let $`wW^c`$. Then for every $`i`$-set $`\mathrm{\Gamma }`$ there is exactly one world $`w^{}`$ such that $`\text{obj}_i(M^c,w^{})=\mathrm{\Gamma }`$ and $`w_iw^{}`$.
###### Proof
Let $`\mathrm{\Sigma }=\text{subj}_i(M^c,w)`$. Since $`\mathrm{\Gamma }`$ consists of $`i`$-objective basic formulas only, $`\mathrm{\Sigma }`$ consists of $`i`$-subjective formulas, and $`\mathrm{\Gamma }`$ and $`\mathrm{\Sigma }`$ are both K45<sub>n</sub>-consistent, by Lemma 3.4, $`\mathrm{\Gamma }\mathrm{\Sigma }`$ is K45<sub>n</sub>-consistent. Let $`w^{}`$ be a maximal consistent set that contains $`\mathrm{\Gamma }\mathrm{\Sigma }`$. Since $`w`$ and $`w^{}`$ agree on $`\mathrm{\Sigma }`$, $`w_iw^{}`$ by Lemma 3.3. The uniqueness of $`w^{}`$ follows by a simple induction argument.
What about the third property? This says that every subset of $`i`$-sets arises as the set of $`i`$-sets associated with the worlds that $`i`$ considers possible in some situation; that is, for every set $`S`$ of $`i`$-sets, there should be some situation $`(M^c,w)`$ such that $`S=\text{Obj}_i(M^c,w)`$. As we now show, this property does not hold in the canonical model. We do this by showing that the set of $`i`$-sets associated with the worlds considered possible in any situation in the canonical model all have a particular property we call limit closure.<sup>8</sup><sup>8</sup>8This turns out to be closely related to the limit closure property discussed in ; a detailed comparison would take us too far afield here though.
###### Definition 3.6
We say that an $`i`$-set $`\mathrm{\Gamma }`$ is a limit of a set $`S`$ of $`i`$-sets if, for every finite subset $`\mathrm{\Delta }`$ of $`\mathrm{\Gamma }`$, there is a set $`\mathrm{\Gamma }^{}S`$ such that $`\mathrm{\Delta }\mathrm{\Gamma }^{}`$. A set $`S`$ of $`i`$-sets is limit closedx if every limit of $`S`$ is in $`S`$.
###### Lemma 3.7
For every world $`w`$ in $`M^c`$, the set $`\text{Obj}_i(M^c,w)`$ is limit closed.
###### Proof
Let $`w`$ be a world in the canonical model and let $`\mathrm{\Gamma }`$ be an $`i`$-set which is a limit of $`\text{Obj}_i(M^c,w)`$. We want to show that $`\mathrm{\Gamma }\text{Obj}_i(M^c,w)`$. Let $`\mathrm{\Sigma }=\{\text{basic }\beta |L_i\beta w\}`$. We claim that $`\mathrm{\Gamma }\mathrm{\Sigma }`$ is K45<sub>n</sub>-consistent. For suppose not. Then there must be a finite subset $`\mathrm{\Delta }\mathrm{\Gamma }`$ such that $`\mathrm{\Delta }\mathrm{\Sigma }`$ is inconsistent. Since $`\mathrm{\Gamma }`$ is a limit of $`\text{Obj}_i(M^c,w)`$, there must be some world $`w^{}`$ in $`\text{Obj}_i(M^c,w)`$ such that $`(M^c,w^{})\mathrm{\Delta }`$. By construction of the canonical model, since $`w^{}𝒦_i^c(w)`$, we must have that $`(M^c,w^{})\mathrm{\Sigma }`$. Thus $`\mathrm{\Delta }\mathrm{\Sigma }`$ is consistent, contradicting our assumption.
Since $`\mathrm{\Gamma }\mathrm{\Sigma }`$ is consistent, there is a world $`w^{}`$ in the canonical model such that $`(M^c,w^{})\mathrm{\Gamma }\mathrm{\Sigma }`$. Clearly $`\mathrm{\Gamma }=\text{obj}_i(M^c,w^{})`$. Moreover, by construction, we must have $`w^{}𝒦_i^c(w)`$. Thus, $`\mathrm{\Gamma }\text{Obj}_i(M^c,w)`$, as desired.
Since there are clearly sets of $`i`$-sets that are not limit closed, it follows that this semantics does not satisfy the third property. One consequence of this is a result already proved in , which we reprove here, using an approach that will be useful for later results.
###### Proposition 3.8
If $`p\mathrm{\Phi }`$ and $`ij`$, then $`^c\neg O_i\neg O_jp`$.
###### Proof
We proceed by contradiction. Our goal is to show that if $`(M^c,w)O_i\neg O_jp`$, then the set $`\text{Obj}_i(M^c,w)`$ must contain a set that includes $`O_jp`$, for otherwise it would not be limit closed. It follows that there is some world $`v𝒦_i^c(w)`$ such that $`(M^c,v)O_jp`$, contradicting the assumption that $`(M^c,w)L_i\neg O_jp`$.
We first need a definition and a lemma. We say that a basic formula $`\psi `$ is (K45)-independent of a basic formula $`\phi `$ if neither $`_{\mathrm{K45}_n}\phi \psi `$ nor $`_{\mathrm{K45}_n}\phi \neg \psi `$ hold.
###### Lemma 3.9
If $`n2`$ (i.e., there are at least two agents) and $`\phi _1,\mathrm{},\phi _m`$ are consistent basic $`i`$-objective formulas, then there exists a basic $`i`$-objective formula $`\psi `$ of the form $`L_j\psi ^{}`$ which is independent of each of $`\phi _1,\mathrm{},\phi _m`$.
###### Proof
Define the depth of a basic formula $`\phi `$, denoted $`d(\phi )`$, inductively:
* $`d(p)=0`$ for a primitive proposition $`p`$,
* $`d(\phi \psi )=\mathrm{max}(d(\phi ),d(\psi ))`$,
* $`d(\neg \phi )=d(\phi )`$,
* $`d(L_i\phi )=1+d(\phi )`$.
Suppose that $`\phi _1,\mathrm{},\phi _m`$ are $`i`$-objective formulas such that $`\mathrm{max}(d(\phi _1(,\mathrm{},d(\phi _k))=K`$. Let $`p`$ be an arbitrary primitive proposition, and suppose $`ji`$. (Such a $`j`$ exists, since we are assuming $`n2`$.) Let $`\psi `$ be the formula $`(L_jL_i)^{K+1}p`$, where by $`(L_jL_i)^{K+1}`$ we mean $`K+1`$ occurrences of $`L_jL_i`$. Standard model theoretic arguments show that $`\psi `$ is independent of $`\phi _1,\mathrm{},\phi _m`$. Very briefly: By results of , we know that $`\phi _i`$ is satisfiable in a treelike structure of depth at most $`d(\phi _i)`$, for $`i=1,\mathrm{},m`$. It is easy to see that this can be extended to two structures, one of which satisfies $`\psi `$, and the other of which satisfies $`\neg \psi `$. Hence, $`\psi `$ is independent of $`\phi _i`$.
Continuing with the proof of Proposition 3.8, suppose by way of contradiction that $`(M^c,w)O_i\neg O_jp`$. Let $`\overline{w}`$ be a world such that $`(M^c,\overline{w})|=O_jp`$ and let $`\mathrm{\Gamma }=\text{obj}_i(M^c,\overline{w})`$. We claim that $`\mathrm{\Gamma }`$ is a limit of $`\text{Obj}_i(M^c,w)`$. To see this, consider any finite subset $`\mathrm{\Delta }`$ of $`\mathrm{\Gamma }`$. Let $`\psi `$ be an $`i`$-objective basic formula of the form $`L_j\psi ^{}`$ which is independent both of the conjunction of the formulas in $`\mathrm{\Delta }`$ and of $`p`$. The existence of such a formula follows from Lemma 3.9. Let $`\mathrm{\Sigma }=\text{subj}_i(M^c,w)`$. By Lemma 3.4, $`\mathrm{\Sigma }\mathrm{\Delta }\{L_j\psi ^{}\}`$ is consistent. Thus, there is some world $`w^{}W^c`$ such that $`(M^c,w^{})\mathrm{\Sigma }\mathrm{\Delta }\{L_j\psi ^{}\}`$. By Lemma 3.4 again, there is some world $`w^{\prime \prime }`$ satisfying $`p\neg L_j\psi ^{}`$ such that $`w^{\prime \prime }_iw^{}`$. Since $`(M^c,w^{})L_j\psi ^{}`$, we cannot have $`w^{\prime \prime }𝒦_j^c(w^{})`$. It follows that $`(M^c,w^{})\neg N_j\neg p`$, and hence $`(M^c,w^{})\neg O_jp`$. Moreover, since $`w^{}`$ and $`w`$ agree on all $`i`$-subjective formulas, the canonical model construction guarantees that $`w^{}_iw`$. Since $`(M^c,w)O_i\neg O_jp`$ , $`(M^c,w^{})\neg O_jp`$, and $`w^{}_iw`$, we must have that $`w^{}𝒦_i^c(w)`$. Thus, $`\text{obj}_i(M^c,w^{})\text{Obj}_i(M^c,w)`$. Moreover, by construction, $`\mathrm{\Delta }\text{obj}_i(M^c,w^{})`$. Since $`\mathrm{\Delta }`$ was chosen arbitrarily, it follows that $`\mathrm{\Gamma }`$ is a limit of $`\text{Obj}_i(M^c,w)`$. By Lemma 3.7, $`\mathrm{\Gamma }\text{Obj}_i(M^c,w)`$. Thus, there is some world $`v𝒦_i^c(w)`$ such that $`\text{obj}_i(M^c,v)=\mathrm{\Gamma }`$. It is a simple property of the canonical-model approach that two worlds that agree on all basic beliefs of an agent also agree on what the agent only believes. Hence, since $`\overline{w}`$ and $`v`$ agree on $`j`$’s basic beliefs, it follows that $`(M^c,v)|=O_jp`$, contradicting the assumption that $`(M^c,w)|=O_i\neg O_jp`$.
It may seem unreasonable that $`\neg O_i\neg O_jp`$ should be valid in the canonical-model approach. Why should it be impossible for $`i`$ to know only that $`j`$ does not only know $`p`$? After all, $`j`$ can (truthfully) tell $`i`$ that it is not the case that all he ($`j`$) knows is $`p`$. We return to this issue in Sections 4 and 5. For now, we focus on a proof theory for this semantics.
### 3.1 A Proof Theory
We now consider an axiomatization for the language. The following axiomatization is exactly like Levesque’s except that axiom A5 now requires K45<sub>n</sub>-consistency instead of merely propositional consistency. For ease of exposition, we use the same names for the axioms as we did in the single-agent case with a subscript $`n`$ to emphasize that we are looking at the multi-agent version.
| A1<sub>n</sub>. | Axioms of propositional logic. |
| --- | --- |
| A2<sub>n</sub>. | $`L_i(\alpha \beta )(L_i\alpha L_i\beta )`$. |
| A3<sub>n</sub>. | $`N_i(\alpha \beta )(N_i\alpha N_i\beta )`$. |
| A4<sub>n</sub>. | $`\sigma L_i\sigma N_i\sigma `$ if $`\sigma `$ is an $`i`$-subjective formula. |
| A5<sub>n</sub>. | $`N_i\alpha \neg L_i\alpha `$ if $`\neg \alpha `$ is a K45<sub>n</sub>-consistent $`i`$-objective basic formula. |
| MP<sub>n</sub>. | From $`\alpha `$ and $`\alpha \beta `$ infer $`\beta `$. |
| --- | --- |
| Nec<sub>n</sub>. | From $`\alpha `$ infer $`L_i\alpha `$ and $`N_i\alpha `$. |
Notice that A5<sub>n</sub> assumes that $`\alpha `$ ranges only over basic $`i`$-objective formulas. We need this restriction in order to appeal to satisfiability in the existing logic K45<sub>n</sub>.<sup>9</sup><sup>9</sup>9Note that this peculiar axiom schema is recursive since satisfiability in propositional K45<sub>n</sub> is decidable . To get a more general version of A5<sub>n</sub>, that applies to arbitrary formulas, we will need to appeal to consistency within the logic that the axioms are meant to characterize. We return to this issue in Section 5. It is not hard to show that these axioms are sound.
###### Theorem 3.10
For all $`\alpha `$ in $`𝒪𝒩_n`$, if $`\alpha `$ then $`^c\alpha `$.
###### Proof
The proof proceeds by the usual induction on the length of a derivation. Here we show only the soundness of A5<sub>n</sub>. Suppose $`\alpha `$ is a basic $`i`$-objective formula such that $`\neg \alpha `$ is K45<sub>n</sub>-consistent. Thus, there is an $`i`$-set containing $`\neg \alpha `$. By Theorem 3.5, it follows that for each world $`wW^c`$, there is a world $`w^{}_iw`$ such that $`(M^c,w^{})\neg \alpha `$. If $`w^{}𝒦_i^c(w)`$, it follows that $`(M^c,w)\neg L_i\alpha `$. If $`w^{}𝒦_i^c(w)`$, then $`(M^c,w)\neg N_i\alpha `$. Thus, $`(M^c,w)\neg L_i\alpha \neg N_i\alpha `$; equivalently, $`(M^c,w)N_i\alpha \neg L_i\alpha `$. It follows that $`N_i\alpha \neg L_i\alpha `$.
We show in Section 4 that this axiomatization is incomplete. In fact, the formula $`\neg O_i\neg O_jp`$ is not provable. Intuitively, part of the problem here is that A5<sub>n</sub> is restricted to basic formulas. For completeness, we would need an analogue of A5<sub>n</sub> for arbitrary formulas. However, we obtain completeness for a restricted language, which we call $`𝒪𝒩_n^{}`$. Roughly, the restriction amounts to limiting what an agent can only know. In particular, only knowing can only be applied to basic formulas.
###### Definition 3.11
$`𝒪𝒩_n^{}`$ consists of all formulas $`\alpha `$ in $`𝒪𝒩_n`$ such that, in $`\alpha `$, no $`N_j`$ may occur within the scope of an $`N_i`$ or $`L_i`$ for $`ij`$.
For example, $`N_iL_i\neg N_ip`$ and $`N_i(L_jpN_i\neg p)`$ are in $`𝒪𝒩_n^{}`$, and $`N_iN_jp`$ and $`N_iL_jN_ip`$ are not, for distinct $`i`$ and $`j`$. Note that formulas such as $`O_i\neg O_jp`$ cannot be expressed in $`𝒪𝒩_n^{}`$.
To prove completeness for the sublanguage $`𝒪𝒩_n^{}`$, we need four preliminary lemmas. The first describes a normal form for formulas. Having such a normal form greatly simplifies the completeness proof.
###### Lemma 3.12
Every formula $`\alpha `$ in $`𝒪𝒩_n`$ is provably equivalent to a disjunction of formulas of the following form:
$$\begin{array}{c}\sigma L_1\phi _{10}\neg L_1\phi _{11}\mathrm{}\neg L_1\phi _{1m_1}\mathrm{}L_n\phi _{n0}\neg L_n\phi _{n1}\mathrm{}\neg L_n\phi _{nm_n}\hfill \\ N_1\psi _{10}\neg N_1\psi _{11}\mathrm{}\neg N_1\psi _{1k_1}\mathrm{}N_n\psi _{n0}\neg N_n\psi _{n1}\mathrm{}\neg N_n\psi _{nk_n},\hfill \end{array}$$
where $`\sigma `$ is a propositional formula and $`\phi _{ij}`$ and $`\psi _{ij}`$ are all $`i`$-objective formulas. Moreover, if $`\alpha `$ in $`𝒪𝒩_n^{}`$, we can assume that $`\phi _{ij}`$ and $`\psi _{ij}`$ are $`i`$-objective basic formulas.
###### Proof
We proceed by induction on the structure of $`\phi `$. The only nontrivial cases are if $`\phi `$ is of the form $`L_i\phi ^{}`$ or $`N_i\phi ^{}`$. If $`\phi `$ is of the form $`L_i\phi ^{}`$, then, since $`\phi 𝒪𝒩_n^{}`$, $`N_j`$ does not appear in $`\phi ^{}`$ for $`ji`$. We use the inductive hypothesis to get $`\phi ^{}`$ into the normal form described in the lemma. Notice that $`N_j`$ does not appear in the normal form for $`ji`$. We now use the the following equivalences to get $`L_i\phi ^{}`$ into the normal form:
* $`L_i(\psi \psi ^{})(L_i\psi L_i\psi ^{})`$
* $`L_i(\psi L_i\psi ^{})(L_i\psi L_i\psi ^{})`$
* $`L_i(\psi \neg L_i\psi ^{})(L_i\psi \neg L_i\psi ^{})`$
* $`L_iL_i\psi L_i\psi `$
* $`(\neg L_i\text{false}L_i\neg L_i\psi )\neg L_i\psi `$
* $`L_iN_i\psi N_i\psi `$
* $`L_i\neg N_i\psi \neg N_i\psi `$.
The first five of these equivalences are standard K45<sub>n</sub> properties; the last two are instances of axiom A4<sub>n</sub>. Similar arguments work in the case that $`\phi `$ is of the form $`N_i\phi ^{}`$. We leave the straightforward details to the reader.
The next two lemmas give us some basic facts about the satisfiability and validity of formulas in $`𝒪𝒩_n^{}`$.
###### Lemma 3.13
If $`S_j`$ is a set of consistent $`j`$-sets, $`j=1,\mathrm{},n`$, and $`\sigma `$ is a consistent propositional formula, then there is a K45<sub>n</sub> situation $`(M,w)`$ such that $`(M,w)\sigma `$ and $`\text{Obj}_j(M,w)=S_j`$.
###### Proof
This follows immediately from part (b) of Proposition 4.2 below.
###### Lemma 3.14
If $`\phi `$ and $`\psi `$ are $`i`$-objective basic formulas such that $`L_i\phi N_i\psi `$ is consistent, then $`\phi \psi `$ is valid.
###### Proof
Suppose that $`\neg \phi \neg \psi `$ is consistent. Then, by Axiom A5<sub>n</sub>, $`N_i(\phi \psi )\neg L_i(\phi \psi )`$ is provable. It follows that $`N_i\psi \neg L_i\phi `$ is provable, contradicting the consistency of $`L_i\phi N_i\psi `$.
The last lemma gives us a sharper condition on when $`w^{}𝒦_i^c(w)`$. This condition will prove useful in the completeness proof.
###### Lemma 3.15
If $`w,w^{}`$ are worlds in the canonical model such that $`w_iw^{}`$ and $`w^{}𝒦_i^c(w)`$, then there is an $`i`$-objective basic formula $`\phi `$ such that $`L_i\phi w`$ and $`\phi w^{}`$.
###### Proof
By the construction of the canonical model, we know that if $`w^{}𝒦_i^c(w)`$, then there is some basic formula $`\phi `$ such that $`L_i\phi w`$ and $`\phi w^{}`$. From Lemma 3.12, it follows that we can assume without loss of generality that $`\phi `$ is in the normal form described by that lemma. Let $`A`$ consist of all subformulas of $`\phi `$ that are of the form $`L_i\psi `$ and do not appear in the scope of any other modal operator. Let $`\phi _A`$ be $`_{\psi Aw}\psi _{\psi Aw}\neg \psi `$. It is easy to see that $`\phi _Aw`$ (since each of its conjuncts is). Since $`w_iw^{}`$, it follows that $`\phi _Aw^{}`$. Let $`\phi ^{}`$ be the result of replacing each subformula $`L_i\psi `$ of $`\phi `$ that is in $`A`$ by either true or false, depending on whether $`L_i\psi w`$. By construction, $`\phi ^{}`$ is $`i`$-objective. It is easy to see that $`\phi _A(\phi \phi ^{})`$ is provable. It follows that $`\phi ^{}w^{}`$. It also follows that $`L_i\phi _A(L_i\phi L_i\phi ^{})`$ is provable. Since $`\phi _A`$ is an $`i`$-subjective formula, $`\phi _AL_i\phi _A`$ is provable. Hence, $`L_i\phi _Aw`$. Since $`L_i\phi w`$, it follows that $`L_i\phi ^{}w`$. This gives us the desired result.
We are finally ready to prove the completeness result.
###### Theorem 3.16
For all $`\alpha 𝒪𝒩_n^{}`$, if $`^c\alpha `$ then $`\alpha `$.
###### Proof
As usual, it suffices to show that if the formula $`\alpha 𝒪𝒩_n^{}`$ is consistent, then it is satisfiable in the canonical model. Without loss of generality, we can assume that $`\alpha `$ is in the normal form described in Lemma 3.12:
$$\begin{array}{c}\sigma L_1\phi _{10}\neg L_1\phi _{11}\mathrm{}\neg L_1\phi _{1m_1}\mathrm{}L_n\phi _{n0}\neg L_n\phi _{n1}\mathrm{}\neg L_n\phi _{nm_n}\hfill \\ N_1\psi _{10}\neg N_1\psi _{11}\mathrm{}\neg N_1\psi _{1k_1}\mathrm{}N_n\psi _{n0}\neg N_n\psi _{n1}\mathrm{}\neg N_n\psi _{nk_n}.\hfill \end{array}$$
Moreover, since $`\alpha 𝒪𝒩_n^{}`$, we can assume that $`\phi _{ij}`$ and $`\psi _{ij}`$ are $`i`$-objective basic formulas. Let $`A_i`$ consist of all the consistent formulas of the form $`\phi _{i0}\psi _{i0}\neg \phi _{ij}`$ or $`\phi _{i0}\psi _{i0}\neg \psi _{ij}`$, $`j1`$. Let $`\xi _i`$ be a formula that is independent of all the formulas in $`A_i`$; such a formula exists by Lemma 3.9. Let $`S_i`$ consist of all $`i`$-sets containing $`\phi _{i0}(\neg \psi _{i0}(\psi _{i0}\xi _i))`$. By Lemma 3.13, there is a K45<sub>n</sub> structure $`(M,w)`$ such that $`\text{Obj}_i(M,w)=S_i`$, $`i=1,\mathrm{},n`$, and $`(M,w)\sigma `$. Thus, there must be a world $`w^{}`$ in the canonical model such that $`w^{}=\{\text{basic }\phi ^{}|(M,w)\phi ^{}\}`$. We claim that $`(M^c,w^{})\alpha `$.
To see this, let $`\alpha ^{}`$ be the formula $`\sigma L_1\phi _{10}\neg L_1\phi _{11}\mathrm{}\neg L_1\phi _{1m_1}\mathrm{}L_n\phi _{n0}\neg L_n\phi _{n1}\mathrm{}\neg L_n\phi _{nm_n}`$. We first show that $`(M,w)\alpha ^{}`$. By construction, we have that $`(M,w)\sigma `$. Furthermore, by definition, each world $`w^{}𝒦_i^M(w)`$ satisfies $`\phi _{i0}`$, so we have that $`(M,w)L_i\phi _{i0}`$. Since $`L_i\phi _{i0}\neg L_i\phi _{ij}`$ is consistent for each $`j1`$, it must be the case that $`\phi _{i0}\neg \phi _{ij}`$ is consistent. Thus, one of $`\phi _{i0}\neg \psi _{i0}\neg \phi _{ij}`$ or $`\phi _{i0}\psi _{i0}\neg \phi _{ij}`$ is consistent. If the latter is consistent, then by the choice of $`\xi `$, $`\phi _{i0}\psi _{i0}\xi _i\neg \phi _{ij}`$ must be consistent as well. Since $`S_i`$ consists of all $`i`$-sets containing $`\phi _{i0}(\neg \psi _{i0}(\psi _{i0}\xi _i))`$, it follows that there must be an $`i`$-set in $`S_i`$ containing $`\neg \phi _{ij}`$. It follows that $`(M,w)\neg L_i\phi _{ij}`$, for $`j1`$. Thus, we have shown that $`(M,w)\alpha ^{}`$. Since $`(M,w)`$ and $`(M^c,w^{})`$ agree on basic formulas, it follows that $`(M^c,w^{})\alpha ^{}`$.
Next, we show that $`(M^c,w^{})N_1\psi _{10}\mathrm{}N_n\psi _{n0}`$. To this end, suppose that $`w^{}_iw^{}`$ and $`w^{}𝒦_i^c(w^{})`$. By Lemma 3.15, there must be some $`i`$-objective basic formula $`\phi ^{}`$ such that $`L_i\phi ^{}w^{}`$ and $`\neg \phi ^{}w^{}`$. Since $`L_i\phi ^{}w^{}`$, it follows that $`(M,w)L_i\phi ^{}`$, and hence $`\phi ^{}`$ is in every $`i`$-set in $`S_i`$. It follows that $`\text{obj}_i(w^{})S_i`$. Now, one of the following four formulas must be in $`\text{obj}_i(w^{})`$: (1) $`\phi _{i0}\psi _{i0}`$, (2) $`\phi _{i0}\neg \psi _{i0}`$, (3) $`\neg \phi _{i0}\psi _{i0}`$, (4) $`\neg \phi _{i0}\neg \psi _{i0}`$. Since $`L_i\phi _{i0}N_i\psi _{i0}`$ is consistent, it cannot be (4), by Lemma 3.14. It cannot be (2), for otherwise $`w^{}`$ would be in $`S_i`$. Thus, it must be (1) or (3), so $`\psi _{i0}\text{obj}_i(w^{})`$. Since this is true for all $`w^{}`$ such that $`w^{}_iw^{}`$ and $`w^{}𝒦_i^c(w^{})`$, it follows that $`(M^c,w^{})N_i\psi _{i0}`$, for $`i=1,\mathrm{},n`$.
Finally, we must show that $`(M^{},w^{})\neg N_1\psi _{ij}`$, for $`i=1,\mathrm{},n`$ and $`j=1,\mathrm{},k_i`$. Clearly $`\psi _{i0}\neg \psi _{ij}`$ is consistent, for otherwise $`N_i\psi _{i0}\neg N_i\psi _{ij}`$ would be inconsistent. Thus, at least one of (1) $`\psi _{i0}\neg \psi _{ij}\neg \phi _{i0}`$ or (2) $`\psi _{i0}\neg \psi _{ij}\phi _{i0}`$ is consistent. In case (2), by choice of $`\xi _i`$, the formula $`\psi _{i0}\neg \psi _{ij}\phi _{i0}\neg \xi _i`$ is consistent. Let $`\beta `$ be $`\psi _{i0}\neg \psi _{ij}\neg \phi _{i0}`$ if it is consistent, and $`\psi _{i0}\neg \psi _{ij}\phi _{i0}\neg \xi _i`$ otherwise. By construction, $`\beta `$ is consistent. By Lemma 3.4, there is a situation $`(M^{},v)`$ such that $`\text{subj}_i(M^{},v)=\text{subj}_i(M^{},w^{})`$ and $`(M^{},v)\beta `$. There is a world $`w^{\prime \prime }`$ in the canonical model which agrees with $`(M^{},v)`$ on the basic formulas. By construction, we have $`w^{\prime \prime }_iw^{}`$. Moreover, since $`(M^c,w)L_i(\phi _{i0}(\neg \psi _{i0}(\psi _{i0}\xi )))`$, it follows that $`(M^c,w^{})L_i\neg \beta `$. Since $`(M^c,w^{})\beta `$, we have that $`w^{}𝒦_i^c(w^{})`$. Moreover, since $`(M^c,w^{})\neg \psi _{ij}`$, it follows that $`(M^c,w^{})\neg N_i\psi _{ij}`$, as desired. This completes the proof.
### 3.2 Discussion
As we have shown, the canonical-model semantics for $`N_i`$ has some attractive features, in particular when restricted to the language $`𝒪𝒩_n^{}`$. It is for this sublanguage that we have a nice proof-theoretic characterization. There is some evidence, however, that the semantics may not have the behavior we desire when we move beyond $`𝒪𝒩_n^{}`$. For one thing, the formula $`\neg O_i\neg O_jp`$ is valid in the canonical model: it is impossible that all $`i`$ knows is that it is not the case that all $`j`$ knows is $`p`$. While it is certainly consistent for $`\neg O_i\neg O_jp`$ to hold, it seems reasonable to have a semantics that allows $`O_i\neg O_jp`$ to be hold as well. As we have seen, the validity of $`\neg O_i\neg O_jp`$ follows from the fact that the canonical-model semantics does not have the third property of Levesque’s semantics in the single-agent case: not all subsets of conceivable states are possible. In the next section, we discuss a different approach to giving semantics to only knowing—essentially that taken in —that has all three of Levesque’s properties, at least as long as we continue to represent an agent’s objective state of affairs using basic formulas only. The approach agrees with the canonical-model approach on formulas in $`𝒪𝒩_n^{}`$, but makes $`O_i\neg O_jp`$ satisfiable. Unfortunately, as we shall see, it too suffers from problems.
## 4 The $`i`$-Set Approach
In the $`i`$-set approach, we maintain the intuition that the set of conceivable states for each agent $`i`$ can be identified with the set of $`i`$-sets. We no longer restrict attention to the canonical model though; we consider all Kripke structures.
We define a new semantics $`^{}`$ as follows: all the clauses of $`^{}`$ are identical to the corresponding clauses for $``$, except that for $`N_i`$. In this case, we have
$$\begin{array}{c}(M,w)^{}N_i\phi \text{ iff }(M^{},w^{})^{}\phi \text{ for all situations }(M^{},w^{})\text{ such that}\hfill \\ \text{Obj}_i(M,w)=\text{Obj}_i(M^{},w^{})\text{ and }\text{obj}_i(M^{},w^{})\text{Obj}_i(M,w)\text{.}\hfill \end{array}$$
Notice that $``$ and $`^{}`$ agree for basic formulas; in general, as we shall see, they differ. We remark that this definition is equivalent to the one given in , except that there, rather than $`i`$-sets, $`i`$-objective trees were considered. We did not want to go through the overhead of introducing $`i`$-objective trees here, since it follows from results in that $`i`$-sets are equivalent to $`i`$-objective trees: every $`i`$-set uniquely determines an $`i`$-objective tree and vice versa.
How well does this approach fare in terms of computing the truth of a formula at a given world. Since we now allow arbitrary structures, not just the canonical model, it seems that for computational reasons, it seems we ought to focus on finite structures. However, when it comes to formulas of the form $`N_i\alpha `$, where $`\alpha `$ is $`i`$-objective, satisfiable formula, it can be shown that $`(M,w)^{}\neg N_i\alpha `$. There are simply too few sets in $`\text{Obj}_i(M,w)`$ if $`M`$ is finite to have $`\neg \alpha `$ hold in all “impossible” situations. This means that to really model interesting situations with only knowing, we must use infinite models, and we are back to the problems discussed in the case of the canonical model. Thus, again we focus on whether this approach gives reasonable semantics to the $`N_i`$ operator.
Notice that to decide if $`N_i\phi `$ holds in $`(M,w)`$, we consider all situations that agree with $`(M,w)`$ on the set of possible states, hence this semantics satisfies the first of the three properties we isolated in the single-agent case. It is also clear that the $`i`$-sets considered in evaluating the truth of $`N_i\phi `$ are precisely those not considered in evaluating the truth of $`L_i\phi `$; hence we satisfy the second property. Finally, as we now show, for every set $`S`$ of $`i`$-sets, there is a situation $`(M,w)`$ such that $`\text{Obj}_i(M,w)=S`$. In fact, we prove an even stronger result.
###### Definition 4.1
Let $`\text{obȷ}_i^+(M,w)`$x consist of all $`i`$-objective formulas (not necessarily just $`i`$-objective basic formulas) true at $`(M,w)`$ (with respect to $`^{}`$) and let $`\text{Obȷ}_i^+(M,w)`$x = $`\{\text{obȷ}_i^+(M,w^{})|w^{}𝒦_i^M(w)\}`$.
###### Proposition 4.2
Let $`\mathrm{\Gamma }`$ be a satisfiable set of $`i`$-objective formulas, let $`S_i`$ be a set of maximal satisfiable sets of $`i`$-objective formulas, $`i=1,\mathrm{},n`$, let $`\mathrm{\Sigma }`$ be a satisfiable set of $`i`$-subjective formulas, and let $`\sigma `$ be a satisfiable propositional formula. Then
1. there exists a situation $`(M_1,w_1)`$ such that $`\mathrm{\Gamma }\text{obȷ}_i^+(M_1,w_1)`$ and $`S_i=\text{Obȷ}_i^+(M_1,w_1)`$.
2. there exists a situation $`(M_2,w_2)`$ such that $`(M_2,w_2)\sigma `$ and $`\text{Obȷ}_j^+(M_2,w_2)=S_j`$, $`j=1,\mathrm{},n`$.
3. there exists a situation $`(M_3,w_3)`$ such that $`(M_3,w_3)\mathrm{\Gamma }\mathrm{\Sigma }`$.
###### Proof
For part (a), we first show that, given an arbitrary situation $`(M,w)`$, we can construct a situation $`(M^{},w^{})`$ such that $`\text{obȷ}_i^+(M^{},w^{})=\text{obȷ}_i^+(M,w)`$ and there are no worlds $`i`$-accessible from $`w^{}`$. The idea is to have $`M^{}`$ be the result of adding $`w^{}`$ to the worlds in $`M`$, where $`w^{}`$ is just like $`w`$ except that it has no $`i`$-accessible worlds and $`w^{}`$ is not accessible from any world. More formally, if $`M=(W,\pi ,𝒦_1,\mathrm{},𝒦_n)`$, we take $`M^{}=(W^{},\pi ^{},𝒦_1^{},\mathrm{},𝒦_n^{})`$, where $`W^{}=W\{w^{}\}`$, $`\pi ^{}(w^{})=\pi (w^{})`$ for $`w^{}W`$, $`\pi ^{}(w^{})=\pi (w)`$, $`𝒦_j^{}=𝒦_j\{(w^{},w^{})|w^{}𝒦_j(w)\}`$ for $`ji`$, and $`𝒦_i^{}=𝒦_i`$. It is easy to see that $`𝒦_j^{}`$ is Euclidean and transitive. By construction, there are no worlds $`i`$-accessible from $`w^{}`$ and $`(w^{},w^{})𝒦_j^{}`$ for all $`w^{}`$ and all $`j`$. Moreover, if $`\psi `$ is an $`i`$-objective formula, we have $`(M^{},w^{})^{}\psi `$ iff $`(M,w)^{}\psi `$, since for $`ji`$, we have $`𝒦_j(w^{})=𝒦_j(w)`$. In particular, this means that $`\text{obȷ}_i^+(M^{},w^{})=\text{obȷ}_i^+(M,w)`$.
For each $`\mathrm{\Delta }S^{}=S_i\{\mathrm{\Gamma }\}`$, there is a situation $`(M^\mathrm{\Delta },w^\mathrm{\Delta })^{}\mathrm{\Delta }`$, where $`M^\mathrm{\Delta }=(W^\mathrm{\Delta },\pi ^\mathrm{\Delta },𝒦_1^\mathrm{\Delta },\mathrm{},𝒦_n^\mathrm{\Delta })`$. By the argument above, we can assume without loss of generality that there are no worlds $`i`$-accessible from $`w^\mathrm{\Delta }`$ and $`w^\mathrm{\Delta }`$ is not accessible from any world. We define $`M_1=(W,\pi ,𝒦_1,\mathrm{},𝒦_n)`$ by taking $`W`$ to be the union of all the worlds in $`W^\mathrm{\Delta }`$, $`\mathrm{\Delta }S^{}`$. (We can assume without loss of generality that these are disjoint sets of worlds.) We define $`\pi `$ so that $`\pi |_{W^\mathrm{\Delta }}=\pi ^\mathrm{\Delta }`$. We define $`𝒦_j=_{\mathrm{\Delta }S^{}}𝒦_j^\mathrm{\Delta }`$ for $`ji`$, and $`𝒦_i`$ to be the least transitive, Euclidean set containing $`_{\mathrm{\Delta }S^{}}𝒦_i^\mathrm{\Delta }\{(w^\mathrm{\Gamma },w^\mathrm{\Delta })|\mathrm{\Delta }S\}`$. It is easy to check that $`\text{obȷ}_i^+(M_1,w^\mathrm{\Delta })=\text{obȷ}_i^+(M^\mathrm{\Delta },w^\mathrm{\Delta })`$ (although this depends on the fact that $`w_\mathrm{\Delta }`$ is not $`j`$-accessible from any world for $`ji`$). Thus, $`\text{Obȷ}_i^+(M_1,w^\mathrm{\Gamma })=S_i`$ and $`\text{obȷ}_i^+(M_1,w^\mathrm{\Gamma })\mathrm{\Gamma }`$. Thus, we can take $`w_1=w^\mathrm{\Gamma }`$, completing the proof of part (a).
To summarize the construction of part (a), we start with an arbitrary situation $`(M,w)`$ satisfying $`\mathrm{\Gamma }`$, convert it to a situation satisfying $`L_i\text{false}\mathrm{\Gamma }`$, essentially by modifying the $`i`$-accessibility relation at $`w`$ so that there are no worlds $`i`$-accessible from $`w`$ and $`w`$ is not accessible from any world, and then again modifying the $`i`$-accessibility relation at $`w`$ so that we get a structure $`(M_1,w_1)`$ such that $`\text{Obȷ}_i^+(M_i,w_i)=S_i`$. Note that in doing this construction, we did not change the propositional formulas true at $`w`$, nor did we change the worlds that were $`j`$-accessible from $`w`$ for $`ji`$. Thus, starting with a situation that satisfies a propositional formula $`\sigma `$, we can repeat this construction for each $`i`$ in turn, for $`i=1,\mathrm{},n`$. The resulting situation is $`(M_2,w_2)`$, and it clearly has the desired properties. This proves part (b).
For part (c), suppose $`(M^{},w^{})\mathrm{\Sigma }`$; let $`\text{Obȷ}_i^+(M^{},w^{})=S_i`$. By part (a), there is a situation $`(M_3,w_3)`$ such that $`\text{Obȷ}_i^+(M_3,w_3)=S_i`$ and $`(M_3,w_3)\mathrm{\Gamma }`$. Since the set of subjective formulas true at a situation $`(M,w)`$ is completely determined by $`\text{Obȷ}_i^+(M,w)`$, and $`\text{Obȷ}_i^+(M^{},w^{})=\text{Obȷ}_i^+(M_3,w_3)`$, it follows that $`(M_3,w_3)\mathrm{\Sigma }`$ as well.
How does this semantics compare to the canonical model semantics? First of all, it is easy to see that the axioms are sound. We write $`^{}\phi `$ if $`(M,w)^{}\phi `$ for every situation $`(M,w)`$. Then we have the following result.
###### Theorem 4.3
For all $`\alpha 𝒪𝒩_n`$, if $`\alpha `$ then $`^{}\alpha `$.
###### Proof
As usual, the proof is by induction on the length of a derivation. All that needs to be done is to show that all the axioms are sound. Again, this is straightforward. The proof in the case of A5<sub>n</sub> proceeds just as that in the proof of Theorem 3.10, using the fact that this semantics satisfies Levesque’s second property.
Moreover, we again get completeness for the sublanguage $`𝒪𝒩_n^{}`$.
###### Theorem 4.4
For all $`\alpha 𝒪𝒩_n^{}`$, $`\alpha `$ iff $`^{}\alpha `$.
###### Proof
As usual, it suffices to show that if $`\alpha `$ is consistent with the axioms, then $`\alpha `$ is satisfiable under the $`^{}`$ semantics. From Theorem 3.16, we know that $`\alpha `$ is satisfiable in the canonical model under the $``$ semantics. Thus, it suffices to show that for all formulas $`\alpha 𝒪𝒩_n^{}`$, we have $`(M^c,w)\alpha `$ iff $`(M^c,w)^{}\alpha `$. By Lemma 3.12, it suffices to consider formulas $`\alpha `$ in normal form. We proceed by induction on the structure of formulas. The only nontrivial case obtains if $`\alpha `$ is of the form $`N_i\alpha ^{}`$. Since $`\alpha `$ is in normal form, we can assume that $`\alpha ^{}`$ is basic. Suppose $`(M^c,w)^{}N_i\alpha ^{}`$. To show that $`(M^c,w)N_i\alpha ^{}`$, we must show that if $`w^{}_iw`$ and $`w^{}𝒦_i^c(w)`$, then $`(M^c,w^{})\alpha `$. By definition, if $`w^{}_iw`$, then $`\text{Obj}_i(M^c,w^{})=\text{Obj}_i(M^c,w)`$. Moreover, we must have $`\text{obj}_i(M^c,w^{})\text{Obj}_i(M^c,w)`$, for otherwise we would have $`w^{}𝒦_i^c(w)`$. Hence, we must have $`(M^c,w^{})^{}\alpha ^{}`$. By the induction hypothesis, we have $`(M^c,w^{})\alpha ^{}`$. Thus, $`(M^c,w)N_i\alpha ^{}`$, as desired.
For the converse, suppose that $`(M^c,w)N_i\alpha ^{}`$. We want to show that $`(M^c,w)^{}N_i\alpha ^{}`$. Suppose that $`(M^{},w^{})`$ is such that $`\text{Obj}_i(M^{},w^{})=\text{Obj}_i(M^c,w)`$ and $`\text{obj}_i(M^{},w^{})\text{Obj}_i(M^c,w)`$. We must show that $`(M^{},w^{})^{}\alpha `$. It is easy to see that for every situation $`(M,w)`$ and basic formula $`\phi `$, we have that $`(M,w)L_i\phi `$ iff $`(M,w)^{}L_i\phi `$ iff $`\phi `$ is in every set in $`\text{Obj}_i(M,w)`$. Thus, it follows that $`\text{subj}_i(M^{},w^{})=\text{subj}_i(M^c,w)`$. There must be a world $`w^{\prime \prime }`$ in $`M^c`$ such that $`(M^c,w^{\prime \prime })`$ agrees with $`(M^{},w^{})`$ on all basic formulas according to the $``$ semantics. Since $`\text{subj}_i(M^c,w^{\prime \prime })=\text{subj}_i(M^c,w)`$, it follows from Lemma 3.3 that $`w_iw^{\prime \prime }`$. Since $`\text{obj}_i(M^{},w^{})\text{Obj}_i(M^c,w)`$ and $`\text{obj}_i(M^{},w^{})=\text{obj}_i(M^c,w^{\prime \prime })`$, it follows that $`w^{\prime \prime }𝒦_i^c(w)`$. Since $`(M^c,w)N_i\alpha ^{}`$, we must have that $`(M^c,w^{\prime \prime })\alpha ^{}`$. And since $`(M^c,w^{\prime \prime })`$ and $`(M^{},w^{})`$ agree on basic formulas, it follows that $`(M^{},w^{})\alpha ^{}`$. Finally, since $``$ and $`^{}`$ agree for basic formulas, we have $`(M^{},w^{})^{}\alpha ^{}`$. This completes the proof that $`(M,w)N_i\alpha ^{}`$.
Although our axiomatization is complete for $`𝒪𝒩_n^{}`$, as we now show, it is not complete for the full language, for neither $``$ nor $`^{}`$. Since the axiomatization is sound for both $``$ and $`^{}`$, to prove incompleteness, it suffices to provide a formula which is satisfiable with respect to $`^{}`$ and not $``$, and another formula which is satisfiable with respect to $``$ and not $`^{}`$. As is shown in Proposition 4.5, $`O_i\neg O_jp`$ is satisfiable with respect to $`^{}`$ and (by Proposition 3.8) not with respect to $``$. On the other hand, it is easy to see that $`L_j\text{false}N_j\neg O_i\neg O_jp`$ is satisfiable with respect to $``$ (in fact, it is equivalent to $`L_j\text{false}`$); as shown in Proposition 4.6, it is not satisfiable with respect to $`^{}`$.
###### Proposition 4.5
$`O_i\neg O_jp`$ is satisfiable under the $`^{}`$ semantics.
###### Proof
Let $`S=\{\text{obȷ}_i^+(M,w)|(M,w)^{}\neg O_jp\}`$. By Proposition 4.2, there is a situation $`(M^{},w^{})`$ such that $`\text{Obȷ}_i^+(M^{},w^{})=S`$. We claim that $`(M^{},w^{})^{}O_i\neg O_jp`$. Clearly $`(M^{},w^{})^{}L_i\neg O_jp`$, since $`\neg O_jp`$ is true at all worlds $`i`$-accessible from $`w^{}`$. To see that $`(M^{},w^{})^{}N_iO_jp`$, suppose that $`\text{Obj}_i(M,w)=\text{Obj}_i(M^{},w^{})`$ and $`\text{obj}_i(M,w)\text{Obj}_i(M^{},w^{})`$. We want to show that $`(M,w)^{}O_jp`$. Suppose that $`(M,w)^{}\neg O_jp`$. By definition, $`\text{obȷ}_i^+(M,w)S`$, so there is some world $`w^{}𝒦_i^M^{}(w^{})`$ such that $`\text{obȷ}_i^+(M^{},w^{})=\text{obȷ}_i^+(M,w)`$. In particular, this means that $`\text{obj}_i(M^{},w^{})=\text{obj}_i(M,w)`$. But this contradicts the assumption that $`\text{obj}_i(M,w)\text{Obj}_i(M^{},w^{})`$. Thus, $`(M,w)O_jp`$ as desired, and $`(M^{},w^{})O_i\neg O_jp`$.
###### Proposition 4.6
There is a formula $`\beta `$ such that $`^{}\beta `$ but $`\neg \beta `$ is satisfiable under the $``$ semantics.
###### Proof
We first show that if $`\phi `$ is an $`i`$-objective formula that is satisfiable under the $`^{}`$ semantics, then $`^{}L_i\text{false}\neg N_i\neg \phi `$. For suppose that $`\phi `$ is satisfiable in a situation $`(M,w)`$. By Proposition 4.2, there is a situation $`(M^{},w^{})`$ such that $`\text{obȷ}_i^+(M^{},w^{})=\text{obȷ}_i^+(M,w)`$ and $`\text{Obȷ}_i^+(M^{},w^{})=\mathrm{}`$. This means that $`(M^{},w^{})^{}\phi L_i\text{false}`$. Now let $`(M^{},w^{})`$ be any situation satisfying $`L_i\text{false}`$. Then $`\text{Obj}_i(M^{},w^{})=\text{Obj}_i(M^{},w^{})=\mathrm{}`$, and $`\text{obj}_i(M^{},w^{})\text{Obj}_i(M^{},w^{})`$. It follows that $`(M^{},w^{})^{}\neg N_i\neg \phi `$. Thus, we have shown that $`^{}L_i\text{false}\neg N_i\neg \phi `$. Since, as we showed in Proposition 4.5, the formula $`O_j\neg O_ip`$ is satisfiable, this means that $`^{}L_i\text{false}\neg N_i\neg O_j\neg O_ip`$. On the other hand, since $`O_j\neg O_ip`$ is not satisfiable with respect to $``$, as we showed in Proposition 3.8, neither is $`\neg N_i\neg O_j\neg O_ip`$, and hence $`L_i\text{false}\neg N_i\neg O_j\neg O_ip`$ is not valid under the $``$ semantics. Indeed, $`L_i\text{false}N_i\neg O_j\neg O_ip`$ is equivalent to $`L_i\text{false}`$ under the $``$ semantics.
We can now show that our axiom system is incomplete for the full language with respect to both the $``$ and $`^{}`$ semantics.
###### Theorem 4.7
There exist formulas $`\alpha `$ and $`\beta `$ in $`𝒪𝒩_n`$ such that $`⊬\alpha `$ and $`\alpha `$, and $`⊬\beta `$ and $`^{}\beta `$.
###### Proof
By Propositions 3.8 and 4.5, we have that $`\neg O_j\neg O_ip`$, but $`\vDash ̸^{}\neg O_j\neg O_ip`$. Since $``$ is sound with respect to $`^{}`$, we cannot have $`\neg O_j\neg O_ip`$ (for otherwise we would have $`^{}\neg O_j\neg O_ip`$). Thus, we can take $`\alpha `$ to be $`\neg O_j\neg O_ip`$. A similar argument shows we can take $`\beta `$ to be $`L_i\text{false}\neg N_i\neg O_jO_ip`$.
The fact that neither $``$ nor $`^{}`$ is complete with respect to the axiomatization described earlier is not necessarily bad. We may be able to find a natural complete axiomatization. However, as we suggested above, the fact that $`\neg O_j\neg O_ip`$ is valid under the $``$ semantics suggests that this semantics does not quite satisfy our intuitions with regards to only-knowing for formulas in $`𝒪𝒩_n𝒪𝒩_n^{}`$. As we now show, $`^{}`$ also has its problems. We might hope that if $`\phi `$ is a satisfiable $`i`$-objective formula, then $`N_i\phi \neg L_i\phi `$ would be valid under the $`^{}`$ semantics. Unfortunately, it is not.
###### Proposition 4.8
The formula $`N_i\neg O_jpL_i\neg O_jp`$ is satisfiable under the $`^{}`$ semantics.
###### Proof
First we show that for any situation $`(M,w)`$ that satisfies $`O_jp`$, there exists another situation $`(M^{},w^{})`$ such that $`(M,w)`$ and $`(M^{},w^{})`$ agree on all basic formulas, but $`(M^{},w^{})^{}\neg O_jp`$. We can construct $`(M^{},w^{})`$ as follows: Choose a particular set $`\mathrm{\Gamma }\text{Obj}_j(M,w)`$. It easily follows from Proposition 4.2 that there is a situation $`(M^{},w^{})`$ such that $`\text{Obj}_j(M^{},w^{})=\text{Obj}_j(M,w)\{\mathrm{\Gamma }\}`$ and $`\text{obj}_j(M^{},w^{})=\text{obj}_j(M,w)`$. We now show that for any basic formula $`\phi `$, we have $`(M,w)^{}\phi `$ iff $`(M^{},w^{})^{}\phi `$. If $`\phi `$ is a $`j`$-objective formula, this is immediate from the construction. Thus, it suffices to deal with the case that $`\phi `$ is of the form $`L_j\phi ^{}`$. By Lemma 3.12, we can assume without loss of generality that $`\phi ^{}`$ is $`j`$-objective. Suppose that $`\phi ^{}`$ is a consistent $`j`$-objective formula. In this case, it is almost immediate from the definitions that if $`(M,w)^{}L_j\phi ^{}`$ then $`(M^{},w^{})^{}L_j\phi ^{}`$. For the converse, suppose that $`(M,w)^{}\neg L_j\phi ^{}`$. Then there is some world $`w^{\prime \prime }𝒦_j^M(w)`$ such that $`(M,w^{\prime \prime })\neg \phi ^{}`$. Since $`(M,w)^{}L_jp`$, we have that $`(M,w^{\prime \prime })^{}p\neg \phi ^{}`$. Let $`\psi `$ be a $`j`$-objective basic formula that is independent of $`p\neg \phi ^{}`$. Let $`\phi ^{\prime \prime }`$ be $`\phi ^{}p\neg \psi `$ if $`\psi \mathrm{\Gamma }`$, and $`\phi ^{}p\psi `$ otherwise. Since $`\psi `$ is independent of $`p\neg \phi ^{}`$, it follows that $`\phi ^{\prime \prime }`$ is consistent. Let $`\mathrm{\Delta }`$ be any $`j`$-set containing $`\phi ^{\prime \prime }`$. It must be the case that $`\mathrm{\Delta }\text{Obj}_j(M,w)`$, for if not, let $`(M^{},w^{})`$ be a situation such that $`\text{Obj}_j(M^{},w^{})=\text{Obj}_j(M,w)`$ and $`\text{obj}_j(M^{},w^{})=\mathrm{\Delta }`$ (such a situation exists by Proposition 4.2). Then $`(M^{},w^{})p`$, contradicting the assumption that $`(M,w)N_j\neg p`$. By construction, $`\mathrm{\Delta }\mathrm{\Gamma }`$. Thus, there is some world $`v𝒦_j^M^{}(w^{})`$ such that $`\text{obj}_j(M^{},v)=\mathrm{\Delta }`$. It follows that $`(M^{},v)^{}\neg \phi ^{}`$, so $`(M^{},w^{})^{}\neg L_j\phi ^{}`$. Thus, $`(M^{},w^{})`$ agrees with $`(M,w)`$ on all basic formulas. However, since $`\mathrm{\Gamma }\text{Obj}_j(M,w)`$, it follows that $`(M^{},w^{})\neg N_j\neg p`$, and hence that $`(M^{},w^{})\neg O_jp`$.
Let $`S=\{\text{obȷ}_i^+(M,w)|(M,w)^{}\neg O_jp\}`$. By Proposition 4.2, there is a situation $`(M^{},w^{})`$ such that $`\text{Obȷ}_i^+(M^{},w^{})=S`$. Clearly $`(M^{},w^{})^{}L_i\neg O_jp`$. We now show that $`(M^{},w^{})^{}N_i\neg O_jp`$ as well. For suppose that $`(M,w)`$ is a situation such that $`\text{Obj}_i(M,w)=\text{Obj}_i(M^{},w^{})`$ and $`\text{obj}_i(M,w)\text{Obj}_i(M^{},w^{})`$. Moreover, suppose, by way of contradiction, that $`(M,w)^{}O_jp`$. By the arguments above, it follows that there is a situation $`(M^{},w^{})`$ such that $`(M^{},w^{})^{}\neg O_jp`$ and $`(M,w)`$ and $`(M^{},w^{})`$ agree on all basic formulas. By construction, $`\text{obȷ}_i^+(M^{},w^{})S=\text{Obȷ}_i^+(M^{},w^{})`$, so $`\text{obj}_i(M^{},w^{})\text{Obj}_i(M^{},w^{})`$. Since $`\text{obj}_i(M,w)=\text{obj}_i(M^{},w^{})`$, we must also have $`\text{obj}_i(M,w)\text{Obj}_i(M^{},w^{})`$, contradicting the choice of $`(M,w)`$. Thus, $`(M,w)\neg O_jp`$, as desired, and so $`(M^{},w^{})L_i\neg O_jpN_i\neg O_jp`$.
### 4.1 Discussion
Proposition 4.8 shows that although the $`i`$-set semantics has the three properties we claimed were appropriate, $`N_i`$ and $`L_i`$ still do not always interact in what seems to be the appropriate way. Intuitively, the problem here is that there is more to $`i`$’s view of a world than just the $`i`$-objective basic formulas that are true there. We should really identify $`i`$’s view of a situation $`(M,w)`$ with the set of all $`i`$-objective formulas that are true there. In the canonical-model approach, the $`i`$-objective basic formulas that are true at a world can be shown to determine all the $`i`$-objective formulas that are true at that world. This is not true at all situations under the $`i`$-set approach.
Indeed, it is no longer true that the $`i`$-set approach has the second of the three properties once we take $`i`$’s view of $`(M,w)`$ to be $`\text{obȷ}_i^+(M,w)`$. For consider the situation $`(M^{},w^{})`$ constructed in the proof of Proposition 4.8. As the proof of that lemma shows, $`\{\text{obȷ}_i^+(M^{},w)|w𝒦_i^M^{}(w^{})\}\{\text{obȷ}_i^+(M^{},w^{})|\text{obj}_i(M^{},w^{})\text{Obj}_i(M^{},w^{})\}`$ does not include all maximal sets of $`i`$-objective formulas. In particular, it does not include those maximal sets that satisfy $`O_jp`$.
To summarize: while the i-set approach arguably has enough worlds, it suffers from the fact that the full complement of worlds is not always taken into account for $`L_i`$ and $`N_i`$, as the above example shows. While the canonical model approach does not run into this problem, it suffers from a perhaps more basic deficiency: there just are not enough worlds to begin with. For example, there is no world where $`O_iO_jp`$ is true.
We consider a different approach in the next section that attempts to deal with both problems.
## 5 What Properties Should Only Knowing Have?
Up to now, we have provided two semantics for only knowing. While both have properties we view as desirable, they also have properties that seem somewhat undesirable. This leads to an obvious question: What properties should only knowing have? Roughly speaking, we would like to have the multi-agent version of Levesque’s axioms, and no more. Of course, the problem here is axiom A5<sub>n</sub>. It is not so clear what the multi-agent version of that should be. The problem is one of circularity: We would like to be able to say that $`N_i\phi \neg L_i\phi `$ should hold for any consistent $`i`$-objective formula. The problem is that in order to say what the consistent formulas are, we need to define the axiom system. In particular, we have to make precise what this axiom should be.
To deal with this problem, we extend the language so that we can explicitly talk about satisfiability and validity in the language. We add a modal operator Val to the language. The formula Val($`\phi `$) should be read “$`\phi `$ is valid”. Of course, its dual Sat($`\phi `$), defined as $`\neg \text{Val(}\neg \phi \text{)}`$, should be read “$`\phi `$ is satisfiable”. With this operator in the language, we can replace A5<sub>n</sub> with
* $`\text{Sat(}\neg \alpha \text{)}(N_i\alpha \neg L_i\alpha )`$ if $`\alpha `$ is $`i`$-objective.
In addition, we have the following rules for reasoning about validity and satisfiability:
$`(\text{Val(}\phi \text{)}\text{Val(}\phi \psi \text{)})\text{Val(}\psi \text{)}`$.
Sat($`\phi `$), if $`\phi `$ is a satisfiable propositional formula.<sup>10</sup><sup>10</sup>10We can replace this by the simpler Sat($`p_1^{}\mathrm{}p_k^{}`$), where $`p_i^{}`$ is a literal—either a primitive proposition or its negation—and $`p_1^{}\mathrm{}p_k^{}`$ is consistent.
$`(\text{Sat(}\alpha \beta _1\text{)}\mathrm{}\text{Sat(}\alpha \beta _k\text{)}\text{Sat(}\gamma \delta _1\text{)}\mathrm{}\text{Sat(}\gamma \delta _m\text{)}\text{Val(}\alpha \gamma \text{)})`$
Sat($`L_i\alpha \neg L_i\neg \beta _1\mathrm{}\neg L_i\neg \beta _kN_i\gamma \neg N_i\neg \delta _1\mathrm{}\neg N_i\neg \delta _m`$),
if $`\alpha ,\beta _1,\mathrm{},\beta _k,\gamma ,\delta _1,\mathrm{},\delta _m`$ are $`i`$-objective formulas.
$`(\text{Sat(}\alpha \text{)}\text{Sat(}\beta \text{)})\text{Sat(}\alpha \beta \text{)}`$ if $`\alpha `$ is $`i`$-objective and $`\beta `$ is $`i`$-subjective.
From $`\phi `$ infer Val($`\phi `$).
Axiom V1 and the rule Nec<sub>V</sub> make Val what is called a normal modal operator. In fact, it can be shown to satisfy all the axioms of S5. The interesting clauses are clearly V2V4, which capture the intuitive properties of validity and satisfiability.
If we restrict to basic formulas, then V3 simplifies to $`(\text{Sat(}\alpha \beta _1\text{)}\mathrm{}\text{Sat(}\alpha \beta _k\text{)})\text{Sat(}L_i\alpha \neg L_i\neg \beta _1\mathrm{}\neg L_i\neg \beta _k\text{)}`$ (we can take $`\gamma ,\delta _1,\mathrm{},\delta _m`$ to be true to get this). The soundness of this axiom (interpreting Sat as satisfiability) follows using much the same arguments as those in the proof of Proposition 4.2. The soundness of V4 if we restrict to basic formulas follows from Lemma 3.4. More interestingly, it follows from the completeness proof given below that these axioms completely characterize satisfiability in K45<sub>n</sub>; together with the K45<sub>n</sub> axioms, they provide a sound and complete language for the language augmented with the Val operator.
Let AX consist of the axioms for $`𝒪𝒩`$ given earlier together with V1V4 and Nec<sub>V</sub>, except that A5<sub>n</sub> is replaced by A5$`{}_{}{}^{}{}_{n}{}^{}`$. AX is the axiom system that provides what we claim is the desired generalization of Levesque’s axioms to the multi-agent case. In particular, A5$`{}_{}{}^{}{}_{n}{}^{}`$ is the appropriate generalization of A5. The question is, of course, whether there is a semantics for which this is a complete axiomatization. We now provide one, in the spirit of the canonical-model construction of Section 3, except that, in the spirit of the extended situations of Section 2, we do not attempt to make the set of worlds used for evaluating $`L_i`$ and $`N_i`$ disjoint.
Let $`𝒪𝒩_n^+`$ be the extension of $`𝒪𝒩_n`$ to include the modal operator Val. For the remainder of this section, when we say “consistent”, we mean consistent with the axiom system AX. We define the extended canonical modelx, denoted $`M^e=(W^e,\pi ^e,𝒦_1^e,\mathrm{},𝒦_n^e,𝒩_1^e,\mathrm{},𝒩_n^e)`$, as follows:
* $`W^e`$ consists of the maximal consistent sets of formulas in $`𝒪𝒩_n^+`$.
* For all primitive propositions $`p`$ and $`wW^e`$, we have $`\pi ^e(w)(p)=\text{true}`$ iff $`pw`$.
* $`(w,w^{})𝒦_i^e`$ iff $`w/L_iw^{}`$.
* $`(w,w^{})𝒩_i^e`$ iff $`w/N_iw^{}`$.
In this canonical model, the semantics for $`L_i`$ and $`N_i`$ is defined in terms of the $`𝒦_i^e`$ and $`𝒩_i^e`$ relations, respectively:
$$\begin{array}{c}(M^e,w)L_i\alpha \text{ if }(M^e,w^{})\alpha \text{ for all }w^{}\text{ such that }(w,w^{})𝒦_i^e\text{.}\hfill \\ (M^e,w)N_i\alpha \text{ if }(M^e,w^{})\alpha \text{ for all }w^{}\text{ such that }(w,w^{})𝒩_i^e\text{.}\hfill \end{array}$$
We define the Val operator so that it corresponds to validity in the extended canonical model:
$`(M^e,w)\text{Val(}\alpha \text{)}`$ if $`(M^e,w^{})\alpha `$ for all worlds $`w^{}`$ in $`M^e`$.
We now want to show that every formula in a maximal consistent set is satisfied at a world in the extended canonical model. To do this, we need one preliminary result, showing that Val and Sat really correspond to provability and consistency in this framework.
###### Proposition 5.1
For every formula $`\phi 𝒪𝒩_n`$, if $`\phi `$ is provable then so is Val($`\phi `$), while if $`\phi `$ is not provable, then $`\neg \text{Val(}\phi \text{)}`$ is provable.
###### Proof
By Nec<sub>V</sub>, it is clear that if $`\phi `$ is provable, so is Val($`\phi `$). Thus, it remains to show that if $`\phi `$ is not provable, then $`\neg \text{Val(}\phi \text{)}`$ is. Using V1, it is easy to see that $`\neg \text{Val(}\phi \text{)}`$ is provably equivalent to Sat($`\neg \phi `$), so it suffices to show that if $`\phi `$ is not provable—i.e., if $`\neg \phi `$ is consistent—then Sat($`\neg \phi `$) is provable. We prove by induction on $`\phi `$ that if $`\phi `$ is consistent, then Sat($`\phi `$) is provable.
If $`\phi `$ is propositional, the result is immediate from V2. For the general case, we first use Lemma 3.12 to restrict attention to formulas in the canonical form specified by the lemma. Using standard modal reasoning (V1 and Nec<sub>V</sub>) it is easy to show that $`\text{Sat(}\phi \psi \text{)}(\text{Sat(}\phi \text{)}\text{Sat(}\psi \text{)})`$. Thus, it suffices to restrict attention to a conjunction in the form specified by the lemma. It is easy to see that if the conjunction is consistent, then each conjunct must be consistent. Using V4, it is easy to see that we can restrict attention to $`i`$-subjective formulas. By applying Lemma 3.12, we can assume without loss of generality that we are dealing with a consistent formula $`\phi `$ of the form $`L_i\alpha \neg L_i\neg \beta _1\mathrm{}\neg L_i\neg \beta _kN_i\gamma \neg N_i\neg \delta _1\mathrm{}\neg N_i\neg \delta _m`$, where $`\alpha ,\beta _1,\mathrm{},\beta _k,\gamma ,\delta _1,\mathrm{},\delta _m`$ are all $`i`$-objective. We can also assume that each of $`\alpha \beta _i`$, $`i=1,\mathrm{},k`$ and $`\gamma \delta _j`$, $`j=1,\mathrm{},m`$ are consistent, for otherwise we could easily show that $`\phi `$ is not consistent. Finally, we can show that $`\alpha \gamma `$ must be provable, for if not, by applying A5$`{}_{n}{}^{}{}_{}{}^{}`$, we can again show that $`\phi `$ is not consistent. We now apply the induction hypothesis to prove the result.
###### Corollary 5.2
Each formula in $`𝒪𝒩_n^+`$ is provably equivalent to a formula in $`𝒪𝒩_n`$.
###### Proof
We proceed by induction on the structure of formulas. The only nontrivial case is for formulas of the form Val($`\phi `$). By the induction hypothesis, $`\phi `$ is provably equivalent to a formula $`\phi ^{}𝒪𝒩_n`$. By straightforward modal reasoning using V1 and Nec<sub>V</sub>, we can show that Val($`\phi `$) is provably equivalent to Val($`\phi ^{}`$). By Proposition 5.1, Val($`\phi ^{}`$) is provably equivalent to either true or false, depending on whether $`\phi ^{}`$ is provable.
Using standard modal logic techniques, we can now prove the following result.
###### Theorem 5.3
$`M^e`$ is a K45<sub>n</sub> structure (that is, $`𝒦_i^e`$ and $`𝒩_i^e`$ are Euclidean and transitive). Moreover, for each world $`wW^e`$, we have $`(M^e,w)\alpha `$ iff $`\alpha w`$.
###### Proof
We leave it to the reader to check that the definition of $`𝒦_i^e`$ guarantees that $`M^e`$ is a K45<sub>n</sub> structure. Given Corollary 5.2, which allows us to restrict attention to $`\alpha 𝒪𝒩_n`$, the proof that $`(M^e,w)\alpha `$ iff $`\alpha w`$ is completely straight forward and follows the same lines as the usual proofs dealing with canonical models (see, for example, ).
We say that $`\alpha `$ is e-validx, denoted $`^e\alpha `$, if $`M^e\alpha `$, that is, if $`(M^e,w)\alpha `$ for all worlds $`wW^e`$. The following result is immediate from Theorem 5.3.
###### Corollary 5.4
$`^e\alpha `$ iff $`AX^{}\alpha `$.
Thus, AX is a sound and complete axiomatization of $`𝒪𝒩_n^+`$ with respect to the $`^e`$ semantics.
While AX is sufficient for our purposes, it comes at the expense of having to explicitly axiomatize validity as part of the logic itself. While we view the ability to axiomatize validity and satisfiability within the logic as a feature in our approach, it is reasonable to ask whether it is really necessary. One of the anonymous referees suggested to us the following interesting variant, which may avoid this complication, although at the expense of an infinite number of axiom schemas.
First we define an infinite sequence of languages $`𝒪𝒩_n^k`$ for $`k=0,1,2\mathrm{}`$:
* $`𝒪𝒩_n^0`$= $`𝒪𝒩_n^{}`$
* $`𝒪𝒩_n^{k+1}`$= $`\{\alpha |`$ $`\alpha `$ is a Boolean combination of formulas of $`𝒪𝒩_n^k`$ together with formulas of the form $`L_i\alpha `$ or $`N_i\alpha `$ for $`\alpha 𝒪𝒩_n^k`$$`\}`$
Roughly, each language adds another level of nestings of only knowing with varying agent indices. For example, $`𝒪𝒩_n^{k+1}`$ contains the formula $`O_{i_0}O_{i_1}\mathrm{}O_{i_{k+1}}p`$, where $`i_ji_{j+1}`$, something that cannot be expressed in $`𝒪𝒩_n^k`$.
Let AX consist of the the axioms A1<sub>n</sub>A4<sub>n</sub> as before together with the following set of axioms
> A5$`{}_{}{}^{k+1}{}_{n}{}^{}`$. $`N_i\alpha \neg Ł_i\alpha `$, where $`\alpha `$ is an i-objective $`𝒪𝒩_n^k`$ formula which is consistent with respect to A1<sub>n</sub>A4<sub>n</sub>,A5$`{}_{}{}^{1}{}_{n}{}^{}`$, A5$`{}_{}{}^{2}{}_{n}{}^{}`$,$`\mathrm{}`$ A5$`{}_{}{}^{k}{}_{n}{}^{}`$.
It is not hard to show that the axioms are sound with respect to the semantics. Whether they are also complete remains an open problem.
Apart from the question of axiomatization, how does the $`^e`$ semantics compare to our earlier two? Clearly, they differ. It is easy to see that the formula $`O_i\neg O_jp`$, which was not satisfiable under $`^c`$, is satisfiable under $`^e`$. In addition, the formula $`N_i\neg O_jpL_i\neg O_jp`$, which is satisfiable under $`^{}`$, is not satisfiable under $`^e`$. In both cases, it seems that the behavior of $`^e`$ is more appropriate. On the other hand, all three semantics agree in the case where our intuitions are strongest, $`𝒪𝒩_n^{}`$. Since the axiom system AX characterizes how our earlier two semantics deal with $`𝒪𝒩_n^{}`$, this is shown by the following result.
###### Theorem 5.5
If $`\phi 𝒪𝒩_n^{}`$, then $`AX\phi `$ iff $`AX^{}\phi `$.
###### Proof
It is easy to see that each axiom of AX is sound in AX. It follows that $`\text{AX}\phi `$ implies AX$`{}_{}{}^{}\phi `$. For the converse, it suffices to show that if $`\phi 𝒪𝒩_n^{}`$ is consistent with AX, then it is also consistent with AX, i.e., that Sat($`\phi `$) holds. We show this by induction on the structure of $`\phi `$, much in the same way we proved Proposition 5.1. We can assume without loss of generality that $`\phi `$ is a conjunction in the normal form described Lemma 3.12. It is easy to see that if we can deal with the case that $`\phi `$ is an $`i`$-subjective formula, then we can deal with arbitrary $`\phi `$ by repeated applications of V4 followed by an application of V2. Thus, suppose that $`\phi `$ is an $`i`$-subjective formula which is consistent with AX. We can assume that $`\phi `$ is of the form $`L_i\alpha \neg L_i\neg \beta _1\mathrm{}\neg L_i\neg \beta _kN_i\gamma \neg N_i\neg \delta _1\mathrm{}\neg N_i\neg \delta _m`$. We must have that $`\alpha \beta _j`$ is consistent for $`j=1,\mathrm{},k`$, and that $`\gamma \delta _l`$ is AX-consistent for $`l=1,\mathrm{},m`$, for otherwise $`\phi `$ would not be AX-consistent. Similarly, by Lemma 3.14, we must have that $`\alpha \gamma `$ is K45<sub>n</sub>-provable, otherwise $`\phi `$ would not be AX-consistent. We can now apply V3 and the inductive hypothesis to show that $`\alpha `$ is AX-consistent.
Thus, we maintain all the benefits of the earlier semantics with this approach. Moreover, while it is just as intractable to compute whether $`(M^e,w)\alpha `$ for a particular world $`w`$ in the extended canonical model as it was in all our other approaches, we can show that the validity problem for this logic is no harder than that for K45<sub>n</sub> alone. It is PSPACE-complete.
###### Theorem 5.6
The problem of deciding if $`AX^{}\phi `$ is PSPACE-complete.
###### Proof
PSPACE hardness follows from the PSPACE hardness of K45<sub>n</sub> .<sup>11</sup><sup>11</sup>11The result in is proved only for $`\mathrm{KD45}_n`$, but the same proof applies to K45<sub>n</sub>. We sketch the proof of the upper bound. First of all, observe that it suffices to deal with the case that $`\phi `$ is in $`𝒪𝒩_n`$, since we can then apply the arguments of Corollary 5.2 to remove all occurrences of Val from inside out. We consider the dual problem of consistency. Thus, we want to check if Sat($`\alpha `$) holds. The first step is to convert $`\alpha `$ to the normal form of Lemma 3.12. Observe that $`\alpha `$ is consistent iff at least one of the disjuncts is consistent. Although the conversion to normal form may result in exponentially many disjuncts, each one is no longer than $`\alpha `$. Thus, we deal with them one by one, without ever writing down the full disjunction. It suffices to show that we can decide if each disjunct is consistent in polynomial space, since we can then erase all the work and start over for the next disjunct (with a little space necessary for bookkeeping). We now proceed much as in the proof of Proposition 5.1. By applying V4 repeatedly and then V2 (as in the previous theorem), it suffices to deal with $`i`$-subjective formulas. We then apply V3 to get simpler formulas, and repeat the procedure. We remark that this gives another PSPACE decision procedure for K45<sub>n</sub>, quite different from that presented in .
To what extent do the three properties we have been focusing on hold under the $`^e`$ semantics? Suppose we take the conceivable states from $`i`$’s point of view to be the maximal consistent sets of $`i`$-objective formulas with respect to AX, or equivalently, the set of $`i`$-objective formulas true at some world in $`M^e`$. Let $`\text{obȷ}_i^e(M^e,w)`$x consist of all the $`i`$-objective formulas true at world $`w`$ in the extended canonical model (under the $`^e`$ semantics) and let $`\text{Obȷ}_i^e(M^e,w)`$x = $`\{\text{obȷ}_i^e(M^e,w^{})|w^{}𝒦_i^e(w)\}`$. It is easy to see that the first two properties we isolated hold under this interpretation of conceivable state. However, it is quite possible that the “possible states” at a world $`(M^c,w)`$, that is, $`\text{Obȷ}_i^e(M^c,w)`$, and the “impossible states”, that is, $`\{\text{obȷ}_i^e(M^c,w^{})|w^{}_iw,w𝒦_i^e(w)\}`$ are not disjoint.
Interestingly, this semantics does not satisfy the third property we isolated. Not all subsets of conceivable states arise as the set of possible states at some situation $`(M^e,w)`$. A proof analogous to that of Lemma 3.7 shows that $`\text{Obȷ}_i^e(M,w)`$ is always limit closed. In the canonical model approach, limit closure prevents an agent from considering certain desirable sets of states. Now this is no longer the case, despite limit closure. Roughly speaking, we avoid problems by having in a precise sense “enough” possibilities. More precisely, given any consistent $`i`$-objective formula $`\alpha `$, it is possible for agent $`i`$ to only know $`\alpha `$ by virtue of considering all maximal sets of $`i`$-objective formulas which contain $`\alpha `$. Thus, as we suggested earlier, the third property we isolated in the beginning turns out to be somewhat too strong—it is sufficient but not necessary once we allow for enough possibilities.
## 6 Multi-Agent Nonmonotonic Reasoning
In this section, we demonstrate that the logic developed in Section 5 captures multi-agent autoepistemic reasoning in a reasonable way. We do this in two ways. First we show by example that the logic can be used to derive some reasonable nonmonotonic inferences in a multi-agent context. We then show that the logic can be used to extend the definitions of stable sets and stable expansions originally developed for single agent autoepistemic logic to the multi-agent setting.
### 6.1 Formal Derivations of Nonmonotonic Inferences
In this section, we provide two examples of how the logic can be used for nonmonotonic reasoning.
###### Example 6.1
Let $`p`$ be agent $`i`$’s secret and suppose $`i`$ makes the following assumption: unless I know that $`j`$ knows my secret assume that $`j`$ does not know it. We can prove that if this assumption is all $`i`$ believes then he indeed believes that $`j`$ does not know his secret. Formally, we can show
$$O_i(\neg L_iL_jp\neg L_jp)L_i\neg L_jp.$$
A formal derivation of this theorem can be obtained as follows. Let $`\alpha =\neg L_iL_jp\neg L_jp`$. The justifications in the following derivation indicate which axioms or previous derivations have been used to derive the current line. PL or K45<sub>n</sub> indicate that reasoning in either standard propositional logic or K45<sub>n</sub>, which are subsumed by AX, is used without further analysis.
| 1. | $`O_i\alpha L_i\alpha `$ | PL |
| --- | --- | --- |
| 2. | $`O_i\alpha N_i\neg \alpha `$ | PL |
| 3. | $`(L_i\alpha \neg L_iL_jp)L_i\neg L_jp`$ | K45<sub>n</sub> |
| 4. | $`N_i\neg \alpha (N_i\neg L_iL_jpN_iL_jp)`$ | K45<sub>n</sub> |
| 5. | Sat($`p`$) | V2 |
| 6. | $`\text{Sat(}p\text{)}\text{Sat(}\neg L_jp\text{)}`$ | V3 |
| 7. | Sat($`\neg L_jp`$) | V1, PL |
| 8. | $`\text{Sat(}\neg L_jp\text{)}(N_iL_jp\neg L_iL_jp)`$ | A5$`{}_{n}{}^{}{}_{}{}^{}`$ |
| 9. | $`N_iL_jp\neg L_iL_jp`$ | PL |
| 10. | $`O_i\alpha \neg L_iL_jp`$ | 2; 4; 9; PL |
| 11. | $`O_i\alpha L_i\neg L_jp`$ | 1; 3; 10; PL |
To see that $`i`$’s beliefs may evolve nonmonotonically given that $`i`$ knows only $`\alpha `$, assume that $`i`$ finds out that $`j`$ has found out about the secret. Then $`i`$’s belief that $`j`$ does not believe the secret will be retracted. In fact, $`i`$ will believe that $`j`$ does believe the secret. Formally, we can show
$$O_i(L_jp\alpha )L_iL_jp.$$
Notice that the logic itself is a regular monotonic logic; the nonmonotonicity of agent $`i`$’s beliefs is hidden within the $`O_i`$-operator.
All the formulas that appear in the proof above are in $`𝒪𝒩_n^{}`$. Thus, we could have used the somewhat simpler proof theory of Section 3.1. To obtain an example where we need the full power of AX, simply replace $`L_jp`$ by $`O_jp`$, that is $`i`$ now uses the default that unless he knows that $`j`$ only knows $`p`$, then he assumes that $`j`$ does not only know $`p`$. In other words, $`i`$ (prudently) makes rather cautious assumptions about $`i`$’s epistemic state and assumes that $`i`$ usually knows more than just $`p`$. The proof is very similar to the one above. The only difference is that we now have to establish that Sat($`\neg O_jp`$) is provable, which is straightforward.
###### Example 6.2
Now let $`p`$ stand for “Tweety flies”.<sup>13</sup><sup>13</sup>13Regular readers of papers on nonmonotonic logic will no doubt be gratified to see Tweety’s reappearance. We want to show that if $`j`$ knows that all $`i`$ knows about Tweety is that by default it flies, then $`j`$ knows that $`i`$ believes that Tweety flies. As before, we capture the fact that all $`i`$ believes is that, by default, Tweety flies, by saying that all $`i`$ believes is that, unless $`i`$ believes that Tweety does not fly, then Tweety flies. Thus, we want to show
$$L_jO_i(\neg L_i\neg pp)L_jL_ip.$$
We proceed as follows:
| 1. | $`O_i(\neg L_i\neg pp)L_ip`$ | as above, with $`L_jp`$ replaced by $`\neg p`$ |
| --- | --- | --- |
| 2. | $`L_j(O_i(\neg L_i\neg pp)L_ip)`$ | 1; Nec<sub>n</sub> |
| 3. | $`L_jO_i(\neg L_i\neg pp)L_jL_ip`$ | 2; K45<sub>n</sub> |
In a sense, $`i`$ is able to reason about $`j`$’s ability to reason nonmonotonically essentially by simulating $`j`$’s reasoning pattern.
A situation where $`i`$ knows that all $`j`$ knows is $`\alpha `$ seems hardly attainable in practice, since an agent usually has at best incomplete information about another agent’s beliefs. It would seem much more reasonable if we could say that $`i`$ knows that $`\alpha `$ is all $`j`$ knows about some relevant subject, say Tweety. This issue is dealt with in , where the canonical-model approach is extended to allow statements of the kind that all agent $`i`$ knows about $`x`$ is $`y`$. It is shown that the forms of nonmonotonic reasoning just described, when restricted to $`𝒪𝒩_n^{}`$, go through just as well with the weaker notion of only knowing about.
### 6.2 $`i`$-Stable Sets and $`i`$-Stable Expansions
Single-agent autoepistemic logic was developed by Moore using the concepts of stable sets and stable expansions. Levesque proved that there is a close relationship between stable sets and only-knowing in the single-agent case. Here we prove an analogous relationship for the multi-agent case. We first need to define a multi-agent analogue of stable sets.
In the single-agent case, it is well known that a stable set is a complete set of formulas that agent $`i`$ could know in some situation; that is, a set $`S`$ is stable if and only if there is a situation $`(W,w)`$ such that $`S=\{\alpha |(W,w)L\alpha \}`$. This is the intuition that we want to extend to the multi-agent setting, where the underlying language is now $`𝒪𝒩_n`$. First we define logical consequence in the extended canonical model in the usual way: If $`\mathrm{\Gamma }`$ is a set of formulas, we write $`M^e\mathrm{\Gamma }`$ if $`M^e\gamma ^{}`$ for each formula $`\gamma ^{}\mathrm{\Gamma }`$. We say that $`\gamma `$ is an e-consequence of $`\mathrm{\Gamma }`$, and write $`\mathrm{\Gamma }^e\gamma `$, exactly if $`M^e\mathrm{\Gamma }`$ implies $`M^e\gamma `$.
###### Definition 6.3
Let $`\mathrm{\Gamma }`$ be a set of formulas in $`𝒪𝒩_n`$. $`\mathrm{\Gamma }`$ is called $`i`$-stable iff
1. if $`\mathrm{\Gamma }^e\gamma `$ then $`\gamma \mathrm{\Gamma }`$,
2. if $`\alpha \mathrm{\Gamma }`$ then $`L_i\alpha \mathrm{\Gamma }`$,
3. if $`\alpha \mathrm{\Gamma }`$ then $`\neg L_i\alpha \mathrm{\Gamma }`$.
Note that the only difference between $`i`$-stable sets and the original definition of stable sets is in condition (a), which requires $`i`$-stable sets to be closed under e-consequence instead of tautological consequence (i.e., logical consequence in propositional logic) as in the single-agent case. Using e-consequence rather than tautological consequence makes no difference in the single-agent case; in the multi-agent case it does. Intuitively, we want to allow the agents to use e-consequence here to capture the intuition that it is common knowledge that all agents are perfect reasoners under the extended canonical model semantics. For example, if agent $`i`$ believes $`\neg L_jp`$ for a different agent $`j`$, then we want him to also believe $`N_j\neg L_jp`$. To do this, we need to close off under e-consequence.
The next theorem shows that $`i`$-stable sets do indeed satisfy the intuitive requirement. We define an $`i`$-epistemicx state to be a set $`\mathrm{\Gamma }`$ of $`𝒪𝒩_n`$-formulas such that for some situation $`(M^e,w)`$ in the extended canonical model, $`\mathrm{\Gamma }=\{\alpha 𝒪𝒩_n|(M^e,w)L_i\alpha \}`$; in this case, we say that $`\mathrm{\Gamma }`$ is the $`i`$-epistemic situation corresponding to $`(M^e,w)`$. For a set of $`𝒪𝒩_n`$-formulas $`\mathrm{\Gamma }`$, let $`\overline{\mathrm{\Gamma }}=\{\gamma |`$ $`\gamma 𝒪𝒩_n`$ and $`\gamma \mathrm{\Gamma }\}`$, $`L_i\mathrm{\Gamma }=\{L_i\gamma |\gamma \mathrm{\Gamma }\}`$, and $`\neg L_i\overline{\mathrm{\Gamma }}=\{\neg L_i\gamma |\gamma \overline{\mathrm{\Gamma }}\}`$.
###### Theorem 6.4
Let $`\mathrm{\Gamma }`$ be a set of $`𝒪𝒩_n`$-formulas. $`\mathrm{\Gamma }`$ is $`i`$-stable iff $`\mathrm{\Gamma }`$ is an $`i`$-epistemic state.
###### Proof
It is straightforward to show that every $`i`$-epistemic state is $`i`$-stable. To show the converse, let $`\mathrm{\Gamma }`$ be $`i`$-stable. We need to show that it is also an $`i`$-epistemic state. Certainly $`L_i\mathrm{\Gamma }`$ is consistent. If $`\mathrm{\Gamma }`$ contains all $`𝒪𝒩_n`$-formulas, that is, if agent $`i`$ is inconsistent, then let $`(M^e,w)`$ be a situation where $`𝒦_i^e(w)=\mathrm{}`$; such a situation clearly exists. Then $`\mathrm{\Gamma }`$ is the $`i`$-epistemic situation corresponding to $`(M^e,w)`$. If $`\mathrm{\Gamma }`$ is a proper subset of the $`𝒪𝒩_n`$-formulas, then $`\mathrm{\Gamma }`$ must be consistent. (If $`\mathrm{\Gamma }`$ were inconsistent, by the first property of stable sets, $`\mathrm{\Gamma }`$ would contain all formulas.) In particular, this means that $`p\neg p\mathrm{\Gamma }`$ for a primitive proposition $`p`$. The second property of stable sets guarantees that $`L_i\mathrm{\Gamma }\mathrm{\Gamma }`$, while the third guarantees that $`\neg L_i(p\neg p)\mathrm{\Gamma }`$. Since $`\mathrm{\Gamma }`$ is consistent, so is $`L_i\mathrm{\Gamma }\{\neg L_i(p\neg p)\}`$. $`W^e`$ consists of all the maximal consistent sets (with respect to AX); thus, there must be some $`w^{}W^e`$ that contains $`L_i\mathrm{\Gamma }\{\neg L_i(p\neg p)\}`$. We claim that $`\mathrm{\Gamma }`$ is the $`i`$-epistemic state corresponding to $`(M^e,w^{})`$. Thus, we must show that $`\phi \mathrm{\Gamma }`$ iff $`(M^e,w^{})L_i\phi `$. To see this, first suppose that $`\phi \mathrm{\Gamma }`$. Thus, $`L_i\phi L_i\mathrm{\Gamma }`$. By construction, $`w^{}`$ contains $`L_i\mathrm{\Gamma }`$. By Theorem 5.3, we have that $`(M^e,w^{})L_i\phi `$. On the other hand, if $`\phi \mathrm{\Gamma }`$, then, since $`\mathrm{\Gamma }`$ is $`i`$-stable, we have that $`\neg L_i\phi \mathrm{\Gamma }`$. By the previous argument, it follows that $`(M^e,w^{})L_i\neg L_i\phi `$. From $`\mathrm{𝐀𝟒}_n`$, it follows that $`(M^e,w^{})\neg L_i\phi `$, and hence $`(M^e,w^{})\vDash ̸L_i\phi `$. This proves the claim.
Moore defined the notion of a stable expansion of a set $`A`$ of formulas in the single-agent case. Intuitively, a stable expansion of $`A`$ is a stable set containing $`A`$ all of whose formulas can be justified, given $`A`$ and the formulas believed in that stable set. Further discussion and justification of the notion of stable expansion can be found in . Rather than discussing this here, we go directly to our multi-agent generalization of Moore’s notion.
###### Definition 6.5
Let $`A`$ be a set of $`𝒪𝒩_n`$-formulas. $`\mathrm{\Gamma }`$ is called an $`i`$-stable expansionx of $`A`$ iff $`\mathrm{\Gamma }=\{\gamma 𝒪𝒩_n|AL_i\mathrm{\Gamma }\neg L_i\overline{\mathrm{\Gamma }}^e\gamma \}`$.
It is easy to see that an $`i`$-stable expansion is an $`i`$-stable set. The definition of $`i`$-stable expansions looks exactly like Moore’s definition of stable expansions except that we again use e-consequence instead of tautological consequence. As in the case of stable sets, this is necessary to capture the fact that it is common knowledge that all agents can do reasoning under the extended canonical model semantics.
We now generalize a result of Levesque’s (who proved it for the single-agent case), showing that the $`i`$-stable expansions of a formula $`\alpha `$ correspond precisely to the different situations where $`i`$ only knows $`\alpha `$. We first need a lemma.
###### Lemma 6.6
Let $`(M^e,w)`$ be a situation with
$$\mathrm{\Sigma }=\{L_i\gamma |(M^e,w)|=L_i\gamma \}\{\neg L_i\gamma |(M^e,w)|=\neg L_i\gamma \}.$$
For any $`\alpha `$, there is an $`i`$-objective formula $`\alpha ^{}`$ such that
1. $`\mathrm{\Sigma }^e\alpha \alpha ^{}`$,
2. $`(M^e,w)(L_i\alpha L_i\alpha ^{})(N_i\alpha N_i\alpha ^{})`$.
###### Proof
Given a formula $`\phi `$, we say that a subformula $`L_i\psi `$ of $`\phi `$ occurs at top level if it is not in the scope of any modal operators. Let $`\alpha ^{}`$ be the result of replacing each top-level subformula of $`\alpha `$ of the form $`L_i\gamma `$ (resp., $`N_i\gamma `$) by true if $`\mathrm{\Sigma }^eL_i\gamma `$ (resp., $`\mathrm{\Sigma }^eN_i\gamma `$) and by false otherwise. Clearly $`\alpha ^{}`$ is $`i`$-objective. Moreover, a trivial argument by induction on the structure of $`\alpha `$ shows that $`\mathrm{\Sigma }\alpha \alpha ^{}`$: If $`\alpha `$ is a primitive proposition then $`\alpha ^{}=\alpha `$; if $`\alpha `$ is of the form $`\alpha _1\alpha _2`$ or $`\neg \alpha ^{}`$, then the result follows easily by the induction hypothesis; if $`\alpha `$ is of the form $`L_j\beta `$ for $`ji`$, then $`\alpha =\alpha ^{}`$, since $`\alpha `$ has no top-level subformulas of the form $`L_i\phi `$; finally, if $`\alpha `$ is of the form $`L_i\beta `$, then $`\alpha ^{}`$ is either true or false, depending on whether $`L_i\beta `$ is in $`\mathrm{\Sigma }`$. Since either $`L_i\beta `$ or $`\neg L_i\beta `$ must be in $`\mathrm{\Sigma }`$, the result is immediate in this case too. Part (b) follows immediately from part (a), since if $`w^{}_iw`$, we must have $`(M^e,w^{})\mathrm{\Sigma }`$, so $`(M^e,w^{})\alpha \alpha ^{}`$.
###### Theorem 6.7
Let $`w`$ be a world in the extended canonical model and let $`\mathrm{\Gamma }`$ be the $`i`$-epistemic state corresponding to $`(M^e,w)`$. Then, for every $`𝒪𝒩_n`$-formula $`\alpha `$, we have that $`(M^e,w)|=O_i\alpha `$ iff (a) $`\mathrm{\Gamma }`$ is an $`i`$-stable expansion of $`\{\alpha \}`$ and (b) $`𝒦_i^e(w)`$ and $`𝒩_i^e(w)`$ are disjoint.
###### Proof
Let $`\mathrm{\Sigma }=L_i\mathrm{\Gamma }\neg L_i\overline{\mathrm{\Gamma }}`$. To prove the “only if” direction, suppose $`(M^e,w)|=O_i\alpha `$. The disjointness of $`𝒦_i^e(w)`$ and $`𝒩_i^e(w)`$ follows immediately from the fact that $`(M^e,w)|=L_i\alpha N_i\neg \alpha `$. To prove that $`\mathrm{\Gamma }`$ is an $`i`$-stable expansion of $`\{\alpha \}`$, it suffices to show that for all $`𝒪𝒩_n`$-formulas $`\beta `$, we have $`(M^e,w)|=L_i\beta \mathrm{iff}\{\alpha \}\mathrm{\Sigma }^e\beta `$.
First suppose that $`\{\alpha \}\mathrm{\Sigma }^e\beta `$. Since $`(M^e,w)|=\{O_i\alpha \}\mathrm{\Sigma }`$, and every formula in $`\mathrm{\Sigma }`$ is of the form $`L_i\gamma `$ or $`\neg L_i\gamma `$, it easily follows that $`(M^e,w)|=L_i\alpha `$ and $`(M^e,w)|=L_i\mathrm{\Sigma }`$. Hence, $`(M^e,w^{})|=\alpha `$ and $`(M^e,w^{})|=\mathrm{\Sigma }`$ for every $`w^{}𝒦_i^e(w)`$. It follows that $`(M^e,w^{})|=\beta `$ and, therefore, $`(M^e,w)|=L_i\beta `$.
For the converse, suppose that $`(M^e,w)|=L_i\beta `$. We want to show that $`\{\alpha \}\mathrm{\Sigma }^e\beta `$. By Lemma 6.6, we can assume without loss of generality that $`\beta `$ is $`i`$-objective. (For if not, we can replace $`\beta `$ by an $`i`$-objective $`\beta ^{}`$ such that $`\mathrm{\Sigma }^e\beta \beta ^{}`$ and $`(M^e,w)L_i\beta L_i\beta ^{}`$, prove the result for $`\beta ^{}`$, and conclude that it holds for $`\beta `$ as well.) To show that $`\{\alpha \}\mathrm{\Sigma }^e\beta `$, we must show that for all worlds $`w^{}`$ such that $`(M^e,w^{})|=\{\alpha \}\mathrm{\Sigma }`$, we have $`(M^e,w^{})\beta `$. By Lemma 6.6, there is an $`i`$-objective formula $`\alpha ^{}`$ such that $`(M^e,w)O\alpha O\alpha ^{}`$ and $`\mathrm{\Sigma }^e\alpha \alpha ^{}`$. Thus, $`(M^e,w)N_i\alpha ^{}L_i\beta `$. By the arguments of Lemma 3.14 (which apply without change to the extended canonical model semantics), it must be the case that $`^e\alpha ^{}\beta `$. Since $`\mathrm{\Sigma }^e\alpha \alpha ^{}`$, it follows that $`\{\alpha \}\mathrm{\Sigma }^e\beta `$. Since $`(M^e,w^{})|=\{\alpha \}\mathrm{\Sigma }`$ by assumption, we have that $`(M^e,w^{})|=\beta `$, as desired.
To prove the “if” direction of the theorem, suppose that $`𝒦_i^e(w)`$ and $`𝒩_i^e(w)`$ are disjoint and that for all $`𝒪𝒩_n`$-formulas $`\beta `$, we have $`(M^e,w)|=L_i\beta \mathrm{iff}\{\alpha \}\mathrm{\Sigma }^e\beta `$. We need to show that $`(M^e,w)|=O_i\alpha `$, that is, $`(M^e,w)|=L_i\alpha N_i\neg \alpha `$.
Since $`\{\alpha \}\mathrm{\Sigma }^e\alpha `$, the fact that $`(M^e,w)|=L_i\alpha `$ follows immediately. To prove that $`(M^e,w)|=N_i\neg \alpha `$, let $`w^{}𝒩_i^e(w)`$ and assume, to the contrary, that $`(M^e,w^{})|=\alpha `$. Since $`w^{}𝒩_i^e(w)`$, it follows that $`\mathrm{\Sigma }w^{}`$. Hence $`w/L_iw^{}`$, from which $`w^{}𝒦_i^e(w)`$ follows, contradicting the assumption that $`𝒦_i^e(w)`$ and $`𝒩_i^e(w)`$ are disjoint.
## 7 Conclusion
We have provided three semantics for multi-agent only-knowing. All agree on the subset $`𝒪𝒩_n^{}`$, but they differ on formulas involving nested $`N_i`$’s. Although a case can be made that the $`^e`$ semantics comes closest to capturing our intuitions for “knowing at most”, our intuitions beyond $`𝒪𝒩_n^{}`$ are not well grounded. It would certainly help to have more compelling semantics corresponding to AX.
On the other hand, it can be argued that semantics does not play quite as crucial a role when dealing with knowing at most as in other cases. The reason is that the structures we must deal with, in general, have uncountably many worlds. For example, whichever of the three semantics we use, there must be uncountably many worlds $`i`$-accessible from a situation $`(M,w)`$ satisfying $`O_ip`$, at least one for every $`i`$-set that includes $`p`$. To the extent that we are interested in proof theory, the proof theory associated with $`^e`$, characterized by the axiom system AX, seems quite natural. The fact that the validity problem is no harder in this setting than that for K45<sub>n</sub> adds further support to its usefulness. Of course, as we suggested above, rather than only knowing, it seems more appropriate to reason about only knowing about a certain topic. Lakemeyer provides a semantics for only knowing about, using the canonical-model approach. It would be interesting to see if this can also be done using the other approaches we have explored here. |
warning/0001/hep-ph0001155.html | ar5iv | text | # Parametrization of polarized parton distribution functions
## 1 Introduction
After the EMC finding of a proton-spin issue, many polarized deep-inelastic scattering (DIS) experiments have been done on spin structure of the nucleon. From these experimental data and theoretical studies, we think that the nucleon spin is carried not only quarks but also gluons and their angular momenta. However, we do not have a clear idea even on the antiquark and gluon contributions, which are difficult to be determined by the present lepton-nucleon DIS data. The situation should become clearer in the near future because RHIC-Spin experiments will provide valuable information on these distributions.
We tried to determine the polarized parton distribution functions (PDFs) by using existing spin asymmetry $`A_1`$ data for understanding the present situation and for suggesting the importance of future experimental studies. The following discussions are based on the work in Ref. 1 with the members of the Asymmetry Analysis Collaboration (AAC). In Sec. 2, we explain how to calculate $`A_1`$ in terms of the unpolarized and polarized parton distributions. Then, the actual parametrization and $`\chi ^2`$-analysis method are discussed in Sec. 3. Our results are shown in Sec. 4 and conclusions are given in Sec. 5.
## 2 Parton model analysis of polarized DIS data
There are many measurements of the spin asymmetry $`A_1`$ for the proton, neutron, and deuteron. To use these experimental data in our analysis, we express $`A_1`$ as
$$A_1(x,Q^2)\frac{g_1(x,Q^2)}{F_1(x,Q^2)}=\frac{2x[1+R(x,Q^2)]}{F_2(x,Q^2)}g_1(x,Q^2),$$
(1)
where $`F_1`$ and $`F_2`$ are unpolarized structure functions. The function $`R(x,Q^2)`$ is given by $`R(x,Q^2)=\sigma _L/\sigma _T`$, where $`\sigma _L`$ and $`\sigma _T`$ are absorption cross sections of longitudinal and transverse photons, and it is determined experimentally in reasonably wide $`Q^2`$ and $`x`$ ranges in the SLAC experiment. $`^\mathrm{?}`$ The polarized structure function $`g_1(x,Q^2)`$ is expressed as
$$g_1(x,Q^2)=\frac{1}{2}\underset{i=1}{\overset{n_f}{}}e_i^2\left\{\mathrm{\Delta }C_q(x,\alpha _s)\left[\mathrm{\Delta }q_i(x,Q^2)+\mathrm{\Delta }\overline{q}_i(x,Q^2)\right]+\mathrm{\Delta }C_g(x,\alpha _s)\mathrm{\Delta }g(x,Q^2)\right\},$$
(2)
where $`e_i`$ is the electric charge of a quark, and the convolution $``$ is defined by $`f(x)g(x)=_x^1\frac{dy}{y}f\left(\frac{x}{y}\right)g(y).`$ The distribution $`\mathrm{\Delta }q_iq_i^{}q_i^{}`$ represents the difference between the number densities of quark with helicity parallel to that of parent nucleon and with helicity anti-parallel. The definitions of $`\mathrm{\Delta }\overline{q}_i`$ and $`\mathrm{\Delta }g`$ are the same. $`\mathrm{\Delta }C_q`$ and $`\mathrm{\Delta }C_g`$ are the coefficient functions. In discussing unpolarized reactions, the structure function $`F_2`$ is usually used rather than $`F_1`$, and $`F_2`$ can be written in terms of unpolarized PDFs, $`q_i`$, $`\overline{q}_i`$, and $`g`$, with coefficient functions in the similar way to $`g_1`$. In the next-to-leading-order (NLO) analysis, we choose the modified minimal subtraction ($`\overline{\mathrm{MS}}`$) scheme.
We provide the polarized parton distributions at $`Q^2=1\mathrm{GeV}^2(\mathrm{Q}_0^2)`$. Then, the distributions are evolved from $`Q_0^2`$ to experimental $`Q^2`$ points by DGLAP equations. In our numerical analysis, we use a modified version of the program in Ref. 3, where the evolution equations are solved by a brute-force method.
## 3 Parametrization of polarized parton distributions
Now, we explain how the polarized parton distributions are parametrized. The unpolarized PDFs $`f_i(x,Q_0^2)`$ and polarized PDFs $`\mathrm{\Delta }f_i(x,Q_0^2)`$ are given at the initial scale $`Q_0^2`$. Here, the subscript $`i`$ represents quark flavors and gluon. In our analysis, we require the positivity condition of the PDFs in order to constrain the forms of the polarized PDFs. Therefore, it is convenient to take the following functional form of the polarized PDFs at $`Q_0^2`$:
$$\mathrm{\Delta }f_i(x,Q_0^2)=A_ix^{\alpha _i}(1+\gamma _ix^{\lambda _i})f_i(x,Q_0^2).$$
(3)
The positivity condition is originated in a probabilistic interpretation of the parton densities: the polarized PDFs should satisfy $`|\mathrm{\Delta }f_i(x,Q_0^2)|f_i(x,Q_0^2).`$ In our analysis, we simply require that this condition should be satisfied not only in the leading-order (LO) and but also in the NLO at $`Q_0^2`$. Thus, we have four parameters ($`A_i`$, $`\alpha _i`$, $`\gamma _i`$ and $`\lambda _i`$) for each $`i`$.
In addition to the positivity condition, we assume the SU(3) flavor symmetry for the sea-quark distributions at $`Q_0^2`$ to reduce the number of free parameters. Then, the first moments of $`\mathrm{\Delta }u_v(x)`$ and $`\mathrm{\Delta }d_v(x)`$, which are written as $`\eta _{u_v}`$ and $`\eta _{d_v}`$, can be described in terms of axial charges for octet baryon, $`F`$ and $`D`$, measured in hyperon and neutron $`\beta `$-decays. They are determined as $`F=0.463\pm 0.008`$ and $`D=0.804\pm 0.008`$, which lead to $`\eta _{u_v}=0.926\pm 0.014`$ and $`\eta _{d_v}=0.341\pm 0.018`$. In this way, we fix these first moments, so that two parameters $`A_{u_v}`$ and $`A_{d_v}`$ are determined by these first moments and other parameter values. Thus, the remaining job is to determine the values of the following 14 parameters, $`A_{\overline{q}},A_g,\alpha _i,\gamma _i,\lambda _i(i=u_v,d_v,\overline{q},g)`$, by the $`\chi ^2`$ analysis of the polarized DIS experimental data.
We determine the values of the 14 parameters by fitting the $`A_1(x,Q^2)`$ data for the proton from E130, E143, EMC, SMC, and HERMES, the neutron from E142, E154, and HERMES, and the deuteron from E143, E155, and SMC. We also use LO and NLO GRV parametrizations for the unpolarized PDFs $`^\mathrm{?}`$ and the SLAC measurement of $`R(x,Q^2)`$$`^\mathrm{?}`$ Then, the best parametrization is obtained by minimizing
$$\chi ^2=\frac{(A_1^{\mathrm{data}}(x,Q^2)A_1^{\mathrm{calc}}(x,Q^2))^2}{(\mathrm{\Delta }A_1^{\mathrm{data}}(x,Q^2))^2},$$
(4)
where $`\mathrm{\Delta }A_1^{\mathrm{data}}`$ represents the error on the experimental data including both systematic and statistical errors. In evolving the distribution functions with $`Q^2`$, we neglect charm-quark contributions to $`A_1(x,Q^2)`$ and take the flavor number $`N_f=3`$, because the contribution is very small in a few $`Q^2`$ region where most experimental data exist. To be consistent with the unpolarized, we use the same scale parameters as the GRV, $`^\mathrm{?}`$ $`\mathrm{\Lambda }_{\mathrm{QCD}}^{(3)}=204\mathrm{MeV}`$ at LO and $`\mathrm{\Lambda }_{\mathrm{QCD}}^{(3)}=299\mathrm{MeV}`$ at NLO.
## 4 Results
We got minimum $`\chi ^2`$: $`\chi ^2`$/d.o.f=322.6/360 for the LO and $`\chi ^2`$/d.o.f=300.4/360 for the NLO. The difference between the LO and NLO $`\chi ^2`$ values is about 7%, which indicates the importance the NLO analysis.
We show the $`Q^2`$ dependence of spin asymmetry $`A_1`$ for the proton in Fig. 1. This figure indicates that the NLO effects become larger in the small $`Q^2`$ region and that there is strong $`Q^2`$ dependence especially in the small $`Q^2`$ region. It is not right to assume $`Q^2`$ independence of the spin asymmetry $`A_1`$ in obtaining $`g_1`$, so that, we have to be careful using PQCD in this region.
Next, we show the polarized parton distributions for the NLO in Fig. 2. As the figure indicates, we obtain negative polarization for $`\mathrm{\Delta }\overline{q}(x)`$, and large positive polarization for $`\mathrm{\Delta }g(x)`$. The first moment for $`\mathrm{\Delta }u_v(x)`$ is fixed at the positive value and the one for $`\mathrm{\Delta }d_v(x)`$ is fixed at the negative value, so that the obtained distributions $`\mathrm{\Delta }u_v(x)`$ and $`\mathrm{\Delta }d_v(x)`$ become positive and negative, respectively. Similar results are obtained in the LO distributions. $`^\mathrm{?}`$ There are slight differences between the LO and NLO distributions $`\mathrm{\Delta }d_v(x)`$. However, the differences are large between the LO and NLO gluon distributions in the wide $`x`$ region. It is caused by the gluon contribution through the coefficient function. We calculate the quark spin content by using the obtained LO and NLO distributions. It is given by $`\mathrm{\Delta }\mathrm{\Sigma }=\eta _{u_v}+\eta _{d_v}+6\eta _{\overline{q}}`$, where $`\eta _{\overline{q}}`$ is the first moment of $`\mathrm{\Delta }\overline{q}(x)`$. Because $`\eta _{u_v}`$ and $`\eta _{d_v}`$ are fixed, only $`\eta _{\overline{q}}`$ affects the spin content in the different analyses. The LO and NLO moments are $`\eta _{\overline{q}}=0.064`$ and $`0.089`$, so that the spin content becomes $`\mathrm{\Delta }\mathrm{\Sigma }=0.201`$ and $`0.051`$, respectively. $`^\mathrm{?}`$ The NLO spin content ($`\mathrm{\Delta }\mathrm{\Sigma }=0.051`$) is significantly smaller than other analysis results. For example, the recent SMC and Leader-Sidrov-Stamenov (LSS) parametrizations $`^{\mathrm{?},\mathrm{?}}`$ obtained $`\mathrm{\Delta }\mathrm{\Sigma }=`$0.19 and 0.28 at $`Q^2`$=1 GeV<sup>2</sup>. In order to investigate the reason for the small $`\mathrm{\Delta }\mathrm{\Sigma }`$ in our analysis, we show each antiquark distribution in Fig. 3.
The NLO antiquark distributions of the SMC, LSS, and AAC analyses are calculated at $`Q^2=`$1 GeV<sup>2</sup>. Because the antiquark distribution is not directly given in the SMC analysis, we may call it as a transformed SMC (“SMC”) distribution. It is calculated by transforming the published distributions by the SMC. All the distributions agree in principle in this region ($`0.01<x<0.1`$) where accurate experimental data exist and the antiquark distribution plays an important role. On the other hand, it is clear that our distribution does not fall off rapidly as $`x0`$ in comparison with the others. This is the reason why our NLO spin content is significantly smaller. In fact, we obtained the parameter $`\alpha `$ for the antiquark distribution as $`\alpha _{\overline{q}}(NLO)=0.32\pm 0.22`$, which controlled the small-$`x`$ behavior of $`\mathrm{\Delta }\overline{q}(x)`$. However, the large error of the parameter $`\alpha _{\overline{q}}`$ suggests that the small-$`x`$ part of $`\mathrm{\Delta }\overline{q}(x)`$ cannot be fixed by the existing data. Actually, there is no data in the small-$`x`$ region $`(x<0.04)`$. Therefore, we had better consider to constrain the parameter $`\alpha _{\overline{q}}`$ by theoretical ideas. We discuss such possibilities by using the Regge theory and the perturbative QCD.
According to the Regge model, the small-$`x`$ behavior of $`g_1`$ is suggested as $`g_1(x)x^\alpha `$ with $`\alpha =0.50.0`$$`^\mathrm{?}`$ Therefore, we expect $`\mathrm{\Delta }\overline{q}(x)x^{0.00.5}`$ as $`x0`$. Because the parametrized function is given by $`\mathrm{\Delta }\overline{q}(x)/\overline{q}(x)`$, we should find out the small-$`x`$ behavior of the unpolarized distribution. The GRV distribution has the property $`x\overline{q}x^{0.14}`$ at $`Q^2`$=1 GeV<sup>2</sup> according to our numerical analysis, the Regge prediction becomes $`\alpha _{\overline{q}}^{Regge}=1.11.6`$, if the theory is applied at $`Q^2=1`$ GeV<sup>2</sup>. The perturbative QCD could also suggest the small-$`x`$ behavior. If we can assume that the singlet-quark and gluon distributions are constants as $`x0`$ at certain $`Q^2`$ ($`Q_1^2`$), their singular behavior is predicted from the evolution equations. $`^\mathrm{?}`$ The singlet distribution behaves like $`\mathrm{\Delta }\mathrm{\Sigma }(x)x^\alpha `$ as $`\alpha =0.120.09`$, if we choose the evolution range from $`Q_1^2=0.3\mathrm{\hspace{0.17em}0.5}`$ GeV<sup>2</sup> to $`Q^2=1`$ GeV<sup>2</sup>. Therefore, the perturbative QCD with the assumption of the above $`Q_1^2`$ range suggests $`\alpha _{\overline{q}}^{pQCD}=1.0`$.
The Regge theory and perturbative QCD suggest the range $`\alpha _{\overline{q}}=1.01.6`$, so that we try the NLO $`\chi ^2`$ analyses by fixing the parameter at $`\alpha _{\overline{q}}=`$0.5, 1.0, and 1.6. The last two values are in the theoretical prediction range, and the first one is simply taken as a slightly singular distribution. The obtained minimum $`\chi ^2`$ values are larger than the NLO fit ($`\chi ^2`$=300.4) by 0.1, 1.8, and 7.7%, and the first moments are $`\mathrm{\Delta }\mathrm{\Sigma }`$=0.123, 0.241 and 0.276 for $`\alpha _{\overline{q}}=`$0.5, 1.0, and 1.6, respectively. The small-$`x`$ falloff for larger $`\alpha _{\overline{q}}`$ changes the $`\eta _{\overline{q}}`$ and $`\mathrm{\Delta }\mathrm{\Sigma }`$ significantly. We show the spin content in the region between $`x_{min}`$ and 1 by calculating $`\mathrm{\Delta }\mathrm{\Sigma }=_{x_{min}}^1\mathrm{\Delta }\mathrm{\Sigma }(x)𝑑x`$ in Fig. 4. Because the LSS and SMC distributions are less singular functions of $`x`$, their spin contents saturate even at $`x_{min}=10^4`$ although the $`\mathrm{\Delta }\mathrm{\Sigma }`$ of our NLO result with free $`\alpha _{\overline{q}}`$ still decreases in this region. If the parameter $`\alpha _{\overline{q}}`$ is taken in the perturbative QCD and Regge theory prediction range, the calculated spin content is within the usually quoted values $`\mathrm{\Delta }\mathrm{\Sigma }=0.10.3`$. In this sense, our results are not inconsistent with the previous analyses. Our results indicate that the spin content cannot be determined uniquely, because the accurate experimental data are not available in small $`x`$ region. The obtained $`\chi ^2`$ value suggests that the $`\alpha _{\overline{q}}`$=1.0 solution could be also taken as one of the good fits. The $`\alpha _{\overline{q}}`$=0.5 distributions are almost the same as the ones in the free-$`\alpha _{\overline{q}}`$ NLO analysis, so that it is redundant to propose it as one of our good fits.
From our analyses, we propose the LO distributions, the NLO ones with free $`\alpha _{\overline{q}}`$ (NLO-1), and those with fixed $`\alpha _{\overline{q}}`$=1.0 (NLO-2) as the longitudinally-polarized parton distributions of the AAC analyses. Useful functional forms are given at $`Q^2`$=1 GeV<sup>2</sup> in Appendix B of Ref. 1 for practical applications.
## 5 Conclusions
From the LO and NLO $`\chi ^2`$ analyses, we obtained good fits to the experimental data. Because the NLO $`\chi ^2`$ is significantly smaller than that of LO, the NLO analysis should be necessarily used in the parametrization studies. It is particularly important for extracting information on $`\mathrm{\Delta }g`$. However, the polarized antiquark and gluon distribution cannot be uniquely determined by the present DIS data. We provide the optimum LO and NLO distributions at $`Q^2=1\mathrm{GeV}^2`$ from our numerical analyses.
## Acknowledgments
M.H. was supported as a Research Fellow of the Japan Society for the Promotion of Science, and he would like to thank S. Kumano and M. Miyama for reading this manuscript. This talk is based on the work with Y. Goto, N. Hayashi, M. Hirai, H. Horikawa, S. Kumano, M. Miyama, T. Morii, N. Saito, T.-A. Shibata, E. Taniguchi, and T. Yamanishi. $`^\mathrm{?}`$
## References |
warning/0001/astro-ph0001153.html | ar5iv | text | # RESOLVING THE JEANS MASS IN HYDRODYNAMIC SIMULATIONS OF HIERARCHICAL CLUSTERING
## 1 INTRODUCTION
Numerical simulations have proved useful for understanding the history and dynamics of cosmological structure (Davis et al. 1992; Katz 1992; Cen et al. 1994; Evrard et al. 1996). Our ability to model collisionless systems has matured to the extent that bulk properties of the dark matter can be predicted with increasing confidence, although important exceptions exist, such as the merger rates of substructure as well as the inner profiles of dark matter halos (Moore et al. 1998, see, for example,).
Combining collisionless N-body techniques with a hydrodynamical method to follow the baryonic component opens the way to a detailed understanding of cosmic structure formation in a variety of models. Modeling the baryonic component removes the need to make assumptions about the relationship between the baryonic material and the underlying collisionless matter. Such techniques have proved to be both popular and powerful (Hernquist & Katz 1989; Evrard 1990; Katz & Gunn 1991; Cen 1992; Navarro & White 1993; Steinmetz & Müller 1993; Ryu et al. 1993; Bryan et al. 1994; Gnedin 1995; Steinmetz 1996). A recent comparison of several hydrodynamic codes (Frenk et al. 1998) has shown that we may also begin to have confidence in the overall properties of the baryonic distribution predicted by these codes. However, several differences which are apparent both between codes and between results of varying resolution suggest that convergence has yet to be established conclusively. We address here the particular issue of convergence with resolution in hydrodynamic simulations in which objects are built up through hierarchical merging. We do not consider radiative cooling in this work.
In the standard model of cosmological structure formation, larger structures are formed via the amalgamation of smaller structures formed at earlier times in a hierarchical fashion. In any simulation which models the formation of such structures in this hierarchical manner, there will be an abundance of objects formed at the resolution limit of the simulation. Numerical effects will be a particular problem for these objects. The propagation of these effects into the larger objects subsequently formed requires characterization both to determine the reliability of the physical model as well as to ensure convergence among simulations performed at various resolutions.
Recently Owen & Villumsen (1997) has argued that it is necessary to resolve the Jeans mass in a simulation if convergence is to be attained with increasing resolution. This requires that a minimum temperature be imposed in the simulation such that the gas in the first bound objects—which form at the resolution limit in a standard simulation of hierarchical clustering—is at a sufficient temperature to support itself via pressure in the gravitational potential well of the first dark matter structures. They argue that if the gas temperature is initially much less than the virial temperature of the first halos (this is a boundary condition frequently adopted in cosmological simulations), the cold gas falling into the first potential wells will suffer spuriously large amounts of shock heating of the gas. This “shocking” will be greater for simulations with lower resolution since the scale over which the shocks occur will be larger. Conversely, if the initial generation of objects is pressure-supported, the interaction of the baryons in merging halos as the hierarchical structure formation proceeds will lead only to weak shocks. (This result follows since the velocity with which lumps merge is similar to the internal velocity dispersion of the individual lumps in a typical hierarchy.) Owen and Villumsen suggests that the difference in the amount of shock heating at early stages will persist through to the final object and prevent convergence of numerical results unless a minimum temperature is imposed.
Owen & Villumsen finds, in scale-independent simulations, that without the imposition of a minimum temperature to pressure support resolution-limited structures, the gas is found to move en masse toward a state of higher density and lower temperature as the resolution of the simulation increases. This behavior is not to be expected during the formation of hydrostatic structures. In the standard model, we expect the virial temperature of a structure to be broadly determined by its mass, following the relationship $`TM^{2/3}`$. Since the bulk of the gas resides in the halos of these structures, the temperature of the bulk of the gas should consequently be determined by the dark matter distribution, a distribution whose convergence among differing resolutions is confirmed in the simulations of Owen & Villumsen.
There are other issues of concern raised by the results of the simulations described in Owen & Villumsen (1997). First, the results suggest paradoxically that the injection of thermal energy into a system at an early time can lead to a cooler system at a later time. The authors claim that this is achieved by reducing the amount of shock-heating as the first, unresolved, objects collapse. However, the injection of thermal energy into the system via the imposition of a minimum temperature pre-heats the gas in a manner which is essentially equivalent to shock-heating the gas in the first objects.
Another concern is the concentration of baryons relative to the dark matter in the cores of structures reported in Owen & Villumsen. In three-dimensional simulations without cooling, spanning a wide range of resolutions, others have found extended distributions of baryonic material, relative to the dark matter (Evrard 1990; Thomas & Couchman 1992; Cen & Ostriker 1993; Kang et al. 1994; Metzler & Evrard 1994; Pearce et al. 1994; Navarro et al. 1995; Anninos & Norman 1996; Lubin et al. 1996; Pildis et al. 1996). Pearce et al. (1994) explains the phenomenon as a result of the merging process in which gas is shocked, permanently removing energy from the dark matter component and passing it to the gas.
A possible explanation for the conflicting results regarding the baryon concentration may be that the simulations of Owen & Villumsen were done using a 2-D code. This allowed them to achieve high spatial resolutions but may have introduced phenomena unique to two-dimensional geometries.
### 1.1 Jeans Mass in a Two-component Fluid
In a self-gravitating baryonic fluid, the Jeans mass is easily obtained, by a perturbatitive analysis, as a sharp threshold between acoustic oscillation and collapse, depending on the density and temperature of the fluid. In a two-component fluid comprised of baryons and collisionless cold dark matter, the situation is more complicated. The low initial velocity dispersion of cold dark matter leads to instability, or growth, in this component on all scales. The relevant question in this case for the baryonic fluid is on what scales do the two components remain coupled adiabatically? On small scales, gas pressure will be sufficient to oppose the gas falling into the small dark matter potential wells, whereas on large scales the initial gas temperature will be lower than the effective virial temperature of the forming dark matter halo and the gas will increase in overdensity with the dark matter. A perturbative analysis of the coupled two-component fluid is relatively straightforward but not very illuminating for our purposes. In particular it can be shown that there is not a single Jeans mass in the sense described above (see Peacock 1999, for example). It is sufficient for our purposes to define the Jeans mass as that which is just sufficient to adiabatically compress the gas to an overdensity which is the same as that of the dark matter.
With this assumption we can approximate the Jeans mass of the gas for a fluid composed of collisionless dark matter of density, $`\rho _{DM}`$, and collisional baryonic material (gas) of density, $`\rho _g`$, and temperature, $`T`$, by
$$M_J\frac{\rho _g}{2}\left(\frac{1}{G\rho _M}\frac{kT}{\mu m_H}\right)^{{}_{}{}^{3}/_{2}^{}},$$
(1)
where the total mass density is given by $`\rho _M=\rho _{DM}+\rho _g`$.
### 1.2 The Minimum Temperature for Pressure Support
The effective mass resolution of a simulation is usually fixed. For simulations such as those described here which use smoothed particle hydrodynamics (SPH) to model the gas forces, this can be taken to be the total mass of the number of particles, $`N_{SPH}`$, over which local quantities are averaged, presuming all gas particles have equal mass. The number of particles is typically on the order of 30. By associating the mass of this number of particles to the Jeans mass (Eq. 1) and assuming a total mass density, $`\rho _M`$ in Eq. 1, equal to the mean density of the simulation volume, $`\overline{\rho }_M`$, the temperature required to support gas in dark matter halos of this size is
$$T_{min}=\left(\frac{2N_{SPH}m_g}{\overline{\rho }_g}\right)^{{}_{}{}^{2}/_{3}^{}}G\overline{\rho }_M\frac{\mu m_H}{k},$$
(2)
where $`m_g`$ is the mass per gas particle. If the gas has a minimum temperature set by this value, then we can say that the simulation is resolving the Jeans mass. This assumption is flawed in a number of ways. Since local density is not in general equal to the mean in the simulation, the minimum temperature should also be a function of position. However, the minimum temperature should vary with the local density as $`T_{min}\rho _M^{{}_{}{}^{1}/_{3}^{}}`$, which is not a particularly sensitive function, especially in those regions in which the first objects are forming. More significant is the decoupling of the components of the fluid, discussed previously. This is likely to occur for the first objects formed since the dark matter halos collapse regardless of whether or not the gas is pressure-supported. For those situations in which, locally, $`\rho _{DM}\rho _g`$, enhancement of the dark matter density without a similar increase in the gas density requires the temperature to rise by an equivalent fraction in order to maintain a constant Jeans mass for the gas. However, during the formation of these dark matter structures, the gas will be allowed to gradually respond, adiabatically, to the increasing gravitational potential, partially or fully offsetting the need for an increased minimum temperature. This would not occur without the initial minimum temperature since the gas in that case would initially collapse along with the dark matter.
### 1.3 The Hypothesis
In this paper, we propose that the effect on the state of the gas of poorly resolving the first baryonic structures does not propagate into larger structures formed through hierarchical merger. The properties of the gas in the larger halos will be determined primarily by the virial properties of the structures which are dominated by the distribution of the collisionless component.
To test this hypothesis, we performed a series of seven simulations of hierarchical structure formation. The simulations were done at three different resolutions. For two of these simulations, the temperature of the gas content was prevented from being less than the temperature required to maintain pressure support for the smallest resolvable object; that is, the Jeans mass was resolved.
Comparison of the data simulated for this study was made using both the total sample of gas in each simulation volume as well as an individual well-resolved structure. Using an individual cluster permits examination of those effects that are frequently the focus of interest in structure formation simulations; the void particles in simulations of this sort are usually ignored due to a lack of resolution in this regime.
We examine the entropy of the gas, which is a measure of the degree of shocking. Also examined is the density distribution of the gas to determine the extent of any bulk movements of the gas to higher densities as resolution is increased. We also explore mean bulk properties of the structures: in particular, the mean dark matter and gas density profiles and the mean baryon fraction profile. The mean gas density and baryon fraction profiles are potentially sensitive to a lack of convergence as well as being of fundamental interest.
The outline of the paper is as follows. The simulations are described in Sec. 2. The results of an attempt to recreate the results of Owen & Villumsen but using a three dimensional geometry are given in Sec. 3.1. The impact of resolving the Jeans mass on the final state of the gas in galaxy clusters is addressed in Sec. 3.2. The results are discussed in Sec. 4. Hereafter, we refer to Owen & Villumsen (1997) as OV97.
Though the simulations are scale invariant, the results on occasion are given in physical units. This is simply a matter of convenience for the purpose of comparison.
## 2 THE SIMULATIONS
The N-body code with SPH, hydra (Couchman et al. 1995), was used for all simulations. hydra calculates the gravitational forces using an adaptive mesh for the large-scale complemented with particle-particle calculations for the short range forces. Smoothed particle hydrodynamics (SPH) is used for the hydrodynamics.
Seven simulations were performed. Five used the original code, and two used a variant of hydra which imposes a minimum temperature on the gas particles. The details of the simulations are outlined in Table 1. Given, there, are the number of each type of particle and the gravitational softening length, $`ϵ`$. The simulations that were evolved to the final expansion factor, $`a=a_f`$, were performed twice, once with and once without a minimum temperature, whereas the first trio were evolved only once, without a minimum temperature. Initial density perturbations are the same for all resolutions within a group, albeit limited by the Nyquist frequency.
The minimum temperature is set to support a mass of particles corresponding to the effective mass resolution of the simulation as described in the introduction. At each epoch of the simulation, structures composed of $`N_{SPH}`$ gas particles with a local density equal to the mean gas density of the volume, $`\overline{\rho }_g`$, and individual masses of $`m_g`$ are supported against collapse by the imposition of a minimum temperature, $`T_{min}`$. The minimum temperature is derived by simply equating the Jeans mass, $`M_J`$, to the structure mass $`N_{SPH}m_g`$, as is done in Sec. 1.2 to derive Eq. 2.
During the two simulations in which a minimum temperature was imposed, $`T_{min}`$ was recalculated at each epoch. Since the minimum temperature scales with the mean density as $`T_{min}\overline{\rho }^{1/3}`$, then it scales with the expansion factor, $`a`$, as $`T_{min}a^1`$. Normally, adiabatic cooling for a uniform monotonic gas would maintain the relation, $`Ta^2`$. Consequently, the imposition of the minimum temperature has the effect of continuously injecting thermal energy into the fraction of the gas whose adiabatic expansion would bring its temperature below the minimum temperature.
The initial conditions for the two-dimensional simulations of OV97 have a density perturbation spectrum with a spectral index equivalent to the three-dimensional $`n=1`$ form used in the simulations presented here. Because the simulations are scale free, the spatial variance of the density fluctuations at the end of the simulation in the larger volume corresponds to an earlier epoch in the smaller volume, after scaling by the box size.
The amplitude of linear density fluctuations, as measured by the power spectrum of the density field, $`P(k)`$, relates to the wavenumber of the fluctuation, $`k`$, measured in comoving coordinates, and the time, $`t`$, via
$$P(k)\left(\frac{\sigma _L}{h}\right)^2k^nt^{{}_{}{}^{4}/_{3}^{}}$$
(3)
where $`\sigma _L`$ is a normalization factor corresponding to the assumed value of the variance, $`\sigma `$, of the present density field on the scale of $`L`$.
Although the simulations discussed in OV97 and those presented here are scale invariant, it is useful to set a physical size to the simulation volume. Indeed, this may be a benefit to those readers familiar with the scale of cosmological structures. Further, this is the only practical way we have of determining the degree of clustering in the simulations of OV97. OV97 uses a volume with a side of length $`128h^1Mpc`$, $`h=0.5`$ (Owen & Villumsen 1996). Since we use a volume with sides $`40h^1Mpc`$ and $`h=0.65`$, this gives a ratio of scales of $`4.16`$ for their simulation scale compared to the simulations described here. For the normalization factor, they have $`\sigma _8=0.6`$ whereas we used $`\sigma _8=0.935`$. This implies an expansion factor of $`0.24a_f`$. That is, the structures in the simulations presented here, scaled by the box size, will have grown in similar magnitude to those of OV97 after an expansion of $`0.24`$ of the final expansion factor of the simulation.
In this study, the data have been compared at the slightly earlier expansion factors of $`0.17a_f`$ and $`0.22a_f`$ to ensure that structures have not undergone more than a couple of merger events. As a check, a volume with the same cosmological parameters as described in OV97 but with the same perturbation waves as found in our smaller boxes was evolved to an expansion factor, $`a_f`$. The data from the the $`40h^1Mpc`$ volumes at an expansion factor of $`0.22a_f`$ are visibly similar to the final output of this test volume, supporting the legitimacy of scaling the numerical results.
It was feasible, computationally, to evolve the volume to the present only at the resolutions of $`2\times 64^3`$ particles and less. Hence, this study will concentrate on a comparison of the simulations of resolution corresponding to this number of particles with simulations of a resolution appropriate to $`2\times 32^3`$ particles. These resolutions will be referred to henceforth as the $`64^3`$ and $`32^3`$ simulations. Each was evolved with and without a minimum temperature. For the earlier epoch corresponding to the expansion factor of $`0.17a_f`$, the volume could also be evolved at the resolution provided by $`2\times 128^3`$ particles. Thus a set of three simulations was performed using the same density perturbations (limited by the Nyquist frequency) at resolutions corresponding to $`128^3`$, $`64^3`$, and $`32^3`$.
Cooling was neglected to maintain the scale invariance. In all simulations, $`\mathrm{\Omega }=1`$, $`\mathrm{\Omega }_{DM}=0.9`$, and $`\mathrm{\Omega }_g=0.1`$.
The initial density perturbations were established by displacing the particle positions from a uniform cubic grid using the Zel’dovich (1970) approximation for growth of density perturbations in the linear regime in the standard way.
The initial redshift for the $`32^3`$ and $`64^3`$ simulations is $`z_{init}=75`$. This was chosen to keep the maximum displacement incurred during the establishment of the density field to less than $`{}_{}{}^{1}/_{2}^{}`$ of the initial grid spacing. This ensures that the correct non-linear field will be attained at later times. The same criterion was used to set the initial redshift for the $`128^3`$ simulation for which $`z_{init}=150`$.
## 3 RESULTS
### 3.1 The State of the Baryons at an Early Epoch
We first attempt to reproduce the results of OV97 by examining the state of the baryons in our simulations at a comparable epoch with that of the final epoch of the simulations of OV97. The smaller physical scale of the simulations presented here requires that we examine our data before the nominal end of the simulation in order that a comparable level of clustering has developed to that present at the final epoch of the simulations of OV97. The required expansion factor is $`a=0.22a_f`$, where $`a_f`$ is the final expansion factor of the simulations.
Following the analysis given by OV97, we look at the distribution of the gas particles in the $`\rho T`$ plane as well as the density distributions for which the density is resampled down to the resolution of the lowest resolution runs as described below. The position within the $`\rho T`$ planes gives an indication of the degree of shocking, with the more highly shocked gas (gas with greater entropy as indicated by $`T\rho ^\gamma `$, $`\gamma ={}_{}{}^{2}/_{3}^{}`$) lying to the upper left (high temperature, low density).
In order to allow a direct comparison of the densities in simulations of different resolution a density resampling technique was employed. The (resolution-degrading) resampling was accomplished by smoothing over equal masses for each particle using an SPH-like density estimate. Consequently, there are approximately eight times as many particles in each estimate for the $`64^3`$ data as in the $`32^3`$ data, but the smoothing radii are approximately the same. This resolution degradation is performed for the high resolution simulations ($`64^3`$) when comparison is made with the lowest resolution simulation ($`32^3`$). This technique was unfortunately not possible with the $`128^3`$ data-set due to the computational effort of summing over the more than 3000 neighbor particles of each gas particle.
It is clear that for density profiles without cores, higher mass and spatial resolutions will lead to more gas at a higher density even if the simulations can be said to have converged to the true solution within resolution limits. If the densities are averaged within volumes encompassing a mass which is resolved in the lowest resolution simulation, then these masses should not vary significantly.
The simulations of OV97 produce more cool, dense gas as the resolution increases. They interpret this to mean that in the higher resolution simulations, the gas in the cores of structures is not as strongly shocked. To search for this effect, three sets of data were evolved at the resolutions provided by $`2\times 32^3`$, $`2\times 64^3`$, and $`2\times 128^3`$ particles, half in each phase of gas and dark matter, as described in Sec. 2. Each set has the same initial perturbations, limited appropriately in the high spatial frequency domain by the Nyquist frequency. These data sets were evolved to an early epoch for which the amplitude of density fluctuations has grown to a level comparable to those in the data presented by OV97.
The evolved data were resampled, in the case of the $`64^3`$ and $`128^3`$ sets, to the same mass resolution as the $`32^3`$ set by randomly selecting $`32^3`$ particles from each phase. This avoids the problem of performing the density estimate over the very large number of neighbor particles required for the $`128^3`$ data set, at the expense of an increase in shot noise. Local SPH density estimates for the sampled particles in these resampled sets were determined allowing comparisons of the data in $`\rho `$-$`T`$ space.
The $`\rho `$-$`T`$ distributions for all the gas in each of these three simulations are illustrated in 1 (top panels). This is comparable to the top panels of Fig. 8 of OV97. It is clear that there is no trend for the bulk of the mass to move toward a cooler and denser state as the resolution increases, in contrast to the results of OV97. Since the bulk of this gas is in the halos of bound structures, the argument given in the introduction is supported; the entropy of the gas is set by the virial properties of the halo which are dictated by the dark matter. There is a trend for the void particles (found to the lower left in the $`\rho `$-$`T`$ plane) to increase in temperature as the resolution increases. There is also a trend for the distributions to become more diffuse, particularly toward the denser and cooler state.
The lower set of panels in 1 give the distribution for particles in a cluster common to all three sets of data. The cluster was selected simply for being one of the larger clusters of particles to have formed by this early stage. The cluster is barely resolved ($`100`$ particles) in the $`32^3`$ simulation, but is well resolved in the $`128^3`$ simulation. Here it is evident that the diffusion of the distributions to a denser and cooler state is due to a small population of particles, and not a bulk movement. Upon examination of the spatial distribution of this population, these particles are found to be located in substructure that is present only in the simulations with greater mass resolution. These clumps are not as shocked as the particles in the halos of the substructure. This is consistent with the phenomenon described by OV97 in nature, though not in extent since it is not a bulk motion, but simply a new population of particles.
As the mass resolution increases, the first bound haloes form at earlier epochs at higher densities leading to increased central densities in the final dark matter halos. This is illustrated in 2 which compares the gas and dark matter mean density profiles for resolved clusters in the simulations, albeit at a later epoch. The dark matter density is double in the $`64^3`$ simulation compard with the $`32^3`$ simulation in the central regions (close to the resolution limit of the $`32^3`$ simulation). This increases the degree of shocking of the halo particles more than offsetting the decreased amount of shocking in the cores of substructure. The mean temperature of this structure thus increases with increasing resolution: in physical units it is $`2.9\times 10^7K`$, $`3.2\times 10^7K`$, $`4.5\times 10^7K`$ for the $`32^3`$, $`64^3`$, and $`128^3`$ simulations, respectively. The total amount of shocking on the whole also increases, but not to such a degree. The amount of shocking, as measured by the mass-weighted mean of the entropy parameter, $`<T/\rho _g^{{}_{}{}^{2}/_{3}^{}}>_M`$, for an individual cluster at the three resolutions evolved to $`a=0.17a_f`$ is given in Table 2. The values not resolution-degraded are given in the second column which reveal a slight increase (10% for every doubling in resolution) in the entropy as the resolution increases. If the densities are recalculated for the resampled data with the intent of comparing similar effective resolutions, then the degree of shocking decreases quite markedly with increasing resolution (Table 2, third column). This occurs simply because the temperatures of the particles are not changing while the resampled densities of the particles in the densest regions decrease due to the use of a larger smoothing radius. It would be erroneous to conclude that the gas is actually less shocked, since the temperature of the halo gas particles and the depth of the potential well are so intimately related.
The gas density distributions for both the $`32^3`$ and $`64^3`$ data were found at an expansion factor of $`0.22a_f`$. The densities of the $`64^3`$ data-set were resampled as described previously. The distribution of gas density at this early epoch (3, left hand panels, dashed lines) is only weakly dependent on resolution, becoming somewhat less peaked at $`10^{28}gcm^3`$ at the higher resolutions. No bulk motion of the gas to higher densities is observed, unlike the simulations described in OV97.
Since no bulk motion of the gas to higher densities and lower degrees of shocking is observed, we might predict that the imposition of a minimum temperature will not significantly affect these results. The distribution of the gas within the $`\rho T`$ plane at the early epoch of $`a=0.22a_f`$ is illustrated in 4. Again, the bulk of the gas is found in the same state in both simulations without a minimum temperature (4, left panels), with a diffusion toward a state of lower entropy. For those simulations in which a minimum temperature is imposed (4, right panels), the gas is displaced within the $`\rho T`$ plane in those regions below $`T_{min}`$ (for obvious reasons) and below the locus that uniform gas would trace if it were allowed to cool to the minimum temperature appropriate at each epoch. In order for gas to be found below this locus it must be initially held at a high temperature and density but with little or no increase in entropy and subsequently allowed to cool adiabatically. The bulk of the gas is above these boundaries. The gas that is above these boundaries is consistently located in similar proportions to the gas found in the same locations in the $`\rho T`$ plane in the simulations without a minimum temperature. There is no gas in the lower resolution simulations which fits these criteria, however.
The resampled density distributions for the simulations with a minimum temperature, illustrated in 3 (right panels, dashed lines), confirm that the minimum temperature does modify the distributions but does not provide convergence. This figure is comparable to Fig. 6d in OV97 which indicates a bulk motion to higher densities as the resolution is increased, even when the data are resampled to the scale of the lowest resolution simulation. In our simulations, for the gas denser than $`10^{27}gcm^3`$, there is actually a decrease in the agreement between the lower and higher resolution data when a minimum temperature is enforced.
The decrease in the degree of shocking reported by OV97 as the resolution of a hydrodynamic simulation is increased is confirmed in the cores of clusters of particles. However, the bulk movement of the gas particles to a less shocked state is not observed in our simulations. Indeed, the mean degree of shocking of cluster particles is essentially conserved, with the mean temperature of particles actually increasing with resolution. The displacement in density distributions reported by OV97 is also not observed, nor is any sort of convergence provided by the imposition of a minimum temperature.
### 3.2 The State of the Baryons at a Late Epoch
After many mergers, it is possible that the gas in the initial structures described above will have ‘forgotten’ the degree of shocking it had at early times. It may have been sufficiently shocked by subsequent mergers that early heating, via strong shocks or otherwise, would no longer be significant to the total thermal content and distribution. To examine this hypothesis, the $`32^3`$ and $`64^3`$ sets of simulations used in Sec. 3.1 were evolved to a further state which involved many levels of merging. It is not presently feasible to evolve a $`128^3`$ to such a state with the hardware on hand in a manageable amount of time.
The gas distributions in the $`\rho `$-$`T`$ plane at this end state are illustrated in 5. At densities less than $`5\times 10^{26}gcm^3`$, the distributions for the $`32^3`$ (top left) and $`64^3`$ (bottom left) simulations are identical in form. At $`5\times 10^{26}gcm^3`$, it is the lower mass-resolution simulation that produces the greatest amount of cool, dense gas.
Evolving these sets of data while enforcing a minimum temperature sufficient to pressure support structures with $`N_{SPH}`$ or fewer particles does little to alter these distributions (5, panels to the right) outside those regions below $`T_{min}`$ and below the line with a slope of $`{}_{}{}^{1}/_{3}^{}`$ which extends from the point of initial density and temperature of the gas. This line corresponds to the path a uniform gas would take due to the imposition of a minimum temperature.
Recall from Sec. 3.1 that individual structures in high resolution simulations have, at early times, less shocked material in their cores than is produced at lower mass resolution. This difference is eliminated after the initial merging period on account of the introduction of new structures falling into the lower-resolution cluster. These structures have material with low degrees of shocking. Compare the lower panels of 1 with 6 (left panels) which plots the $`\rho `$-$`T`$ distribution for a cluster from a $`32^3`$ and $`64^3`$ simulation at late times (note that the figures have different scales). At early times, there is a population of less shocked gas in the high resolution runs as described earlier. This leads to the spread in temperature at high densities. Although the distribution of substructures differs among the simulations of differing resolution, the presence of substructures at the termination of the simulations leads to a similar spread in temperatures, regardless of resolution. Taken as singular objects, the bulk properties of individual clusters do not vary substantially between simulations of differing resolutions. The degree of shocking, as parameterized by the entropy factor $`<T\rho ^\gamma >_M`$ (a mass-weighted mean) is found to be essentially constant for the cluster simulated at either resolution. It is $`(1.12\pm 0.02)\times 10^{26}Kcm^2g^{{}_{}{}^{2}/_{3}^{}}`$ for both the $`32^3`$ and $`64^3`$ simulations. With the imposition of a minimum temperature, the entropy is essentially unchanged at $`(1.13\pm 0.03)\times 10^{26}Kcm^2g^{{}_{}{}^{2}/_{3}^{}}`$. The trend for the mean temperature of structures to increase with increasing resolution is maintained, with the gas temperature rising from $`(2.9\pm 0.1)\times 10^7K`$ to $`(3.2\pm 0.1)\times 10^7K`$ between the $`32^3`$ and $`64^3`$ simulations. This is understood to be a by-product of the deeper potential wells in the cores of structures as described earlier.
The density distribution confirms the similarity of the states of the gas in the simulations of differing resolution (3, solid line). In the absence of a minimum temperature (3, left panels), the gas density distribution varies marginally with resolution with a decrease in the amount of gas at a density of $`10^{29}gcm^3`$ compensated by an increase at the higher density of $`10^{27}gcm^3`$ as the spatial resolution is doubled. The addition of a minimum temperature skews the distribution even more significantly to lower densities for the lower-resolution simulation. The shift in the simulations without a minimum temperature can be explained as a result of the greater abundance of substructure in the dark matter distribution which locally compresses and heats the gas. The gas at the highest densities in the cores of structures ($`\rho >10^{26}gcm^3`$), is not significantly affected.
Finally, a useful diagnostic is the mean density profile (2) of the clusters at the final epoch, $`a=a_f`$. It confirms that the presence of a minimum temperature does not change the morphology of the clusters. The mean profiles are calculated from a sample of clusters selected by their size, to ensure that the profiles are resolved. The inner cores are excluded from the calculation of the means, since they are not resolved.
The simulations of OV97 display a marked tendency for the gas to become enriched in the cores of clusters, illustrated in Fig. 9 of OV97. This is contrary to the results found by others (Evrard 1990; Thomas & Couchman 1992; Cen & Ostriker 1993; Kang et al. 1994; Metzler & Evrard 1994; Pearce et al. 1994; Navarro et al. 1995; Anninos & Norman 1996; Lubin et al. 1996; Pildis et al. 1996). Indeed, it is most commonly found that the gas becomes more dispersed than the dark matter. This result is supported by the simulations we performed. The baryon fraction parameter, $`\mathrm{{\rm Y}}`$, which is the ratio of the baryon density to total mass density normalized by the global mean ratio, is plotted in 7 for the clusters. The profiles given are the mean profiles for the same sample of clusters for which the mean density profiles have been calculated. It is also clear from this figure that the presence of a minimum temperature has no effect on the baryon fraction distribution. It also shows clearly that resolution does have an effect in the cores.
## 4 DISCUSSION
It has been suggested (Owen & Villumsen 1997) that convergence with increasing mass resolution in numerical simulations of hierarchical structure formation requires the imposition of a minimum temperature in order that gas is pressure supported in structures with masses similar to the effective resolution of the simulation. Our results do not support this contention, either at early epochs after only a few mergers, or at later epochs after a significant number of mergers have occurred and the minimum temperature has dropped well below the virial temperature of the typical structures. Our results indicate that the problem of the gas becoming over-shocked in simulations in which the Jeans Mass is not resolved is not as severe as suggested by OV97. The larger degree of shocking of the bulk of the gas noted by OV97 as resolution is lowered is not reproduced at any epoch of the simulations presented here. Consequently, the imposition of a minimum temperature is not necessary to achieve convergence.
We do observe a phenomenon which is similar to the decrease in shocking with increasing resolution as described in OV97. At an early epoch, isolated halos are shown to possess a large amount of substructure in their high resolution simulations that is not present in their low resolution counterparts. This substructure contains cool, dense, weakly shocked gas. However, the presence of this population of cool and dense gas in the high resolution simulations is offset by an increase in the temperature of the halo gas surrounding these substructures. The net effect is not to lower the mean temperature of the halo but to marginally increase it ($`10\%`$ for each doubling in resolution) as well as to increase the entropy of the halo by an equivalent fraction. At a later epoch, after subsequent accumulation of material, the halos in the lower resolution simulations develop a comparable population of cooler, unshocked gas within substructure.
Further attempts were made to reproduce the results of OV97 by modifying our code. Since it was believed that the gravity calculations for the simulations of OV97 were done using a particle-mesh method, we disabled the particle-particle contributions of the gravitational force as well as the refinements in the hydra code. This was unsuccessful in reproducing their results.
The results of imposing a minimum temperature sufficient to pressure-support the resolution-limited gas structures were examined for adverse effects. At the final epoch, the gas density profiles of the clusters were not altered. The degree of shocking increased by less than $`5\%`$. This increase may be attributed primarily to an increase in temperature and decrease in density of the gas lying near the cores of structures. The bulk of the gas, which lies in the halo, was unaffected.
The concentration of baryonic material in the cores of structures reported by OV97 is also not observed. Instead, an anti-bias is found, in accordance with the substantial majority of numerical simulations reported in the literature as noted previously. Pearce et al. (1994) explains the phenomenon as a result of the merging process of gas-dark-matter halo systems in which gas is shocked, permanently removing the energy from the dark matter component. For the gas to be more biased in clusters, without an energy-loss mechanism operating such as cooling, the gas must transfer energy to the dark matter component.
It is unclear what causes the differences between our results and those of OV97. The most striking remaining difference between our investigations and those of OV97, which we have not explored, is the difference in dimensionality of the numerical simulations. OV97 used two-dimensional simulations, ours were three-dimensional. Structures in two dimensions are rods, which have essentially the same cross sections per mass as spheres, but one less degree of freedom in which to avoid collisions. This may strongly affect the efficiency of shock heating, with more head-on collisions per merging event leading to more violently shocked gas. However, naïve interpretation of this observation would suggest that the gas halos should be more extended, and more anti-biased relative to the dark matter, contrary to their results.
We conclude that resolution of the Jeans mass of the smallest structures is consequently not important in studies of the gas in structures after merging has occurred. It may be important for the first objects formed, but for any study other than the statistics of position and number, these first objects should be ignored since they are, by definition, poorly resolved. Clearly, this result may not hold for simulations in which there is cooling and/or star formation feedback, but these will have an effect only in the densest parts of the clusters, and should not affect global properties.
We thank Rob Thacker for helpful discussions. HMPC acknowledges the support of NSERC. |
warning/0001/hep-th0001018.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Since it was shown by Randall and Sundrum (RS) that gravity in the background of non-dilatonic domain wall with the exponentially decreasing warp factor is effectively compactified, some efforts have been made to understand gravitating objects living in such domain wall, e.g., Refs. . We showed that dilatonic domain walls also effectively compactify gravity if the dilaton coupling parameter is sufficiently small. The main motivation to consider dilatonic domain walls was that the consistency of equations of motion requires the cosmological constant term to have dilaton factor in order for the domain wall spacetime to admit charged brane solutions.
In our previous work , we constructed completely localized solutions describing extreme charged branes living in the worldvolume of extreme dilatonic domain walls for the purpose of understanding charged branes in the RS type model. Unexpectedly, it is found out that a charged $`p`$-brane is not effectively compactified to the charged $`p`$-brane in one lower dimensions. It is speculated that this is due to the unusual properties of the Kaluza-Klein (KK) modes of gauge fields and also presumably form fields that the zero mode is not localized on the lower-dimensional hypersurface of the domain wall and the massive modes strongly couple to the fields on the brane . So, following the result of Ref. that the Schwarzschild black hole in four-dimensional world within a domain wall should be regarded as an uncharged black string in five dimensions, we speculated that a charged $`p`$-brane in one lower dimensions might have to be regarded as a charged $`(p+1)`$-brane where one of its longitudinal directions is along the transverse direction of the domain wall.
It is the purpose of this paper to construct a solution describing a non-extreme charged dilatonic $`(p+1)`$-brane in an extreme dilatonic domain wall in $`D`$ dimensions where one of the longitudinal directions of the brane is along the transverse direction of the domain wall and to study its properties in relation to the RS type model. We find that in the case of an uncharged branes, physics of the uncharged $`p`$-brane in one lower dimensions is reproduced by the uncharged $`(p+1)`$-brane in the domain wall. However, when the $`(p+1)`$-brane is charged, the dynamics of a test particle in the background of the charged $`p`$-brane in one lower dimensions is not reproduced. This is due to the fact that generally the transverse (to the domain walls) component of the spacetime metrics of charged branes in the domain walls has non-trivial dependence on the longitudinal coordinates of the domain walls. Note, the original RS model assumes that the perturbations of the domain wall metric should be along the longitudinal directions of the domain wall, only, in order for the lower-dimensional gravity to be reproduced. So, in order for the RS type model to admit wide variety of gravitating objects which reproduce physics in one lower dimensions, one has to somehow modify the model.
The paper is organized as follows. In section 2, we present the solution describing a nonextreme dilatonic $`(p+1)`$-brane intersecting an extreme dilatonic domain wall in $`D`$-dimensions where one of the longitudinal directions of the brane is along the direction transverse to the domain wall. In section 2, we study the dynamics of the probe $`(p+1)`$-brane in such source background, comparing with the dynamics of the probe $`p`$-brane in the background of the $`(D1)`$-dimensional source $`p`$-brane. In section 3, we repeat the same analysis with a test particle. The conclusion is given in section 4.
## 2 Brane-World Solitons
In this section, we discuss the $`D`$-dimensional solution describing a non-extreme dilatonic $`(p+1)`$-brane with the dilaton coupling parameter $`a_{p+1}`$ intersecting extreme dilatonic domain wall with the dilaton coupling parameter $`a`$ such that one of the longitudinal directions of the $`(p+1)`$-brane is along the direction transverse to the domain wall. The configuration is given in the following table.
| | $`t`$ | $`𝐰`$ | $`𝐱`$ | $`y`$ |
| --- | --- | --- | --- | --- |
| brane | $``$ | $``$ | | $``$ |
| domain wall | $``$ | $``$ | $``$ | |
Here, $`t`$ is the time coordinate, $`𝐰=(w_1,\mathrm{},w_p)`$ and $`y`$ are the longitudinal coordinates of the $`(p+1)`$-brane, and $`𝐰`$ and $`𝐱=(x_1,\mathrm{},x_{Dp2})`$ are the longitudinal coordinates of the domain wall. The solution for such configuration solves the equations of motion of the following action:
$$S=\frac{1}{2\kappa _D^2}𝑑x^D\sqrt{g}\left[\frac{4}{D2}(\varphi )^2\frac{1}{2(p+3)!}e^{2a_{p+1}\varphi }F_{p+3}^2+e^{2a\varphi }\mathrm{\Lambda }\right].$$
(1)
The solution has the following form:
$`ds^2`$ $`=`$ $`H^{\frac{4}{(D2)\mathrm{\Delta }}}[H_{p+1}^{\frac{4(Dp4)}{(D2)\mathrm{\Delta }_{p+1}}}(fdt^2+dw_1^2+\mathrm{}+dw_p^2)`$ (3)
$`+H_{p+1}^{\frac{4(p+2)}{(D2)\mathrm{\Delta }_{p+1}}}(f^1dx^2+x^2d\mathrm{\Omega }_{Dp3}^2)]+H^{\frac{4(D1)}{(D2)\mathrm{\Delta }}}H^{\frac{4(Dp4)}{(D2)\mathrm{\Delta }_{p+1}}}_{p+1}dy^2,`$
$`e^{2\varphi }`$ $`=`$ $`H^{\frac{(D2)a}{\mathrm{\Delta }}}H_{p+1}^{\frac{(D2)a_{p+1}}{\mathrm{\Delta }_{p+1}}},`$ (4)
$`A_{tw_1\mathrm{}w_py}`$ $`=`$ $`{\displaystyle \frac{2}{\sqrt{\mathrm{\Delta }_{p+1}}}}{\displaystyle \frac{\mu \mathrm{cosh}\delta _{p+1}\mathrm{sinh}\delta _{p+1}}{x^{Dp4}}}H_{p+1}^1H^{\frac{4}{\mathrm{\Delta }}},`$ (5)
where the harmonic functions $`H_{p+1}`$ and $`H`$ for the $`(p+1)`$-brane and the domain wall, and the non-extremality function $`f`$ are given by
$$H_{p+1}=1+\frac{\mu \mathrm{sinh}^2\delta _{p+1}}{x^{Dp4}},H=1+Q|y|,f=1\frac{\mu }{x^{Dp4}},$$
(6)
and the parameters $`\mathrm{\Delta }`$’s in the solutions are defined as
$`\mathrm{\Delta }_{p+1}`$ $`=`$ $`{\displaystyle \frac{(D2)a_{p+1}^2}{2}}+{\displaystyle \frac{2(p+2)(Dp4)}{D2}},`$ (7)
$`\mathrm{\Delta }`$ $`=`$ $`{\displaystyle \frac{(D2)a^2}{2}}{\displaystyle \frac{2(D1)}{D2}}.`$ (8)
Here, $`x|𝐱|`$ is the radial coordinate of the transverse space of the $`(p+1)`$-brane, $`\mu >0`$ is the non-extremality parameter, and $`Q`$ is related to the cosmological constant $`\mathrm{\Lambda }`$ as $`\mathrm{\Lambda }=2Q^2/\mathrm{\Delta }`$. The extreme limit of the $`(p+1)`$-brane is achieved by taking $`\mu 0`$ such that $`\mu e^{2\delta _{p+1}}`$ is a non-zero constant. The consistency of the equations of motion requires that the dilaton coupling parameters satisfy the following constraint:
$$aa_{p+1}=\frac{4(Dp4)}{(D2)^2}.$$
(9)
Note, this constraint (and the intersection rules arising from this constraint) is different from the ordinary one which is satisfied by the intersecting branes whose harmonic functions depend on the overall transverse coordinates. The configuration under consideration in this paper does not have an overall transverse direction and each constituent is localized along the relative transverse directions. So, the solution (5) can be regarded as a particular case of general semi-localized solutions describing intersecting branes in which each constituent is localized along the relative transverse directions and delocalized along the overall transverse directions. Such semi-localized intersecting brane solutions are constructed in Ref. for the extreme case and in Ref. for the non-extreme case, and satisfy different intersection rules <sup>2</sup><sup>2</sup>2No-force requirement on the probe brane in the source brane background yields the ordinary intersection rules, only. So, the constraint (9) on the dilaton coupling parameters is different from the one we expected through the no-force requirement in our previous work .. An example is intersecting two NS5-branes with one-dimensional intersection , rather than the three-dimensional one. It is interesting to note that when $`p=D4`$ there is no constraint on one of the dilaton coupling parameters if the other one is zero. So, in this case, the bulk background of non-dilatonic domain walls ($`a=0`$) can admit charged $`(p+1)`$-branes with an arbitrary dilaton coupling parameter $`a_{p+1}`$. Also, in this case, the bulk background of both dilatonic ($`a0`$) and non-dilatonic ($`a=0`$) domain walls can admit non-dilatonic charged $`(p+1)`$-branes ($`a_{p+1}=0`$). However, in this case the harmonic function $`H_{p+1}`$ is logarithmic.
## 3 Probing Brane-World Solitons
In this section, we repeat the probe dynamics analysis of our previous work with the source background (5) of a brane-world soliton, comparing with the probe dynamics in the following source background of a $`(D1)`$-dimensional non-extreme dilatonic $`p`$-brane:
$`ds^2`$ $`=`$ $`H_p^{\frac{4(Dp4)}{(D3)\mathrm{\Delta }_p}}\left[fdt^2+dw_1^2+\mathrm{}+dw_p^2\right]+H_p^{\frac{4(p+1)}{(D3)\mathrm{\Delta }_p}}\left[f^1dx^2+x^2d\mathrm{\Omega }_{Dp3}^2\right],`$ (10)
$`e^{2\varphi }`$ $`=`$ $`H_p^{\frac{(D3)a_p}{\mathrm{\Delta }_p}},A_{tx_1\mathrm{}x_p}={\displaystyle \frac{2}{\sqrt{\mathrm{\Delta }_p}}}{\displaystyle \frac{\mu \mathrm{cosh}\delta _p\mathrm{sinh}\delta _p}{x^{Dp4}}}H_p^1,`$ (11)
where
$`H_p`$ $`=`$ $`1+{\displaystyle \frac{\mu \mathrm{sinh}^2\delta _p}{x^{Dp4}}},f=1{\displaystyle \frac{\mu }{x^{Dp4}}},`$ (12)
$`\mathrm{\Delta }_p`$ $`=`$ $`{\displaystyle \frac{(D3)a_p^2}{2}}+{\displaystyle \frac{2(p+1)(Dp4)}{D3}}.`$ (13)
We regard this solution as being obtained by compactifying a $`D`$-dimensional $`(p+1)`$-brane with a dilaton coupling parameter $`a_p`$ (i.e., the solution (5) with $`H=1`$) along a longitudinal direction (which is the $`y`$-direction in the notation of Eq. (5)) on $`S^1`$. Since the parameter $`\mathrm{\Delta }_p`$ is invariant under reductions or oxidations which do not involve field truncation, one can see that $`a_p`$ is related to $`a_{p+1}`$ of the solution in one higher dimension, namely that of the $`(p+1)`$-brane in Eq. (5), as
$$(D3)a_p^2=(D2)a_{p+1}^2+4\frac{(Dp4)^2}{(D2)(D3)}.$$
(14)
### 3.1 Probing with a $`(p+1)`$-brane
The worldvolume action for a dilatonic $`(p+1)`$-brane with the following bulk action:
$$S_E=\frac{1}{2\kappa _D^2}d^Dx\sqrt{g}\left[\frac{4}{D2}(\varphi )^2\frac{1}{2(p+3)!}e^{2a_{p+1}\varphi }F_{p+3}^2\right]$$
(15)
has the following form:
$`S_\sigma `$ $`=`$ $`T_{p+1}{\displaystyle }d^{p+2}\xi [e^{a_{p+1}\varphi }\sqrt{\mathrm{det}_aX^\mu _bX^\nu g_{\mu \nu }}`$ (17)
$`+{\displaystyle \frac{\sqrt{\mathrm{\Delta }_{p+1}}}{2}}{\displaystyle \frac{1}{(p+2)!}}ϵ^{a_1\mathrm{}a_{p+2}}_{a_1}X^{\mu _1}\mathrm{}_{a_{p+2}}X^{\mu _{p+2}}A_{\mu _1\mathrm{}\mu _{p+2}}],`$
where $`T_{p+1}`$ is the tension of the probe $`(p+1)`$-brane and the target space fields $`g_{\mu \nu }`$, $`\varphi `$ and $`A_{\mu _1\mathrm{}\mu _{p+2}}`$ are the background fields (produced by the source brane) in which the probe $`(p+1)`$-brane with the target space coordinates $`X^\mu (\xi ^a)`$ ($`\mu =0,1,\mathrm{},D1`$) and the worldvolume coordinates $`\xi ^a`$ ($`a=0,1,\mathrm{},p+1`$) moves.
In the static gauge, in which $`X^a=\xi ^a`$, the pull-back fields for the probe $`(p+1)`$-brane, oriented in the same way as the source $`(p+1)`$-brane, take the following forms:
$`\widehat{G}_{ab}`$ $``$ $`g_{\mu \nu }_aX^\mu _bX^\nu =g_{ab}+g_{ij}_aX^i_bX^j,`$ (18)
$`\widehat{A}_{a_1\mathrm{}a_{p+2}}`$ $``$ $`A_{\mu _1\mathrm{}\mu _{p+2}}_{a_1}X^{\mu _1}\mathrm{}_{a_{p+2}}X^{\mu _{p+2}}=A_{a_1\mathrm{}a_{p+2}},`$ (19)
where the indices $`i,j=1,\mathrm{},Dp2`$ label the transverse space of the probe $`(p+1)`$-brane, i.e., $`(X^i)=(x_1,\mathrm{},x_{Dp2})`$ in the notation of Eq. (5). So, the worldvolume action (17) takes the following form:
$$S_\sigma =T_{p+1}d^{p+2}\xi \left[e^{a_{p+1}\varphi }\sqrt{\mathrm{det}\left(g_{ab}+g_{ij}_aX^i_bX^j\right)}+\frac{\sqrt{\mathrm{\Delta }_{p+1}}}{2}A_{01\mathrm{}p+1}\right].$$
(20)
From now on, we assume that the target space transverse coordinates $`X^i`$ for the probe $`(p+1)`$-brane depend on the time coordinate $`\tau =\xi ^0`$ only, i.e., $`X^i=X^i(\tau )`$.
By substituting the explicit expressions (5) for the source background fields of the brane-world $`(p+1)`$-brane into the general expression (20) for the probe $`(p+1)`$-brane action, one obtains the following:
$`S_\sigma `$ $`=`$ $`T_{p+1}{\displaystyle }d^{p+2}\xi H^{\frac{4}{\mathrm{\Delta }}}[H_{p+1}^1f^{\frac{1}{2}}\sqrt{1H_{p+1}^{\frac{4}{\mathrm{\Delta }_{p+1}}}\left\{f^2\left({\displaystyle \frac{dx}{d\tau }}\right)^2+x^2f^1\mu _m^2\left({\displaystyle \frac{d\varphi _m}{d\tau }}\right)^2\right\}}`$ (22)
$`+{\displaystyle \frac{\mu \mathrm{cosh}\delta _{p+1}\mathrm{sinh}\delta _{p+1}}{x^{Dp4}}}H_{p+1}^1].`$
On the other hand, the probe $`p`$-brane action in the source background (11) of the $`(D1)`$-dimensional $`p`$-brane is
$`S_\sigma `$ $`=`$ $`T_p{\displaystyle }d^{p+1}\xi [H_p^1f^{\frac{1}{2}}\sqrt{1H_p^{\frac{4}{\mathrm{\Delta }_p}}\left\{f^2\left({\displaystyle \frac{dx}{d\tau }}\right)^2+x^2f^1\mu _m^2\left({\displaystyle \frac{d\varphi _m}{d\tau }}\right)^2\right\}}`$ (24)
$`+{\displaystyle \frac{\mu \mathrm{cosh}\delta _p\mathrm{sinh}\delta _p}{x^{Dp4}}}H_p^1].`$
Here, the angular coordinates $`0\varphi _m<2\pi `$ ($`m=1,\mathrm{},[(Dp2)/2]`$) are associated with $`[(Dp2)/2]`$ rotation planes in the transverse space of the branes (with the coordinates $`𝐱`$) and the index $`m`$ is summed over $`m=1,\mathrm{},[(Dp2)/2]`$. The remaining angular coordinates, which determine the direction cosines $`\mu _m`$, are constant due to the conservation of the direction of the angular momentum.
We see that the probe actions (22) and (24) have the same form except that the probe $`(p+1)`$-brane action (22) has an additional overall factor $`H^{4/\mathrm{\Delta }}`$. So, the dynamics of the probe $`(p+1)`$-brane in the background of the brane-world $`(p+1)`$-brane is identical to that of the probe $`p`$-brane in the background of the $`p`$-brane in one lower dimensions. (Note, $`\mathrm{\Delta }_p=\mathrm{\Delta }_{p+1}`$, provided the $`(D1)`$-dimensional $`p`$-brane is obtained by compactifying the $`D`$-dimensional $`(p+1)`$-brane on $`S^1`$ without field truncation.) The effect of the overall factor $`H^{4/\mathrm{\Delta }}`$ in the former case is to effectively increase \[decrease\] the tension of the probe $`(p+1)`$-brane when $`\mathrm{\Delta }>0`$ \[$`\mathrm{\Delta }<0`$\], namely $`T_{p+1}^{\mathrm{eff}}=T_{p+1}H^{4/\mathrm{\Delta }}`$.
Since the probe actions for the two cases have the same form, one also expects that the first law of black brane thermodynamics of the $`p`$-brane in one lower dimensions can be extracted from that of the $`(p+1)`$-brane in the domain wall, and vice versa. The first law of black brane thermodynamics of the latter brane with the solution given by Eq. (5) is
$$\delta M_{p+1}=T_H^{p+1}\delta S_{p+1}+\mathrm{\Phi }_{p+1}\delta Q_{p+1},$$
(25)
where $`M_{p+1}`$ and $`S_{p+1}`$ are the ADM mass and the entropy of the source $`(p+1)`$-brane per unit $`(p+1)`$-brane worldvolume, $`Q_{p+1}`$ is the source $`(p+1)`$-brane charge normalized to take integer values (i.e., the number of elementary $`(p+1)`$-branes with unit charge), and the Hawking temperature $`T_H^{p+1}`$ and the chemical potential $`\mathrm{\Phi }_{p+1}`$ of the source $`(p+1)`$-brane are given by
$$T_H^{p+1}=\frac{Dp4}{4\pi \mu ^{\frac{1}{Dp4}}\mathrm{cosh}^{\frac{4}{\mathrm{\Delta }_{p+1}}}\delta _{p+1}},\mathrm{\Phi }_{p+1}=\frac{2}{\sqrt{\mathrm{\Delta }_{p+1}}}T_{p+1}H^{\frac{4}{\mathrm{\Delta }}}\mathrm{tanh}\delta _{p+1}.$$
(26)
In the case of the source $`p`$-brane in $`(D1)`$-dimensions with the solution given by Eq. (11), the temperature $`T_H^p`$ and the chemical potential $`\mathrm{\Phi }_p`$ are respectively related to those of the $`(p+1)`$-brane as $`T_H^p=T_H^{p+1}`$ and $`\mathrm{\Phi }_p/T_p=\mathrm{\Phi }_{p+1}/(T_{p+1}H^{4/\mathrm{\Delta }})`$, if we let $`\delta _p=\delta _{p+1}`$. Note, $`\mathrm{\Delta }_{p+1}`$ in Eq. (5) and $`\mathrm{\Delta }_p`$ in Eq. (11) are the same, if the $`p`$-brane is obtained from the $`(p+1)`$-brane through the dimensional reduction without field truncation. One can think of the changes $`\delta M_{p+1}`$, $`\delta S_{p+1}`$ and $`\delta Q_{p+1}`$ as being due to an addition of the probe $`(p+1)`$-brane to the source $`(p+1)`$-brane . Namely, we bring the probe $`(p+1)`$-brane with a unit brane charge from spatial infinity ($`x=\mathrm{}`$) to the source brane horizon ($`x=x_H=\mu ^{1/(Dp4)}`$). Then, one can interpret $`T_H^{p+1}\delta S_{p+1}`$ as the heat released by the probe $`(p+1)`$-brane while it falls inside the source $`(p+1)`$-brane, which is just the difference in static potential energy $`V_{p+1}(x)`$ of the probe $`(p+1)`$-brane, i.e., $`T_H^{p+1}\delta S_{p+1}=V(\mathrm{})V(x_H)`$ . From the above probe actions (22) and (24), one obtains the following static potentials on the probe branes:
$$V_{p+1}=T_{p+1}H^{\frac{4}{\mathrm{\Delta }}}\left(f^{\frac{1}{2}}+\frac{\mu \mathrm{cosh}\delta _{p+1}\mathrm{sinh}\delta _{p+1}}{x^{Dp4}}\right)H_{p+1}^1,$$
(27)
for the probe $`(p+1)`$-brane, and
$$V_p=T_p\left(f^{\frac{1}{2}}+\frac{\mu \mathrm{cosh}\delta _p\mathrm{sinh}\delta _p}{x^{Dp4}}\right)H_p^1,$$
(28)
for the probe $`p`$-brane. As expected, in the extreme limit ($`\mu 0`$ with $`\mu e^{2\delta }`$ finite constant), the potentials are constant in accordance with the no-force condition for extreme branes. We see that the two static potentials are related as $`V_{p+1}/(T_{p+1}H^{4/\mathrm{\Delta }})=V_p/T_p`$ and therefore $`T_H^{p+1}\delta S_{p+1}/(T_{p+1}H^{4/\mathrm{\Delta }})=T_H^p\delta S_p/T_p`$, if $`\delta _p=\delta _{p+1}`$. The probe actions (22) and (24) with $`\delta _p=\delta _{p+1}`$ imply that the mass density changes are related as $`\delta M_{p+1}/(T_{p+1}H^{4/\mathrm{\Delta }})=\delta M_p/T_p`$. Since the probe branes have unit charges, $`\delta Q_p=1=\delta Q_{p+1}`$. Gathering all the above, one can bring the first law of the black $`(p+1)`$-brane thermodynamics (25) to the following first law of thermodynamics of the black $`p`$-brane in $`D1`$ dimensions:
$$\delta M_p=T_H^p\delta S_p+\mathrm{\Phi }_p\delta Q_p.$$
(29)
### 3.2 Probing with a test particle
In analyzing the dynamics of a test particle in a curved spacetime background, it is convenient to utilize the symmetry of the spacetime. The Killing vectors of the spacetime metrics of both of the solutions (5) and (11) are $`/t`$, $`/w_i`$ and $`/\varphi _m`$. Contracting these Killing vectors with the velocity $`U^\mu =dx^\mu /d\lambda `$ of the test particle along the geodesic path $`x^\mu (\lambda )`$ parametrized by an affine parameter $`\lambda `$, one obtains the following constants of motion for the test particle:
$`E`$ $`=`$ $`g_{\mu \nu }\left({\displaystyle \frac{}{t}}\right)^\mu U^\nu =g_{tt}{\displaystyle \frac{dt}{d\lambda }},`$ (30)
$`p^i`$ $`=`$ $`g_{\mu \nu }\left({\displaystyle \frac{}{w_i}}\right)^\mu U^\nu =g_{ii}{\displaystyle \frac{dw_i}{d\lambda }},`$ (31)
$`J^m`$ $`=`$ $`g_{\mu \nu }\left({\displaystyle \frac{}{\varphi _m}}\right)^\mu U^\nu =g_{\varphi _m\varphi _m}{\displaystyle \frac{d\varphi _m}{d\lambda }}.`$ (32)
In addition, there is another constant of motion associated with metric compatibility along the geodesic path:
$$ϵ=g_{\mu \nu }\frac{dx^\mu }{d\lambda }\frac{dx^\nu }{d\lambda },$$
(33)
where $`ϵ=1,0`$ respectively for a massive particle (i.e., a timelike geodesic) and a massless particle (i.e., a null geodesic).
For the simplicity of the calculation, we shift the transverse coordinate $`y`$ of the domain wall so that the harmonic function for the domain wall takes the form $`H=Qy`$, where we restrict to the region $`y0`$. Then, we apply the following change of coordinate:
$$y=\left(\frac{\mathrm{\Delta }+2}{\mathrm{\Delta }}Q^{\frac{2}{\mathrm{\Delta }}}z\right)^{\frac{\mathrm{\Delta }}{\mathrm{\Delta }+2}}$$
(34)
to bring the domain wall metric to the conformally flat form. In this new coordinate, the metric in Eq. (5) takes the following form:
$`ds^2`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{\Delta }+2}{\mathrm{\Delta }}}Qz\right)^{\frac{4}{(D2)(\mathrm{\Delta }+2)}}[H_{p+1}^{\frac{4(Dp4)}{(D2)\mathrm{\Delta }_{p+1}}}(fdt^2+dw_1^2+\mathrm{}+dw_p^2)`$ (36)
$`+H_{p+1}^{\frac{4(p+2)}{(D2)\mathrm{\Delta }_{p+1}}}(f^1dx^2+x^2d\mathrm{\Omega }_{Dp3}^2)+H_{p+1}^{\frac{4(Dp4)}{(D2)\mathrm{\Delta }_{p+1}}}dz^2].`$
For the test particle moving along the $`z`$-direction, i.e., only the $`z`$-component of $`U^\mu `$ is non-zero, the geodesic motion is described by the following equation, derived from the geodesic equation and Eq. (33):
$$\frac{d}{d\lambda }\left[z^{\frac{4}{(D2)(\mathrm{\Delta }+2)}}\frac{dz}{d\lambda }\right]=\frac{2ϵ}{(D2)(\mathrm{\Delta }+2)}\left(\frac{\mathrm{\Delta }+2}{\mathrm{\Delta }}Q\right)^{\frac{4}{(D2)(\mathrm{\Delta }+2)}}H_{p+1}^{\frac{4(Dp4)}{(D2)\mathrm{\Delta }_{p+1}}}\frac{1}{z}.$$
(37)
The geodesic path $`z(\lambda )`$ for a massless test particle (i.e., the $`ϵ=0`$ case) is $`z=\mathrm{constant}`$ or
$$z=z_0\lambda ^{\frac{(D2)(\mathrm{\Delta }+2)}{(D2)(\mathrm{\Delta }+2)+4}},$$
(38)
where $`z_0`$ is an arbitrary constant. The general explicit expression for the timelike geodesic path (i.e., the $`ϵ=1`$ case) is hard to obtain. So, in the following we shall consider only the case of the null geodesic motion along the $`z`$-direction. The null geodesic path $`z=\mathrm{constant}`$ simply corresponds to the motion constrained along the longitudinal directions of the domain wall.
Making use of the constants of the motion in Eq. (32), one can put Eq. (33) into the following form:
$$\left(\frac{dx}{d\lambda }\right)^2+\frac{g_{zz}}{g_{xx}}\left(\frac{dz}{d\lambda }\right)^2+\frac{J_m^2}{g_{xx}g_{\varphi _m\varphi _m}}+\frac{E^2}{g_{xx}g_{tt}}+\frac{ϵ}{g_{xx}}=0,$$
(39)
where the index $`m`$ is summed over $`m=1,\mathrm{},[(Dp2)/2]`$. The angular coordinates of the transverse space of the branes (with the coordinates $`𝐱`$) in the spherical coordinates, except for the ones $`\varphi _m`$ associated with the angular momenta $`J_m`$, are constant due to the conservation of the direction of the angular momentum. And we are considering the motion of the test particle with the longitudinal coordinates $`w_i`$ of the branes constant, which is possible due to the conservation of the linear momenta $`p^i`$ along those directions. By plugging the explicit expression (36) for the background metric of the source brane into the general expression (39), we obtain
$`\left({\displaystyle \frac{dx}{d\lambda }}\right)^2`$ $`+`$ $`𝒞^2\left[\left\{𝒞^2H_{p+1}^{\frac{4}{\mathrm{\Delta }_{p+1}}}\left({\displaystyle \frac{dz}{d\lambda }}\right)^2+{\displaystyle \frac{𝒥^2}{H_{p+1}^{\frac{8(p+2)}{(D2)\mathrm{\Delta }_{p+1}}}x^2}}\right\}fE^2H_{p+1}^{\frac{4(D2p6)}{(D2)\mathrm{\Delta }_{p+1}}}\right]`$ (40)
$`+`$ $`ϵ_{p+1}𝒞^1H_{p+1}^{\frac{4(p+2)}{(D2)\mathrm{\Delta }_{p+1}}}f=0,`$ (41)
where $`𝒞=\left(\frac{\mathrm{\Delta }+2}{\mathrm{\Delta }}Qz\right)^{\frac{4}{(D2)(\mathrm{\Delta }+2)}}`$ is the conformal factor in the metric (36), $`z=z(\lambda )`$ for the null geodesic motion is constant or is given by Eq. (38), and $`ϵ_{p+1}=0,1`$ respectively for the massless and the massive test particle. On the other hand, for the test particle in the source background (11) of the $`(D1)`$-dimensional $`p`$-brane, the geodesic motion is described by
$$\left(\frac{dx}{d\lambda }\right)^2+\left[ϵ_pH_p^{\frac{4(p+1)}{(D3)\mathrm{\Delta }_p}}+\frac{𝒥^2}{H_p^{\frac{8(p+1)}{(D3)\mathrm{\Delta }_p}}x^2}\right]fE^2H_p^{\frac{4(D2p5)}{(D3)\mathrm{\Delta }_p}}=0,$$
(42)
where $`ϵ_p=1,0`$ respectively for the timelike and the null geodesic. Here, $`𝒥`$ in the above is defined in terms of the conserved angular momenta $`J^m`$ of the test particle as
$$𝒥^2\underset{m=1}{\overset{[\frac{Dp2}{2}]}{}}\frac{(J^m)^2}{\mu _m^2},$$
(43)
where the direction cosines $`\mu _m`$ specifying the direction of $`x`$ are constant due to conservation of the direction of angular momentum (therefore $`𝒥`$ is also constant).
First, when the $`(p+1)`$-brane is uncharged (i.e., $`H_{p+1}=1`$), one can bring the equation (41) for the null geodesic motion ($`ϵ_{p+1}=0`$) to the form of the equation (42) for the time-like geodesic motion ($`ϵ_p=1`$) in the $`(D1)`$-dimensional uncharged $`p`$-brane background. To see this, we consider the following equation obtained from Eq. (41) by setting $`ϵ_{p+1}=0`$ and $`H_{p+1}=1`$, and substituting (38):
$`\left({\displaystyle \frac{dx}{d\lambda }}\right)^2+\left({\displaystyle \frac{\mathrm{\Delta }+2}{\mathrm{\Delta }}}Qz_0\right)^{\frac{8}{(D2)(\mathrm{\Delta }+2)}}\lambda ^{\frac{8}{(D2)(\mathrm{\Delta }+2)+4}}[\{\left({\displaystyle \frac{\mathrm{\Delta }+2}{\mathrm{\Delta }}}Qz_0\right)^{\frac{8}{(D2)(\mathrm{\Delta }+2)}}`$ (44)
$`\times \left({\displaystyle \frac{(D2)(\mathrm{\Delta }+2)z_0}{(D2)(\mathrm{\Delta }+2)+4}}\right)^2+{\displaystyle \frac{𝒥^2}{x^2}}\}(1{\displaystyle \frac{\mu }{x^{Dp4}}})E^2]=0.`$ (45)
By redefining the radial coordinate, constants of the motion and parameters in the following way:
$`\stackrel{~}{x}`$ $``$ $`𝒜x,\stackrel{~}{E}𝒜E,\stackrel{~}{J}𝒜^2J,\stackrel{~}{\mu }𝒜^{Dp4}\mu ,`$ (46)
$`\stackrel{~}{\lambda }`$ $``$ $`{\displaystyle \frac{(D2)(\mathrm{\Delta }+2)+4}{(D2)(\mathrm{\Delta }+2)}}\left({\displaystyle \frac{\mathrm{\Delta }+2}{\mathrm{\Delta }}}Qz_0\right)^{\frac{2}{(D2)(\mathrm{\Delta }+2)}}\lambda ^{\frac{(D2)(\mathrm{\Delta }+2)}{(D2)(\mathrm{\Delta }+2)+4}},`$ (47)
where
$$𝒜\left(\frac{(D2)(\mathrm{\Delta }+2)z_0}{(D2)(\mathrm{\Delta }+2)+4}\right)^1\left(\frac{\mathrm{\Delta }+2}{\mathrm{\Delta }}Qz_0\right)^{\frac{4}{(D2)(\mathrm{\Delta }+2)}},$$
(48)
one can bring Eq. (45) into the following form:
$$\left(\frac{d\stackrel{~}{x}}{d\stackrel{~}{\lambda }}\right)^2+\left(1+\frac{\stackrel{~}{𝒥}^2}{\stackrel{~}{x}^2}\right)\left(1\frac{\stackrel{~}{\mu }}{\stackrel{~}{x}^{Dp4}}\right)=\stackrel{~}{E}^2.$$
(49)
This reproduces the equation for the timelike geodesic motion in the background of uncharged $`(D1)`$-dimensional $`p`$-brane, i.e., Eq. (42) with $`ϵ_p=1`$ and $`H_p=1`$. This result generalizes the result of Ref. to the case of an uncharged black brane in a dilatonic domain wall in arbitrary spacetime dimensions.
Next, we consider the charged branes. The equation for the null geodesic motion (with nontrivial lightlike motion along the $`z`$-direction) in the background of the brane-world charged $`(p+1)`$-brane, i.e., Eq. (41) with $`ϵ_{p+1}=0`$ and Eq. (38) substituted, reduces to the following form after the quantities are redefined as in Eq. (47):
$$\left(\frac{d\stackrel{~}{x}}{d\stackrel{~}{\lambda }}\right)^2+\left[\stackrel{~}{H}_{p+1}^{\frac{4}{\mathrm{\Delta }_{p+1}}}\frac{\stackrel{~}{𝒥}^2}{\stackrel{~}{H}_{p+1}^{\frac{8(p+2)}{(D2)\mathrm{\Delta }_{p+1}}}\stackrel{~}{x}^2}\right]\stackrel{~}{f}\stackrel{~}{E}^2\stackrel{~}{H}_{p+1}^{\frac{4(D2p6)}{(D2)\mathrm{\Delta }_{p+1}}}=0,$$
(50)
where
$$\stackrel{~}{H}_{p+1}=1+\frac{\mu \mathrm{sinh}^2\delta _{p+1}}{(𝒜\stackrel{~}{x})^{Dp4}},\stackrel{~}{f}=1\frac{\stackrel{~}{\mu }}{\stackrel{~}{x}^{Dp4}}.$$
(51)
This equation is different from the equation for the timelike geodesic in the background of the $`(D1)`$-dimensional $`p`$-brane, i.e., Eq. (42) with $`ϵ_p=1`$, since the powers of the harmonic functions are different in the two equations. This difference might be attributed to the fact that when one compactifies the Einstein-frame metric for the $`D`$-dimensional $`(p+1)`$-brane (i.e., Eq. (36) without the $`z`$-dependent conformal factor) along one of its longitudinal directions (i.e., the $`z`$-direction) by using the KK metric Ansatz without the Weyl-scaling factor in the $`(D1)`$-dimensional part of the metric (i.e., $`g_{\mu \nu }=\mathrm{diag}(\overline{g}_{\overline{\mu }\overline{\nu }},\overline{\phi })`$ with $`\mu ,\nu =0,1,\mathrm{},D1`$ and $`\overline{\mu },\overline{\nu }=0,1,\mathrm{},D2`$), one gets non-Einstein-frame metric for the $`(D1)`$-dimensional $`p`$-brane. However, as can be seen from Eq. (39), even in such non-Einstein-frame spacetime in $`D1`$ dimensions, the equation for the timelike geodesic will look different because of the dependence of the $`(z,z)`$-component of the $`D`$-dimensional metric on $`x`$. (Cf. The second term on the LHS of Eq. (39) for the null geodesic motion in $`D`$ dimensions is identified with the last term on the LHS of Eq. (39) for the timelike geodesic motion in $`D1`$ dimensions.) We just write down the equation for the timelike geodesic motion in such $`(D1)`$-dimensional background for comparison:
$$\left(\frac{dx}{d\lambda }\right)^2+\left[H_{p+1}^{\frac{4(p+2)}{(D2)\mathrm{\Delta }_{p+1}}}+\frac{𝒥^2}{H_{p+1}^{\frac{8(p+2)}{(D2)\mathrm{\Delta }_{p+1}}}x^2}\right]fE^2H_{p+1}^{\frac{4(D2p6)}{(D2)\mathrm{\Delta }_{p+1}}}=0.$$
(52)
As mentioned, the first terms in the square brackets of Eqs. (50) and (52) are different. So, only when the probe motion along the $`z`$-direction is trivial, i.e., $`z=\mathrm{constant}`$, the null geodesic motion in the background of the $`D`$-dimensional $`(p+1)`$-brane (described by Eq. (41) with $`ϵ_{p+1}=0`$ and $`z=\mathrm{constant}`$) reproduces the null geodesic motion in the background of the $`p`$-brane in one lower dimensions (described by Eq. (42) with $`ϵ_p=0`$). In the case of the uncharged branes, the above problems did not arise because the $`(z,z)`$-component of the $`D`$-dimensional metric is independent of $`x`$ and the dimensional reduction of the Einstein-frame metric for the $`D`$-dimensional uncharged $`(p+1)`$-brane along a longitudinal direction (without the Weyl-scaling term in the metric) leads to the Einstein-frame metric for the $`(D1)`$-dimensional uncharged $`p`$-brane. It is interesting to note that for the $`p=D4`$ case the equation (50) for the null geodesic motion (with nontrivial motion along the $`z`$-direction) in the background of the brane-world $`(p+1)`$-brane reproduces the equation for the timelike geodesic motion in the background of the $`p`$-brane in one lower dimensions, i.e., Eq. (42) with $`ϵ_p=1`$ and Eq. (52). In such case, the transverse (to the domain wall) component of the metric (36), i.e., the $`(z,z)`$-component, is independent of the longitudinal coordinates of the domain wall. This is in accordance with our speculation on the source of disparity in the probe particle dynamics in the backgrounds of charged branes.
In fact, one of the assumptions of the RS model is that the $`D`$-dimensional conformally flat metric (of the domain wall) should have the perturbation around the flat metric along the longitudinal directions of the domain wall, only. Indeed, as mentioned in the above, one can see from Eq. (39) that the extra space component of the metric, i.e., $`g_{zz}`$, should be independent of the longitudinal coordinates of the domain wall, in order for the null geodesic motion (with non-trivial lightlike motion along the $`z`$-direction) in $`D`$ dimensions to reproduce the timelike geodesic motion in the background of the source $`p`$-brane in $`D1`$ dimensions. So, for the $`p=D4`$ case, even if the branes are charged, the null geodesic motion in the source background of the brane-world $`(p+1)`$-brane reproduces the timelike geodesic motion in the source background of the $`p`$-brane in one lower dimensions, because the extra space component $`g_{zz}`$ of the brane-world $`(p+1)`$-brane metric is independent of the longitudinal coordinates of the domain wall. Also, recently, it is shown that the metric perturbation in the direction transverse to the domain wall is not localized on the brane. On the other hand, in general, the extra space component $`g_{zz}`$ of the spacetime metrics for charged brane solutions in the domain wall depends on the longitudinal coordinates $`𝐱`$ of the domain wall. So, it seems inevitable that all the charged branes in brane worlds do not reproduce physics in one lower dimensions. This might be the indication of either the need to modify physics in one lower dimensions (which seems unlikely) or the limitation of the current RS model that needs to be modified so that it can accommodate, for example, charged black holes or branes that will reproduce lower-dimensional physics. (Also, the non-dilatonic domain wall of the RS model in Refs. does not admit even charged black string solutions, which are supposed to be identified as charged black holes in four dimensions, because of the constraint on the dilaton coupling parameters.)
## 4 Conclusion
We studied the non-extreme dilatonic $`(p+1)`$-brane in the bulk of extreme dilatonic domain wall, where one of the longitudinal directions of the brane is along the transverse direction of the wall. Such $`(p+1)`$-brane is expected to be the domain wall bulk counterpart to the ordinary $`p`$-brane observed on the hypersurface of the domain wall. We studied the probe dynamics on such background for the purpose of seeing whether the effective compactification of such $`(p+1)`$-brane through the RS type domain wall leads to the physics of $`p`$-brane in one lower dimensions obtained through the ordinary KK compactification or not. In this paper, we found partial agreement of the probe dynamics in the two backgrounds. Namely, in the case of the probe $`(p+1)`$-brane in the background of the source $`(p+1)`$-brane in the bulk of the domain wall, its dynamics reproduces the dynamics of the probe $`p`$-brane in the background of the source $`p`$-brane originated from the source $`(p+1)`$-brane (without the domain wall) through the ordinary KK compactification. However, in the case of the test particle moving in the same source backgrounds, we found agreement only for the case when the source branes are uncharged. We have attributed this difference to the fact that the metric for the charged brane solutions does not satisfy the RS gauge condition. Namely, RS showed that the gravity in one lower dimensions is reproduced only for the case where the domain wall metric has perturbations along the longitudinal directions of the wall, only, whereas the transverse (to the wall) component of the metric for charged branes in the bulk of domain walls in general has non-trivial dependence on the domain wall worldvolume coordinates. Perhaps, in some cases RS type models might give rise to the effective theory in one lower dimensions different from the one that would have been obtained through the ordinary KK compactification. On the other hand, the RS gauge for the domain wall metric perturbation used in Ref. is applicable only for the case when there are no additional fields in the action. Namely, when additional fields are added to the action, one might have to modify the RS gauge condition of Ref. , possibly, to include the transverse perturbation as well. Furthermore, the trick used in studying the geodesic motion of the test particle with non-trivial motion along the extra space, which was devised in Ref. and also applied in this paper, rather seems not to be rigorous. Definite answer to this point cannot be given until one would solve the full coupled nonlinear geodesic equations $`\text{}^\mu +\mathrm{\Gamma }_{\nu \rho }^\mu \dot{x}^\nu \dot{x}^\rho =0`$. We defer the above unanswered questions to our future work. |
warning/0001/astro-ph0001210.html | ar5iv | text | # The stellar populations of early-type galaxies in the Fornax cluster
## 1 INTRODUCTION
Great efforts have been made in the last few years to develop evolutionary stellar population synthesis models \[Bruzual & Charlot 1993, Worthey 1994, Weiss, Peletier & Matteucci 1995, Vazdekis et al. 1996, Kodama & Arimoto 1997\] in order to analyze the integrated light of galaxies and derive estimates of their mean ages and metal abundances. One of the main obstacles in the interpretation has been the age/metallicity degeneracy in old stellar populations. As pointed out by Worthey (1994) the integrated spectral energy distribution (SED) of an old ($`>2`$ Gyrs) stellar population looks almost identical when the age is doubled and total metallicity reduced by a factor of three at the same time. Therefore two galaxies with almost identical broad-band colours can have significantly different ages and metallicities. In the optical wavelength range, only a few narrow band absorption line-strength indices and the 4000 Å break (see also Gorgas et al. 1999) have so far been identified which can break this degeneracy.
One of the most successful and widely used methods for measuring the strength of age/metallicity discriminating absorption features is the Lick/IDS system \[Burstein et al. 1984, Worthey et al. 1994, Trager et al. 1998\] which has been used by many authors \[González 1993, Davies, Sadler & Peletier 1993, Fisher, Franx & Illingworth 1995, Fisher, Franx & Illingworth 1996, Ziegler & Bender 1997, Longhetti et al. 1998, Mehlert 1998, Jørgensen 1999\]. In contrast with high resolution index systems \[Rose 1994, Jones & Worthey 1995\], which promise a better separation of age and metalllitiy, the Lick/IDS system allows the investigation of dynamically hot galaxies that have intrinsically broad absorption lines. By plotting an age sensitive index and a metallicity sensitive index against each other, one can (partially) break the age/metallicity degeneracy and estimate, with the help of model predictions, the luminosity weighted age and metallicity of an integrated stellar population (see Figure 4). Most recently, Jørgensen (1999) used this methodology to investigate the stellar populations of a large sample of early-type galaxies in the Coma cluster. She concluded that there are real variations in both the ages and the abundances while an anti-correlation between the mean ages and the mean abundances makes it possible to maintain a low scatter in scaling relations such as Mg–$`\sigma _0`$. Colless et al. (1999) present similar conclusions from the analysis of a combination of the Mg–$`\sigma _0`$ relation and the Fundamental Plane in a large sample of cluster early-type galaxies.
The spread in the ages for early-type galaxies and the anti-correlation of age and metallicity found by the previous authors supports the hierarchical picture for the construction of galaxies in which galaxies form via several mergers involving star-formation \[Baugh et al. 1996, Kauffmann 1996\]. However, the results are inconsistent with the conventional view that all luminous elliptical galaxies are old and coeval. In this picture the global spectrophotometric relations observed for ellipticals, for example the colour-magnitude relation \[Visvanathan & Sandage 1977, Bower, Lucey & Ellis 1992, Terlevich 1998\] are explained by the steady increase in the abundance of heavy elements with increasing galaxy mass. This increase arises naturally in galactic wind models such as that of Arimoto & Yoshii (1987) and Kodama & Arimoto (1997).
Although with line-strength indices we can (partially) break the age/metallicity degeneracy this is by no means the last obstacle to overcome on our way to fully understand the stellar populations of early-type galaxies and the cause of scaling ralations. Since the late 70’s evidence has been accumulating that abundance ratios in galaxies are often non-solar. In particular the Magnesium to Iron ratio seems to be larger in luminous early-type galaxies when compared to solar neighbourhood stars \[O’Connell 1976, Peletier 1989, Worthey, Faber & González 1992, Davies, Sadler & Peletier 1993, Henry & Worthey 1999, Jørgensen 1999\]. However, with only a very few exceptions (e.g., Weiss, Peletier & Matteucci 1995), non-solar abundance ratios have not yet been incorporated in the model predictions. Among other issues this seems to be the most important single problem which prevents us from deriving accurate absolute age and metallicity estimates from integrated light spectroscopy \[Worthey 1998\]. Nevertheless with the current models and high S/N data we are able to study relative trends in ages and abundances as well as start to investigate the effects of non-solar abundance ratios for individual elements \[Worthey 1998, Peletier 2000\].
In this paper, high S/N nuclear spectra of a complete sample of early-type galaxies in the Fornax cluster brighter than $`M_B=17`$ are analyzed in the Lick/IDS system. The early results of this study have already been presented in a letter to this journal \[Kuntschner & Davies 1998\]. This paper is organized as follows. Section 2 describes the sample selection and basic data reduction. The calibration of the line-strength indices to the Lick/IDS system is presented in Section 3 and the Appendix. Section 4 presents a consistency test for the model predictions and our measured line-strength indices. In Section 5 several index combinations are compared to model predictions. In particular the effects of non-solar abundance ratios, composite stellar populations and age/metallicity estimates of the integrated light are discussed. In Section 6 observed index–$`\sigma _0`$ relations are presented. Relations between derived parameters such as age, metallicity and \[Mg/Fe\] with the central velocity dispersion are investigated in Section 7. We then discuss the implications of our results in Section 8 and present the conclusions in Section 9. The fully corrected Lick/IDS indices of our sample are tabulated in the Appendix.
## 2 THE OBSERVATIONS AND DATA REDUCTION
### 2.1 The sample
Our sample of 22 early-type galaxies has been selected from the catalogue of Fornax galaxies (Ferguson 1989, hereafter F89), in order to obtain a complete sample down to $`B_T=14.2`$ or $`M_B=17.`$<sup>1</sup><sup>1</sup>1Adopting a distance modulus of $`mM=31.2`$; based on I-band surface brightness fluctuations \[Jensen, Tonry & Luppino 1998\]. This corresponds to H$`{}_{0}{}^{}80`$ km s<sup>-1</sup> Mpc<sup>-1</sup> for a flat Universe. We have adopted the morphological classifications given by F89 and checked them with images we obtained on the Siding Spring 40$`\mathrm{}`$ telescope. From these we noted a central dust lane in ESO359-G02 and a central disc in ESO358-G59 which led us to classify them as lenticular galaxies. NGC 1428 was not observed because of a bright star close to its centre. We also added the elliptical galaxy IC2006 to our sample, as it was not classified by F89. The bona-fide elliptical NGC3379 was observed as a calibration galaxy. The observations were carried out with the AAT (3.9m) on the nights of 1996 December 6-8 using the RGO spectrograph. The characteristics of the detector and the instrument set-up is given in Table 1.
Typically, exposure times were between 300 and 1800 sec per galaxy (see Table 2 for a detailed listing). For most of the observations the slit was centred on the nucleus at $`\mathrm{P}A=90^{}`$. The seeing was generally better than one arcsec. Additionally we observed 15 different standard stars (mainly K-giants) during twilight to act as templates for velocity dispersion measurements as well as to calibrate our line-strength indices to the Lick/IDS system \[Worthey et al. 1994\]. The spectrophotometric standard stars GD 108 and L745-46A were observed to enable us to correct the continuum shape of our spectra. Table 3 lists all observed standard stars with their spectral types (obtained from SIMBAD, operated by CDS, Strasbourg) and also comments on their use as Lick/IDS standard, velocity standard or spectrophotometric standard.
### 2.2 Basic data reduction
Most of the basic data reduction steps have been performed with packages under IRAF. For each night individually the science frames were overscan corrected and a bias frame was subtracted. A few bad columns were removed by linear interpolation. From several domeflats and skyflats a final flatfield accounting for the pixel-to-pixel variations and vignetting was constructed and applied to the frames. Cosmic rays were removed using the cleanest task in the REDUCEME package \[Cardiel et al. 1998\]. This task automatically detects and removes cosmic rays via a sophisticated deviation algorithm, while at the same time one can interactively inspect potential cosmic rays in sensitive areas such as close to the galaxy centre. The wavelength solution was determined from Th-Ar-lamp spectra which were taken before and after most of the science observations. The rms residual in the wavelength fit was typically 0.1 - 0.2 Å. Finally the sky was subtracted.
The central spectrum for each galaxy was extracted by fitting a low order polynomial to the position of the centre along the wavelength direction, re-sampling the data in the spatial direction and finally co-adding the spectra within a 5 pixels aperture yielding an effective aperture of 2$`\stackrel{}{.}`$3$`\times `$3$`\stackrel{}{.}`$85. Multiple exposures of the same galaxy were combined. The resulting S/N in the spectra ranges from $`30`$ \[Å<sup>-1</sup>\] for the faintest galaxies to more than 100 \[Å<sup>-1</sup>\] for the brightest ones (measured in a $``$100 Å wide region just bluewards of the Mg b feature). For stars we used the IRAF task apall to extract 1d-spectra. All galaxy and stellar spectra were logarithmically rebinned to a common wavelength range and increment. Finally the continuum shape of our spectra was corrected to a relative flux scale with the help of the spectrophotometric standard stars.
### 2.3 Kinematics
In order to correct the line-strength indices for velocity dispersion broadening and to construct index–$`\sigma _0`$ relations we need to measure the central velocity dispersion for each galaxy. Estimates were derived with the Fourier correlation quotient (FCQ, version 8) method \[Bender 1990, Bender, Saglia & Gerhard 1994\]. For the FCQ analysis the spectra were rebinned to twice the original spectral sampling and a wavelength range of 4876 to 5653 Å was extracted. Note that the $`\mathrm{H}\beta `$ feature is excluded from the wavelength range as it proved to be a source of severe template mismatch for galaxies with strong Balmer absorption. As we only consider central spectra in this paper we fit a pure Gaussian profile to the broadening function, neglecting higher order terms. To check the reliability of the FCQ analysis we used eight different G & K-giant template stars. For galaxies with a central velocity dispersion of $`\sigma _070`$ km s<sup>-1</sup>, all eight template stars give very similar results and an average value was adopted. The rms scatter between different template stars is 0.007 in log units for galaxies with $`\sigma _0100`$ km s<sup>-1</sup>. For galaxies with $`70\sigma _0<100`$ km s<sup>-1</sup> the rms scatter increases to 0.024 and for galaxies with $`\sigma _0<70`$ km s<sup>-1</sup> we find an rms scatter of 0.074. The uncertainty introduced by different template stars was comparable or larger than the internal error estimates of the FCQ program. Note, that for galaxies with $`\sigma _0<70`$ km s<sup>-1</sup> some template stars gave a poor fit to the broadening function and were excluded from the template sample. Only remaining measurements were averaged.
Using this procedure, velocity dispersions as low as $`50`$ km s<sup>-1</sup> could be recovered, although systematic errors will start to dominate for $`\sigma _0<90`$ km s<sup>-1</sup>. As our spectral resolution is rather low compared to velocity dispersions of $`5060`$ km s<sup>-1</sup> we emphasize that for these faint galaxies our velocity dispersions are only rough estimates. The final velocity dispersion errors for galaxies with $`\sigma 70`$ km s<sup>-1</sup> ($`\mathrm{\Delta }\mathrm{log}\sigma _0=0.022`$) were derived by a literature comparison (see Appendix, Figure 13). For galaxies with $`\sigma _0<70`$ km s<sup>-1</sup> we adopt the mean rms scatter of the template stars ($`\mathrm{\Delta }\mathrm{log}\sigma _0=0.074`$).
## 3 LICK/IDS CALIBRATION
The wavelength range of our spectra covers 16 different line-strength indices, such as Mg<sub>2</sub>, $`\mathrm{H}\beta `$ and $`\mathrm{H}\gamma _\mathrm{A}`$, in the Lick/IDS system which is described in detail in Worthey (1994, hereafter W94), Worthey & Ottaviani (1997, hereafter WO97) and Trager et al. (1998). In the following analysis we use an updated version of the W94 models which is available from Dr. G. Worthey’s home page. The updates affect only models where $`[\mathrm{Fe}/\mathrm{H}]1.0.`$ and are most noticable for the $`\mathrm{H}\beta `$ index. For a recent study of the behaviour of the Balmer indices at low metallicities see Poggianti & Barbaro (1997) and Maraston, Greggio & Thomas (1999). Before one can compare the measured indices with model predictions by e.g., W94 and Vazdekis et al. (1996, hereafter V96), the measurements have to be carefully calibrated to the Lick/IDS system. Generally there are three effects to account for: (a) the difference in the spectral resolution between the Lick/IDS system and our set-up, (b) the internal velocity broadening of the observed galaxies and (c) small systematic offsets caused by e.g., continuum shape differences.
(a) In order to account for diffrences in spectral resolution we broadened the spectra with a Gaussian of wavelength dependent width, such that the Lick/IDS resolution was best matched at each wavelength (see Figure 7 in WO97). After this step our spectra should resemble very well the general properties of the original spectra obtained by the Lick group.
(b) In a second step we need to correct the indices for velocity dispersion broadening. The observed spectrum of a galaxy is the convolution of the integrated spectrum of its stellar population(s) by the instrumental broadening and the distribution of line-of-sight velocities of the stars. These effects broaden the spectral features, in general reducing the observed line-strength compared to the intrinsic values. In order to compare the raw index measurements for galaxies with model predictions we calibrate the indices to zero velocity dispersion. Spectra of 15 different G9-K4 giant stars were first broadened to the Lick/IDS resolution and then further broadened using a Gaussian to $`\sigma =20360`$ km s<sup>-1</sup> in steps of 20 km s<sup>-1</sup>. The indices are then measured for each star and $`\sigma `$-bin and a correction factor, $`C(\sigma )`$, such that $`C(\sigma )=`$Index(0)/Index($`\sigma `$) is determined.
Figure 15 in the Appendix shows the dependence of the correction factor on $`\sigma `$ for all 16 indices. Note that for the molecular indices Mg<sub>1</sub> and Mg<sub>2</sub> and the index $`\mathrm{H}\gamma _\mathrm{F}`$<sup>2</sup><sup>2</sup>2This index is actually not a molecular index but typical index values are close to zero hence a correction factor can degenerate. the correction factor is defined as $`C(\sigma )=`$Index(0) – Index($`\sigma `$). The scatter in $`C(\sigma )`$ at 360 km s<sup>-1</sup> was $`<5\%`$ for all indices but $`\mathrm{H}\beta `$.
It is worth looking in detail why the $`\mathrm{H}\beta `$ velocity dispersion correction seems to be so insecure. The derived correction factors are only useful if the stars used for the simulations resemble the galaxy spectra. In principle one might expect a dependence of the correction factor on line-strength – but most indices do not show such a behaviour. In fact $`\mathrm{H}\beta `$ is the only index where we find a significant influence of line-strength on the correction factor at a given $`\sigma `$. It turns out that stars which exhibit $`\mathrm{H}\beta `$-absorption $`1.1`$ Å lead to correction factors of $`C(\sigma )<1.0`$ and stars with $`\mathrm{H}\beta `$-absorption $`>1.1`$ Å imply positive corrections. In the Fornax sample there are no galaxies with $`\mathrm{H}\beta `$ absorption line-strength of less than 1.4 Å, hence only stars with a $`\mathrm{H}\beta `$ index greater 1.1 Å have been used to evaluate the correction factor (scatter $`<5\%`$).
Another way to check the accuracy of the velocity dispersion corrections is to use galaxy spectra with small internal velocity dispersions as templates and treat them in the same way as stars. The galaxies NGC1373, NGC1380A, NGC1336, IC1963 and ESO358-G59 were used for this purpose. They span a range in $`\mathrm{H}\beta `$ absorption of $`1.73`$ Å and a range in central velocity dispersion of $`\sigma _0=5496`$ km s<sup>-1</sup>. In Figure 15 the galaxies are represented by open circles and they agree very well with the stellar correction for most of the indices. As expected for $`\mathrm{H}\beta `$, the galaxies match the results from stars with a $`\mathrm{H}\beta `$ absorption $`>1.1`$ Å.
The final correction factors are derived by taking the mean of 15 stars and the five galaxies in each $`\sigma `$-bin (solid line in Figure 15). The velocity dispersion corrections are applied by a FORTRAN program which reads in the raw index-measurements from continuum corrected and resolution corrected galaxy spectra. For each galaxy and index it applies a correction for velocity dispersion. The program linearly interpolates between $`\sigma `$-bins and also adds the error from the velocity-dispersion correction factor to the raw Poisson error of the spectra. As the error in the correction factor is much bigger than any error caused by uncertainties in $`\sigma `$, we assumed the velocity dispersion of the galaxies to be error free.
(c) Although we have matched very well the spectral resolution of the Lick system, small systematic offsets of the indices introduced by continuum shape differences are generally present (note that the original Lick/IDS spectra are not flux calibrated). To establish these offsets we compared our measurements for stars in common with the Lick/IDS stellar library. In total we observed 13 different Lick/IDS stars. Figure 14 in the Appendix shows the difference between Lick/IDS measurements and ours after the mean offset has been removed. The mean offsets and associated errors for each index are summarized in Table 4. The star HD221148 was excluded from the offset analysis because our index measurements proved to be very different from the original Lick/IDS measurements – possibly due to its variable nature (see Table 3). The formal error in the offset is evaluated by the mean standard deviation of stars with respect to the mean offset divided by $`\sqrt{n_{stars}1}`$.
Most off the indices show small offsets to the Lick/IDS system, similar to the ones quoted in (WO97, Table 9). The rather large offset in Mg<sub>2</sub> is due to a well known difference in continuum shape.
Recently Trager et al. (1998) published the Lick/IDS library of extragalactic objects including 7 galaxies in the Fornax cluster and NGC3379. We can check our previous offset evaluation by comparing our galaxy measurements with Trager et al. For this purpose we extracted a 3 pixels central aperture (2$`\stackrel{}{.}`$3$`\times `$2$`\stackrel{}{.}`$44) for our galaxies matching the Lick/IDS standard aperture of 1$`\stackrel{}{.}`$4$`\times `$4″. Our indices are then corrected for velocity dispersion as described in paragraph (b) and the offsets from Table 4 are applied. The results are overplotted in Figure 14 in the Appendix (filled symbols). The galaxies show for all indices more scatter around the mean offset than the stars which is somewhat reflected in the bigger error bars, but there are also some outliers. This is not surprising as seeing effects and aperture differences will introduce some non-reproduceable offsets for individual galaxies. Furthermore we note that the Lick group had to observe the Fornax galaxies at a very high airmass. With the possible exception of the indices G4300 and Fe4383 the offset inferred from the galaxy comparison are consistent with the stellar comparison.
The offsets listed in Table 4 were applied to all indices after the correction for velocity dispersion. Note that the Lick/IDS-system offset-error is a constant value and does not depend on the velocity dispersion of the galaxy itself. Therefore we did not include this error in the individual index errors but rather quote for each index a common offset error (see also Table 4). The final corrected central (2$`\stackrel{}{.}`$$`\times `$3$`\stackrel{}{.}`$85) index measurements and associated errors for the Fornax galaxies and NGC3379 are presented in Table D2 in the Appendix. For each galaxy we give the index measurement in the first row and the 1$`\sigma `$ error in the second row.
Note that for the galaxies NGC1381 and NGC1427 we combined three exposures yielding a very high S/N spectrum. Here our index-error estimation taking into account only the Poisson error becomes invalid because of other error sources such as the wavelength calibration, continuum correction and aperture effects. By comparing individual exposures we established that 1.5 $`\times `$ the original Poisson error estimate is a good indicator of the true error. This adjusted error was adopted in Table D2 and for any further analysis.
## 4 A consistency test of the model predictions
For the following analysis of the nuclear stellar populations (Section 5) it is extremely important that our index measurements are accurately calibrated onto the Lick/IDS system which is based on the Lick/IDS stellar library \[Worthey et al. 1994\]. Here we investigate the accuracy and consistency of our calibration and the model predictions by presenting index–index plots which are almost degenerate in age and metallicity. In that way the model predictions cover only a small “band” of the parameter space and they should trace the relation of the galaxies if the models describe accurately the galaxy properties and our calibration is accurate. Figure 1 shows the relation between $`<`$Fe$`>`$<sup>3</sup><sup>3</sup>3$`<`$Fe$`>`$=(Fe5270+Fe5335)/2, Fe5015 and Fe5406 for our sample of galaxies (filled circles) and the original Lick/IDS sample of galaxies (Trager et al. 1998, small dots). Overplotted are models by W94 (black lines) and V96 (grey lines). The plots show a good agreement between index measurements and the model predictions. The reduced Poisson noise of our data set compared to the Lick/IDS measurements can be clearly seen (see also Figure caption). We note that the model predictions of Vazdekis and Worthey are in good agreement.
A similar analysis of the three Mg indices is shown in Figure 2. Here we find a significant deviation of the measured index values compared to the model predictions for metal rich and/or old stellar populations. The deviations are seen in the Fornax sample as well as in the original Lick/IDS galaxy spectra (see also Worthey 1992, Figure 5.12 & 5.13). We therefore note that this discrepancy is inherent to the Lick/IDS system & models and any models which use the Lick/IDS fitting functions are likely to show the same offset. In Figure 3 we present the Balmer indices $`\mathrm{H}\beta `$, $`\mathrm{H}\gamma _\mathrm{A}`$ and $`\mathrm{H}\gamma _\mathrm{F}`$. Here we find generally good agreement with small deviations between model predictions and data at low values of $`\mathrm{H}\gamma _\mathrm{A}`$ vs $`\mathrm{H}\gamma _\mathrm{F}`$ which is present in our data and the original Lick/IDS measurements.
Figures 1 to 3 suggest that our Lick/IDS calibration is very consistent with the original galaxy measurements of the Lick group. However, we note that small, systematic offsets exist between the parameter space covered by galaxies and the model predictions for Magnesium at high index values and for Balmer lines at low index values.
## 5 THE NUCLEAR STELLAR POPULATIONS
The aim of this section is to derive estimates of the mean (luminosity weighted) ages and metal abundances of early-type galaxies in the Fornax cluster. As pointed out by W94, the determination of the ages and metallicities of old stellar populations is complicated by the similar effects that age and metallicity have on the integrated spectral energy distributions. However, this degeneracy can be partially broken by plotting a particular age sensitive index, such as one of the Balmer line indices, against a more metallicity sensitive index. The usefulness of this approach has been demonstrated by many authors \[González 1993, Fisher, Franx & Illingworth 1995, Mehlert 1998, Kuntschner & Davies 1998, Jørgensen 1999\]. However, as we will see in this Section, among other issues the treatment of non-solar abundance ratios is a crucial parameter in the determination of absolute age and metallicity estimates. We will also investigate the effects of nebular emission and composite stellar populations on age/metallicity estimates in Sections 5.2 & 5.3 respectively, before we present our best age/metallicity estimates of the Fornax early-type galaxies in Section 5.4.
### 5.1 Non-solar abundance ratios
In Figure 4 we present age/metallicity diagnostic diagrams of six metallicity sensitive indices (Mg<sub>2</sub>, C<sub>2</sub>4668, Ca4455, Fe3, $`<`$Fe$`>`$ & Fe5406) plotted against the age sensitive Balmer line indices $`\mathrm{H}\beta `$ and $`\mathrm{H}\gamma _\mathrm{A}`$ (the new index Fe3 is defined in Equation 1). Figure 4h is a reproduction from Kuntschner & Davies (1998) with minor data up-dates. Overplotted are model predictions from W94, WO97 (black lines) and V96 (grey lines). Solid lines represent lines of constant age and the dashed lines are lines of constant metallicity. The Worthey models span a range in age of 1.5–5 Gyr with \[Fe/H\]=-0.225 to 0.5 and 8–17 Gyr with \[Fe/H\]=-2 to 0.5. The V96 models span a range in age of 1–17.4 Gyr with \[Fe/H\]=-0.7 to 0.4. The direction of increasing age and metallicity is indicated in Figure 4a by arrows.
Our previous result from Kuntschner & Davies (1998), that Fornax ellipticals form a sequence in metallicity at old ages and that the S0s spread to younger ages is confirmed in all diagrams. Examining the diagrams in detail one can see that the mean age and metallicity of the sample changes from diagram to diagram; e.g., the ellipticals appear older and more metal poor in the $`<`$Fe$`>`$ vs $`\mathrm{H}\gamma _\mathrm{A}`$ diagrams compared to the Mg<sub>2</sub> vs $`\mathrm{H}\gamma _\mathrm{A}`$ diagram. This effect was previously reported and recently reviewed by Worthey (1998). It is now widely accepted that this discrepancy in the model predictions is caused by non-solar abundance ratio effects. For example Mg as measured by the Mg<sub>2</sub> index is overabundant compared to Fe in luminous elliptical galaxies, i.e., $`[Mg/Fe]>0`$ \[O’Connell 1976, Peletier 1989, Worthey, Faber & González 1992, Davies, Sadler & Peletier 1993, Weiss, Peletier & Matteucci 1995, Jørgensen 1997, Jørgensen 1999\].
The Mg overabundance can be examined in a Mg-index vs Fe-index plot \[Worthey, Faber & González 1992\]. In such a diagram the model predictions cover only a narrow band in the parameter space as effects of age and metallicity are degenerate. Figure 5 shows plots of $`<`$Fe$`>`$ & Fe5270 vs Mg<sub>2</sub> for the Fornax sample. Overplotted are model predictions from W94, V96 and Weiss et al. (1995). We assume that the models reflect solar abundance ratios if not stated otherwise, i.e., $`[Mg/Fe]=0`$. If the model predictions accurately resemble the galaxy properties they should trace the observed relation. The measured line-strength of most of the S0s agrees with the model predictions, perhaps 3–4 galaxies having slightly low Mg<sub>2</sub> absorption compared to Fe5270 & $`<`$Fe$`>`$. However, for most of the ellipticals and the S0 NGC1380, the models predict too little Mg-absorption at a given Fe-absorption strength. Additionally the most metal rich galaxies are the furthest away from the model grids. Using the Mg-overabundance correction by Greggio (1997, see Figure 5a) and the models for \[Mg/Fe\]=0.45 by Weiss et al. (1995), we conclude that the stellar populations of Fornax ellipticals and the bulge of NGC1380 are Mg-overabundant relative compared to Fe. The overabundance ranges between \[Mg/Fe\]=0.0 to $``$0.4. We note that there exists considerable spread in overabundance at a given Fe-line strength in our sample.
Non-solar abundance ratios do exist not only in elliptical galaxies but also in our own galaxy where stars show an overabundance for $`\alpha `$-elements<sup>4</sup><sup>4</sup>4$`\alpha `$ includes the elements O, Mg, Si, S, Ca, and Ti at \[Fe/H\]$``$0.0 \[Edvardsson et al. 1993, McWilliam 1997\]. Of course if those stars are incorporated in a stellar library which in turn is used for model predictions of integrated stellar populations, the predictions will be somewhat $`\alpha `$-element overabundant at low metallicities. We therefore note that models which use the Lick/IDS fitting functions are probably $`\alpha `$-element overabundant at low metallicities which makes it more difficult to interpret trends in diagrams such as Figure 5.
Several indices covered by our wavelength range show deviations from the model predictions when compared to the average Fe index: Mg<sub>1</sub>, Mg<sub>2</sub>, Mg b, Fe5709 & C<sub>2</sub>4668. Fe5709 is a very weak index and its correction for velocity dispersion broadening may well be insecure, so we cannot draw any firm conclusions. C<sub>2</sub>4668 is an important index because it shows the strongest total metallicity sensitivity in the Lick/IDS system \[Worthey 1998\] and is therefore preferentially used in age/metallicity diagnostic diagrams. In Figure 6 we present a plot of C<sub>2</sub>4668 vs Fe3. Fe3 is a combination of three prominent Fe lines, thus maximizing its sensitivity to Fe while minimizing the Poisson errors:
$$\mathrm{F}e3=\frac{\mathrm{F}e4383+\mathrm{F}e5270+\mathrm{F}e5335}{3}$$
(1)
As a consequence of the extreme metallicity sensitivity of C<sub>2</sub>4668, the models are not as degenerate as in the previous plots. Nevertheless it is clear that for a C<sub>2</sub>4668 absorption strength in excess of $``$6 Å the model predictions do not follow the observed trend (see also Kuntschner 1998). Hence, we conclude that C<sub>2</sub>4668, or better at least one of the species contributing to the index, is overabundant compared to Fe in metal rich Fornax galaxies. Can this overabundance be caused by Mg as seen in the Fe vs Mg plot (Figure 5)? Due to the proximity of metal absorption lines in the optical wavelength region non of the Lick/IDS indices measures the abundance of only a particular element such as Fe or Mg. There are always contributions from other elements or molecules to an index \[Tripicco & Bell 1995\]. In particular, the C<sub>2</sub>4668 index has a relatively wide central bandpass (86.25 Å) including a wide range of metal lines. The most dominant species here is carbon in form of C<sub>2</sub>-bands which blanket the central bandpass \[Tripicco & Bell 1995\]. Yet, more important here is that according to Tripicco & Bell (1995; Table 6, cool giants) the C<sub>2</sub>4668 index decreases when the Mg-abundance (or Oxygen-abundance) is increased (at fixed abundances of all other elements). Therefore the ’overabundance’ of C<sub>2</sub>4668 cannot be caused by Mg. What exactly drives the overabundance of C<sub>2</sub>4668 compared to Fe remains to be seen.
The overabundance of certain elements compared to Fe has a profound effect on the use of age/metallicity diagnostic diagrams if the model predictions reflect only solar abundance ratios. Recalling Figure 4f (see also Kuntschner & Davies 1998) we find that not only are the metallicities (measured as \[Fe/H\]) overestimated due to the C<sub>2</sub>4668 overabundance but furthermore much of the age upturn at high metallicities is caused by the overabundance. The effect of changing the age estimates is caused by the residual age/metallicity degeneracy which is still present in all diagrams in Figure 4. Only if the index-combination breaks the degeneracy completely, i.e., if lines of constant age and constant metallicity are perpendicular, would the age estimates not be affected by non-solar abundance ratios. We further note that trends in abundance ratios within a dataset such as our Fornax sample (increasing Mg/Fe with galaxy mass) can lead to artificial relative age trends in diagrams such as in Figure 4f. Taking into account not only the uncertainties introduced by non-solar abundance ratios but also other model parameters such as which isochrone library to use, it seems very insecure to derive absolute age estimates from the currently available stellar population models.
Introducing non-solar abundance ratios in model predictions is rather complicated, as accurate model predictions do need a stellar library covering the whole parameter space of T<sub>e</sub>, log g, \[Fe/H\] and \[Mg/Fe\]. Furthermore, new isochrone calculations may well be needed for each \[Mg/Fe\] bin: recently Salaris & Weiss (1998) suggested that scaled-solar isochrones cannot be used to replace Mg-enhanced ones at the same total metallicity. The latter will not only change the model predictions for indices such as Mg<sub>2</sub> but may affect all indices and in particular the age sensitive ones such as $`\mathrm{H}\beta `$ and $`\mathrm{H}\gamma _\mathrm{A}`$ (see also Worthey 1998). However, note that Weiss, Peletier & Matteucci (1995) concluded in their study that scaled solar isochrones are sufficient to calculate model predictions for non-solar abundance ratios.
Another way to examine non-solar abundance ratios is to compare the metallicity estimates derived from different metal lines using the same age indicator. Figure 7 compares the metallicity estimates taken from Mg<sub>2</sub>, C<sub>2</sub>4668, Fe5406, Ca4455 vs $`\mathrm{H}\beta `$ diagrams with the estimates taken from a Fe3 vs $`\mathrm{H}\beta `$ diagram. Here Fe3 serves as our mean Fe-abundance indicator. The metallicity estimates are derived from the V96 models<sup>5</sup><sup>5</sup>5These models have a bimodal IMF which is very similar to Salpeter for $`\mathrm{M}>0.6\mathrm{M}_{}`$. Note that for an age of $``$17 Gyrs V96 models predict $`0.10.2`$ Å less $`\mathrm{H}\beta `$-absorption compared to W94 models. This of course will affect the absolute age estimates but has little affect on the metallicity estimates.. To get more accurate estimates, the age/metallicity–grid was expanded to a step size of 0.025 in \[Fe/H\] by linear interpolation. Furthermore the diagram was extrapolated to \[Fe/H\]=+0.7 by linear extrapolation. The age range of 1 to 17.4 Gyrs is covered by 18 grid points. Errors on the metallicity estimates were derived by adding and subtracting the index error for each galaxy individually (Poisson error and Lick/IDS offset error added in quadrature) and re-deriving the metallicity estimates. The final uncertainty displayed in Figure 7 was taken to be 0.7 times the maximum change in \[Fe/H\].
In panel (a) of Figure 7 we can clearly see that for elliptical galaxies Mg<sub>2</sub> gives metallicity estimates which are larger than those derived from Fe3 and there is a trend that the Mg-overabundance increases with increasing metallicity. Most of the S0s are consistent with solar or slightly less than solar abundance ratios of Mg/Fe. However, the (more luminous) S0s NGC1380 & NGC1381 show a weak overabundance of Mg. The index C<sub>2</sub>4668 (panel b) gives on average high metallicity estimates compared to Fe3. Although three galaxies with just above solar metallicity show solar abundance ratios. As expected the Fe-index Fe5406 (panel c) is in good agreement with the estimates derived from Fe3. The Ca4455 index (panel d) gives marginally higher metallicity estimates compared to the Fe3 indicator. We note that the Ca4455 index is more sensitive to a mix of heavy elements than to Calcium on its own despite its name \[Tripicco & Bell 1995\].
In conclusion we can confirm our previous results that Mg and C<sub>2</sub>4668 are overabundant compared to Fe. The Mg-overabundance follows a trend where metal rich (and luminous) Es show a stronger overabundance than less luminous and metal poor galaxies.
### 5.2 Nebular emission in early-type galaxies
So far we have concentrated on breaking the age/metallicity degeneracy and the treatment of non-solar abundance ratios. A further important issue when estimating ages and metallicities from line strength indices is nebular emission. Elliptical galaxies normally contain much less dust and ionized gas than spirals, in fact, for a long time they were regarded as dust and gas free. However, spectroscopic surveys of large samples of early-type galaxies revealed that about 50-60% of the galaxies show weak optical emission lines \[Phillips et al. 1986, Caldwell 1984\]. Typically the reported strength of emission lines such as \[OII\], \[H$`\alpha `$\] and \[NII\] $`\lambda 6584`$ indicates the presence of only $`10^310^5M_{}`$ of warm ionized gas in the centre. A more recent study of 56 bright elliptical galaxies by Goudfrooij et al. (1994) detected ionized gas in 57% of their sample and confirms the amount of ionized gas present. Additionally, HST images of nearby bright early-type galaxies revealed that approximately 70-80% show dust features in the nucleus \[van Dokkum & Franx 1995\]. Stellar absorption line-strength measurements can be severely affected if there is emission present in the galaxy which weakens the stellar absorption \[Goudfrooij & Emsellem 1996\]. For example, nebular $`\mathrm{H}\beta `$-emission on top of the integrated stellar $`\mathrm{H}\beta `$-absorption weakens the $`\mathrm{H}\beta `$-index and leads therefore to wrong, i.e., too old age estimates.
The spectrum of ESO358-G25 shows clear emission in $`\mathrm{H}\beta `$ and H$`\gamma `$ along with weak \[O III\] emission (see Kuntschner & Davies 1998, Figure 3). As a consequence, the age is overestimated in Figures 4 (a)–(f). The arrow attached to ESO358-G25 indicates a rough emission correction. However, it is extremely difficult to accurately correct the $`\mathrm{H}\beta `$-index in individual galaxies for emission contamination. A much better method to reduce emission contamination is to use higher order Balmer lines such as $`\mathrm{H}\gamma `$ as they are less affected by nebular emission \[Osterbrock 1989\]. Indeed, in Figures 4 (g)–(l) the galaxy ESO358-G25 moves to much younger ages. As none of the other galaxies move significantly to younger ages we conclude that nebular emission is not very prominent in our Fornax sample. This is supported by the absence of strong \[OIII\]$`\lambda 5007`$ emission. Only 5 galaxies show emission above our detection limit of $``$0.2 Å. The strongest emission is detected in ESO358-G25 with 0.7 Å equivalent width (for details see Kuntschner 1998).
### 5.3 Effects of composite stellar populations
Most of the S0s in our sample have luminosity weighted young stellar populations with some of them having also high metallicities when compared to single-burst stellar population (SSP) models. However, these galaxies show only a central young stellar population on top of an underlying older one as opposed to be entirely young. It is not straight forward to compare composite stellar populations with SSP models \[de Jong and Davies 1997\]. So, how reliable are the age, metallicity and abundance ratio estimates taken directly from SSP models for these young galaxies? In order to explore this issue we calculated model predictions based on V96 for simple composite stellar populations. Two representative tracks are shown in Figure 8 for two age/metallicity diagnostic diagrams and a plot of Fe3 vs Mg<sub>2</sub> in order to explore the behaviour of abundance ratios: Model A is a 15 Gyr old (90% in mass) stellar population plus a burst (10% in mass) of varying ages from 0.1 Gyr to $`3`$ Gyrs. Both populations have solar metallicity. For Model B we reduced the burst fraction to 1% (in mass) while the other parameters are the same as in Model A.
Overall one can see in the age/metallicity diagnostic diagrams that for a short time the burst population will dominate the integrated light leading to strong $`\mathrm{H}\beta `$ absorption and weak metal-line absorption. Then the underlying old population becomes more and more important and after $``$3 Gyrs the galaxy is almost back to its original place in the diagram. However, the burst strength influences the exact track which the galaxy takes in the diagram. For a burst of 10% or 20% (not shown) in mass the tracks follow roughly the solar metallicity line in the normal SSP models. Yet, for a small burst (1% in mass) the integrated light looks for a short while as having metallicities well above solar. This effect is more pronounced for Mg<sub>2</sub> than for Fe3. Of course, this in turn leeds to an artificially created overabundance when these galaxies are compared to SSP models (see Figure 8c). For bursts stronger than a few percent the abundance ratios are not significantly affected.
In summary we find that composite stellar populations and in particular small (in mass) bursts, such as used in our simple models, can lead to an overestimation of the metallicity in the context of SSP models. Abundance ratios can be affected in the sense that the Mg/Fe ratio is too strong. Our model calculations show that these conclusions qualitatively hold if the metallicity is changed or different metallicities are combined. A more thorough investigation of these issues would be very valuable but is beyond the scope of this paper (see Hau, Carter & Balcells (1999) for a more detailed analysis).
### 5.4 Best age and metallicity estimates
Having examined some of the fundamental problems with applying stellar population model predictions to observed line-strength indices we present in Figure 9 what we consider our best age/metallicity diagnostic diagram. A mean Fe-index (Fe3) is plotted against an emission robust higher order Balmer line ($`\mathrm{H}\gamma _\mathrm{A}`$). Due to the lack of model predictions with non-solar abundance ratios we decided to avoid indices which are affected by overabundance problems (e.g., Mg & C<sub>2</sub>4668). Instead we use here a combination of Fe-indices (Fe3) as metal indicator which will bias our results towards the Fe abundance. We note however, that our metallicities are not to be understood as total metallicity but rather as a good estimate of the Fe-abundance. Any non-solar abundance ratios which affect $`\mathrm{H}\gamma _\mathrm{A}`$ are ignored. Model predictions by W94 and V96 are overplotted in Figure 9.
The ellipticals form a sequence of metallicity at roughly constant age. The centres of the bright S0s NGC1380 & NGC1381 follow the sequence of Es. The remaining S0s cover a large range in metallicity and spread to much younger luminosity weighted ages than the Es. We emphasize that these age and metallicity estimates are central luminosity weighted estimates and for apparently young galaxies the derived parameters are somewhat more insecure (see previous discussion about the effects of composite stellar populations). The age and metallicity gradients within the galaxies will be discussed in a future paper. Fornax A, a bright peculiar S0, shows strong Balmer lines and strong Fe absorption which translates into a luminosity weighted young and metal rich stellar population. As we will see in the next section all other young or metal poor S0s have velocity dispersions of $`\sigma _070`$ km s<sup>-1</sup>. The two galaxies with the weakest metal lines and strong $`\mathrm{H}\beta `$ & $`\mathrm{H}\gamma _\mathrm{A}`$ absorption (ESO359-G02, cross and ESO358-G25, open triangle) appear to be different from the rest of the sample. These galaxies are likely to be post-starburst or starburst galaxies respectively. They have remarkable spectra for early-type galaxies, showing blue continua, strong Balmer lines, and weak metal lines. These galaxies are amongst the faintest in our sample and are $``$3° away from the centre of the cluster (see also Kuntschner & Davies 1998).
## 6 LINE–STRENGTH INDICES AND THE CENTRAL VELOCITY DISPERSION
The central velocity dispersion $`\sigma _0`$ of early-type galaxies is known to correlate strongly with colours \[Bower, Lucey & Ellis 1992\] and the absorption strength of the Mg-absorption feature at 5174 Å \[Terlevich et al. 1981, Burstein et al. 1988, Bender, Burstein & Faber 1993, Jørgensen 1997, Colless et al. 1999\]. The relatively small scatter about these relations imply that the dynamical properties of galaxy cores are closely connected with their stellar populations. However, analysing the Mg–$`\sigma _0`$ relation for a sample of 736 mostly early-type galaxies in 84 clusters, the EFAR group \[Colless et al. 1999\] finds a rather large dispersions in age (40%) and in metallicity (50%) at fixed velocity dispersion using the constraints from the Mg–$`\sigma _0`$ relation and the Fundamental Plane. Correlations of other metal indices, such as $`<`$Fe$`>`$, with the central velocity dispersion have long been expected but so far relations have shown a large scatter and only weak correlations \[Fisher, Franx & Illingworth 1996, Jørgensen 1997, Jørgensen 1999\]. However, we will demonstrate that galaxies in the Fornax cluster do show a clear correlation between Fe-indices and central velocity dispersion.
Following Colless et al. (1999) we find it more convenient to express the “atomic” indices in magnitudes like the “molecular” index Mg<sub>2</sub>. The new index is denoted by the index name followed by a prime sign \[\], e.g., Mg b. Note that by using only the logarithm of the atomic index, one introduces a non linear term in comparison to the magnitude definition. Furthermore negative index values such as for the $`\mathrm{H}\gamma _\mathrm{A}`$ index cannot be put on a simple logarithmic scale. A priori it is not clear whether $`\mathrm{log}`$ index or index$`^{}`$ correlates better with $`\mathrm{log}\sigma _0`$, but as the classical Mg–$`\sigma _0`$ relation was established with Mg<sub>2</sub> measured in mag we adopt this approach here for all other indices as well. The conversion between an index measured in Å and magnitudes is
$$\mathrm{i}ndex^{}=2.5\mathrm{log}\left(1\frac{\mathrm{i}ndex}{\mathrm{\Delta }\lambda }\right)$$
(2)
where $`\mathrm{\Delta }\lambda `$ is the width of the index bandpass (see e.g., WO97 and Trager et al. 1998 for a list of bandpass definitions). Fe3 is defined as
$$\mathrm{F}e3^{}=\frac{\mathrm{F}e4383^{}+\mathrm{F}e5270^{}+\mathrm{F}e5335^{}}{3}.$$
(3)
Figure 10 shows index–$`\sigma _0`$ relations for eight different metal indices and two Balmer-line indices. The best fitting linear relations and the scatter are summarized in Table 5 for all indices considered in this paper. For the fits we used an ordinary least square method, minimizing the residuals in y-direction \[Isobe et al. 1990, hereafter OLS(Y$`|`$X)\]. Included in the fit are all galaxies with old stellar populations, i.e., all Es plus the bright S0s NGC1380 & NGC1381; in total 13 galaxies. The 1-$`\sigma `$ scatter around the relation was robustly estimated by deriving a value which includes 9 out of 13 galaxies (69%). A correlation coefficient derived from a (non parametric) Spearman rank-order test is given in the lower right corner of each panel in Figure 10. The probability that the parameters are not correlated is given in brackets.
For the galaxies with old stellar populations the Mg–$`\sigma _0`$ relation is in excellent agreement with the literature (Jørgensen 1997; Colless et al. 1999). Remarkably, the Fe-line-indices also show a clear positive correlation with the central velocity dispersion and little scatter. This is the first time such strong correlations have been found at a significant level. We note that all the Fe-line and Ca4455–$`\sigma _0`$ relations show a slope consistent with a value of $`0.035`$. In contrast the slope of the Mg-lines and C<sub>2</sub>4668 are significantly steeper (see dot-dashed line in Figure 10 and Table 5).
Although our Mg<sub>2</sub>$`\sigma _0`$ relation agrees well with the literature values we find significant differences for other $`\mathrm{log}(index)`$$`\sigma _0`$ relations compared to the data of Jørgensen (1997,1999). Table 6 shows a comparsion of the slopes. The $`\mathrm{log}(<`$Fe$`>)`$$`\sigma _0`$ relation seems to be far steeper in the Fornax cluster whereas the $`\mathrm{log}\mathrm{H}\beta _G`$$`\sigma _0`$ relation is shallower compared to Coma. The $`\mathrm{log}\mathrm{C}_24668`$$`\sigma _0`$ relation in Fornax is marginally consistent with Jørgensen (1997). It is not clear why the $`\mathrm{log}(index)`$$`\sigma _0`$ relations for H$`\beta _G`$ and $`<`$Fe$`>`$ should be different to the Coma cluster. We will present a possible explanantion at the end of this Section and in Section 8 where we discuss our results.
The centres of the two bright and old S0s NGC1380 & NGC1381 follow generally well the relation set by the elliptical galaxies. The lower luminosity S0s have velocity dispersions $`\sigma _070`$ km s<sup>-1</sup> and show a large scatter about the mean relation of the old galaxies. However, it is worth noting that they exhibit generally weak Mg absorption and some of the faint S0s show as much Fe absorption as L ellipticals. Fornax A, the brightest galaxy in our sample has a central velocity dispersion of $`\sigma _0220`$ km s<sup>-1</sup> which is too low compared to ellipticals of this luminosity in the Faber-Jackson relation (see Figure 11). It also departs significantly from the Mg–$`\sigma _0`$ relations in the sense that it shows too weak Mg-absorption. As Fornax A is regarded as the product of a recent merger \[Schweizer 1980, Schweizer 1981, Mackie & Fabbiano 1998\] we interpret our results as strong indications of at least one young stellar component in this galaxy.
One would expect the young stars in this galaxy to produce strong Balmer absorption lines (as seen in Figure 9) and to dilute (or weaken) the metal lines of the underlying older stellar component. However, if the burst mass is not too small, the relative abundances of metal lines should in first order not be affected (see discussion of composite stellar populations in Section 5.3). Yet, we find that Fornax A deviates only from the Mg–$`\sigma _0`$ relation and not from any of the other metal-index–$`\sigma _0`$ relations (Figure 10). We interpret this as good evidence that the underlying older stellar population of Fornax A is significantly different from ellipticals at this velocity dispersion, i.e.,the \[Mg/Fe\] ratio is lower, close to solar.
Two of the ellipticals stand out from the normal metal index–$`\sigma _0`$ relation: NGC1373 and IC2006 (labelled in Figure 10). These galaxies always show stronger metal line absorption than what would be expected from the mean relation. This is most prominent in the Fe3$`\sigma _0`$ diagram (panel c). We note however, that the galaxies follow the mean Faber-Jackson relation (Figure 11). There is little known about the galaxy NGC1373; perhaps the best explanation why this (elliptical) galaxy is somewhat off the mean relation is to regard it as a transition galaxy between the sequence of Es and the faint S0s. However, IC2006 has been studied in detail by Schweizer et al. (1989). They found a large counter-rotating ring of neutral hydrogen (HI) associated with faint optical features and suggest that the HI ring may have formed during a merger which created IC2006. Franx et al. (1994) re-analysed the optical photometry of Schweizer et al. taking into account the inclination of the galaxy and concluded that it probably has a large disc in the outer parts which is seen almost face on and therefore difficult to detect. They suggest that it should be classified as E/S0 rather than a bona fide elliptical.
It seems plausible that the (perhaps peculiar) merger history of this galaxy is the reason for its deviation from the index–$`\sigma _0`$ relations. However from our data, it is not clear whether the stellar populations of IC2006 are too metal rich or whether the central velocity dispersion is reduced compared to other elliptical galaxies of this mass. If indeed this type of galaxy is more frequent in other clusters, such as the Coma cluster (see Jørgensen 1999), it would explain why previous authors did not find a clear correlation of Fe-lines with $`\sigma _0`$. A detailed analysis of the kinematics and stellar population of this galaxy could be very valuable for our understanding of how todays early-type galaxies were created.
Panels (i) & (j) in Figure 10 show the index–$`\sigma _0`$ relations for two Balmer lines. Both indices show negative correlations. Elliptical galaxies and the bulges of NGC1380 and NGC1381 show little spread around the mean relation whereas the younger galaxies, most remarkably NGC1316, tend to have significantly stronger Balmer absorption at a given $`\sigma _0`$. We emphasize here that the slope in the relation of the galaxies with old stellar populations is mainly caused by a metallicity effect (metal poorer galaxies have stronger Balmer absorption) and has little to do with age differences. The two “metal-rich” galaxies IC2006 and NGC1373 are deviant from the main $`\mathrm{H}\gamma _\mathrm{A}`$$`\sigma _0`$ relation in the sense of lower $`\mathrm{H}\gamma _\mathrm{A}`$ line strengths. This is caused by the residual metallicity sensitivity of $`\mathrm{H}\gamma _\mathrm{A}`$. The side-bands of this index are located on metal lines which lower the pseudo-continuum level and thus weaken the index.
## 7 Global relations
In this section we investigate the relations between our best age, metallicity, \[Mg/Fe\] estimates and the central velocity dispersions. Figure 12 presents the results. The ages and metallicities were estimated from a Fe3 vs $`\mathrm{H}\gamma _\mathrm{A}`$ age/metallicity diagnostic diagram (Figure 9) in combination with V96 models. The errors are evaluated following the procedure outlined in Section 5.1 but only including the Poisson error for individual galaxies. Some of the young galaxies are at the edge or outside the range of the model predictions which prevents an accurate error evaluation. For the latter galaxies we do not plot error bars. Notice that the ages and metallicities are derived parameters which carry all the caveats discussed in the previous sections. For example the independent measurement errors of the line-strength indices translate into correlated errors in the age – metallicity plane due to the residual age/metallicity degeneracy in the Fe3 vs $`\mathrm{H}\gamma _\mathrm{A}`$ diagram. We note that the results presented in the following paragraphs would not change significanlty if $`\mathrm{H}\beta `$ is used as an age indicator (see Figure 16 and 17 in the Appendix for a comparison of the age and metallicity estimates derived from $`\mathrm{H}\gamma _\mathrm{A}`$ & $`\mathrm{H}\beta `$). The Mg-overabundance is estimated by evaluating the difference in metallicity estimate between a Mg<sub>2</sub>$`\mathrm{H}\beta `$ and a Fe3–$`\mathrm{H}\beta `$ diagram (see Figure 18 in the Appendix for an estimation of \[Mg/Fe\] using $`\mathrm{H}\gamma _\mathrm{A}`$).
Our age estimates of ellipticals do not show a significant correlation with $`\mathrm{log}\sigma _0`$ (panel a). With the exception of Fornax A all galaxies with $`\sigma _0>70`$ km s<sup>-1</sup> show roughly the same age whereas the younger galaxies populate the low velocity dispersion range. However, there is a hint that the two dynamically hottest galaxies are younger than their smaller brethren.
For galaxies with old stellar populations there is a clear correlation between the central metallicity and the central velocity dispersion $`\sigma _0`$ (panel b). Consistent with our findings for the Fe3–$`\sigma _0`$ relation the galaxies IC2006 and NGC1373 show a stronger metal content than what would be expected from the mean relation. The young S0s spread over the whole metallicity range. The best fitting OLS(Y$`|`$X) relation (solid line, Jack-Knife error analysis) to galaxies with old stellar populations gives:
$$[Fe/H]=(0.56\pm 0.20)\mathrm{log}\sigma (1.12\pm 0.46).$$
(4)
A correlation coefficient derived from a Spearman rank-order test (including all ellipticals and the two large S0s) is given in the lower right corner of each panel in Figure 12. The probability that the parameters are not correlated is given in brackets. Excluding NGC1373 and IC2006 from the fit gives the following relation (dashed line in panel b):
$$[Fe/H]=(0.82\pm 0.18)\mathrm{log}\sigma (1.72\pm 0.40).$$
(5)
In the age-metallicity plane (panel c) we find a statistically significant relation in the sense that the more metal rich (and also more luminous) galaxies are younger. The slope of this relation is similar to what Jørgensen (1999) found for the Coma cluster (see also Worthey, Trager & Faber 1995), yet the Fornax galaxies with velocity dispersion $`\sigma _0>70`$ km s<sup>-1</sup> span a much smaller range in age and \[Fe/H\]. We note that the non-treatment of non-solar abundance ratios in combination with correlated errors could be the sole reason for the trend found in Fornax. The direction and magnitude of correlated errors for a galaxy of solar metallicity and 8 Gyrs age are shown in panel (c), top right corner. Following on from the age–metallicity relation Jørgensen (1999) established for the Coma cluster an age–\[Mg/H\]–$`\sigma _0`$ relation. It would be very interesting to see whether such a correlation exists also in Fornax. However, the small number of galaxies combined with a rather small spread in age makes such an analysis very insecure and has therefore not been attempted.
The Mg-overabundance shows a weak positive correlation with central velocity dispersion and \[Fe/H\] in the sense that dynamically hotter and more metal rich galaxies are more overabundant (panel d, f). In the Fornax cluster significant overabundances are found for galaxies with $`\sigma _0100`$ km s<sup>-1</sup> or $`[\mathrm{Fe}/\mathrm{H}]0.0`$ (panel f, g). The best fitting linear relation between \[Mg/Fe\] and $`\mathrm{log}\sigma _0`$ is:
$$[Mg/Fe]=(0.49\pm 0.18)\mathrm{log}\sigma (0.80\pm 0.41).$$
(6)
This relation is qualitatively in agreement with the results from the Coma cluster (Jørgensen 1999, long-dashed line in Figure 12d). The scatter about the \[Mg/Fe\]–$`\sigma _0`$ relation in Fornax is consistent with the errors for \[Mg/Fe\], but there seems to be a rather large spread in the \[Mg/Fe\] ratio at a given metallicity (panel f). Although the latter is in good agreement with our findings in the Fe-Mg<sub>2</sub> diagram (Figure 5), we note that the errors are heavily correlated in the \[Mg/Fe\]–\[Fe/H\] diagram. There is no significant correlation of the \[Mg/Fe\] ratio with $`\mathrm{log}\mathrm{a}ge`$ where the young S0s show solar or slightly less than solar \[Mg/Fe\] ratios.
## 8 DISCUSSION
In this study, great care was taken to calibrate the line-strength measurements to a standard system in which we can compare the results with theoretical model predictions (Section 3). The accuracy of this calibration is vital when one wants to derive absolute age and metallicity estimates. Although we were able to demonstrate the high quality of our calibration some unresolved issues such as the systematic offset in the Mg<sub>2</sub> vs Mg b diagram, the rather large rms error in the original Lick/IDS stellar library and perhaps most important of all the largely unknown effects of non-solar abundance ratios prevent us from deriving accurate absolute age and metallicity estimates. However, for the discussion of relative differences in the stellar populations of early-type galaxies our data set and current models are of excellent use.
In this paper we have made use of two stellar population models provided by Worthey (1994) and Vazdekis et al. (1996). Both models make use of the Lick/IDS fitting functions but have otherwise somewhat different prescriptions to predict line-strength indices of integrated single-burst stellar populations (SSP). The predictions of the two models are consistent and our conclusions would not change if only one of them had been used for the analysis. To our knowledge, this would be also true if we had used any other model which makes use of the Lick/IDS fitting functions.
One of the most important results from this study is the homogeneity of the stellar populations in dynamically hot early-type galaxies in the Fornax cluster. Apart from Fornax A all early-type galaxies (Es & S0s) with $`\sigma _0>70`$ km s<sup>-1</sup> are of roughly the same age and their central metallicity scales with $`\mathrm{log}\sigma _0`$. The homogeneity is reflected in tight relations of observables such as Mg–$`\sigma _0`$ and Fe–$`\sigma _0`$ and a clear correlation of \[Fe/H\] with the central velocity dispersion. The existence of the latter is reassuring in terms of our current understanding of the colour-magnitude-relation (CMR) in clusters being mainly a result of increasing metallicity with increasing luminosity \[Kodama & Arimoto 1997, Terlevich et al. 1999\].
Previous authors \[Fisher, Franx & Illingworth 1996, Jørgensen 1997, Jørgensen 1999\] pointed out that the lack of a correlation of Fe-absorption strength with central velocity dispersion would give evidence for a second parameter or conspiracy of age, metallicity and \[Mg/Fe\] ratio which keeps the CMR tight. For example in the Coma cluster Jørgensen (1999) did not find a strong correlation of $`<`$Fe$`>`$ with central velocity dispersion and hence her \[Fe/H\]–$`\sigma _0`$ relation is also not significant. However, both the $`<`$Mg$`>`$$`\sigma _0`$ and \[Mg/H\]–$`\sigma _0`$ relation are clearly seen in Coma. In this context it is important to note that the slope of the Fe–$`\sigma _0`$ relation which one would expect from the change of metallicity in the CMR \[Kodama & Arimoto 1997\] is quite shallow and therefore only detectable with high S/N data. In contrast, the Mg–$`\sigma _0`$ relation is steeper, and therefore easier to detect. The reason for this is a combination of a larger dynamical range in the Mg indices compared to the average Fe-index and an increasing Mg overabundance with central velocity dispersion giving a steeper slope than what would be expected from the change in metallicity only.
An alternative explanation for the lack of a Fe–$`\sigma _0`$ relation in Coma could be based on galaxies such as IC2006, which do not follow the Fe–$`\sigma _0`$ relation very well. If this type of galaxy is more frequent in the Coma cluster than in Fornax it would be impossible to find a clear Fe–$`\sigma _0`$ relation. In summary we find that in the Fornax cluster there is no need for a second parameter such as age, metallicity or \[Mg/Fe\] to keep the CMR tight. Indeed, we favour an interpretation where small variations of age, metallicity and/or \[Mg/Fe\] at any given $`\sigma _0`$ are responsible for some real scatter in the scaling relations for the Fornax cluster. However, we emphasize that this may not be true for other (larger?) clusters.
In addition to the population of old, dynamically hot early-type galaxies, we find a sizeable fraction of young, dynamically colder ($`\sigma _070`$ km s<sup>-1</sup>) systems within our magnitude limited survey. Some of the young S0s (NGC1375, ESO359-G02 and ESO358-G25) fit in remarkably well with the predictions of galaxy harassment in clusters (Moore, Lake, Katz 1998; Lake, Katz, Moore 1998). In this scenario, medium sized disc galaxies (Sc-type) fall into a cluster environment and get “harassed” by high speed encounters with cluster galaxies. The end-products are small spheroidal galaxies where some gas of the disc is driven into the centre of the galaxy. This gas is likely to be turned into stars in a central stellar burst. We note, that most of these young galaxies are in the periphery of the Fornax cluster consistent with having been “accreted” onto the cluster from the field.
Two of the S0s which show young populations in the centre, also have extended discs (NGC1380A and IC1963). This seems to be in contradiction with the harassment picture. However we emphasize, that the existing harassment simulations do not include spirals with a substantial bulge component. Here the bulge is likely to stabilize the disc and the end-products may be able to keep substantial disc components (Ben Moore, private communication). The existence of a population of dynamically colder galaxies with young stellar populations in the nuclear regions is in agreement with a typical (nearby) cluster CMR where one finds a tail of blue galaxies towards the faint end (e.g., Terlevich 1998). Furthermore Terlevich et al. (1999) demonstrate for the Coma cluster, using line-strength analysis, that these blue galaxies contain young stellar populations rather than being metal poor.
It seems that in the Fornax cluster significant amounts of young stellar populations are predominantly found in low luminosity (lenticular) systems. However, for a sample of Coma cluster early-type galaxies Mehlert (1998) found that relatively bright S0s spread over the whole range in age (Es, excluding the cDs, are found to be old). This of course raises the question whether morphology is the driving parameter for young stellar populations (only S0s are younger) or whether luminosity is the important parameter (low luminosity E & S0 galaxies are on average younger). Taking the results from Coma and Fornax together we would like to argue that in clusters it is only the lenticular galaxies which show signs of recent star formation and that low luminosity lenticular systems are more likely to do so. The latter may be just caused by the recent accretion of these low luminosity systems onto the cluster.
So far we have addressed the age and metallicity distributions in the Fornax cluster with the help of line-strength indices. However there is more detailed information on the star formation (SF) processes to be gained if one investigates the \[Mg/Fe\] abundance ratios. When new stars are formed chemical enrichment is predominantly driven by the ejecta of SN Ia (main producer of Fe peak elements) and SN II (producing mainly alpha elements). However, SN Ia are delayed compared to SN II which explode on short time-scales of $`10^610^7`$ yr. Taking this into account there are mainly two mechanisms which determine the Mg/Fe ratio in galaxies: (i) the star formation time scale and (ii) the fraction of high mass stars, i.e., the initial mass function (IMF) (see e.g., Worthey et al. 1992). As re-confirmed in this study, the majority of cluster early-type galaxies show a trend of increasing \[Mg/Fe\] ratio with central velocity dispersion. Galaxies with young stellar populations and/or low luminosity galaxies show roughly solar abundance ratios. Given that the most luminous galaxies are also the metal richest, we emphasize that any realistic star-formation models have to be able to produce metal rich and Mg-overabundant stars at the same time.
In principle one can reproduce the observed trends of overabundances and metallicity with varying star-formation time-scales: large galaxies form within shorter time-scales than smaller galaxies \[Bressan, Chiosi, & Tantalo 1996\]. However, this leeds to extremely short star-formation time-scales for the most massive galaxies. A plausible way to solve this dilemma would be a varying IMF where massive galaxies have a top heavy IMF and low luminosity galaxies show a more Salpeter like IMF. For further discussions of the matter see also Peletier (1999) and Tantalo, Chiosi & Bressan (1998).
Recently, Thomas & Kauffmann (1999, see also Thomas 1999) presented preliminary results from their semi-analytic galaxy formation models for the distribution of \[Mg/Fe\] in galaxies as a function of luminosity. In this scenario luminous ellipticals are the last to form and hence Thomas & Kauffmann find a trend that the \[Mg/Fe\] ratio decreases with increasing luminosity, opposite to the observed trend. In general it seems very difficult with the current stellar population models to reproduce the observed Magnesium strength, and therefore \[Mg/Fe\] values ($`[\mathrm{Mg}/\mathrm{Fe}]0.4`$) of luminous ellipticals (Greggio 1997, but also see Sansom & Proctor 1998).
## 9 CONCLUSIONS
We have measured the central line strength indices in a magnitude limited sample of early-type galaxies brighter than $`\mathrm{M}_B=17`$ in the Fornax cluster and have applied the models of Worthey (1994), Worthey & Ottaviani (1997) and Vazdekis et al. (1996) to estimate their ages, metallicities and abundance ratios. We find that:
1. Elliptical Galaxies appear to be roughly coeval forming a sequence in metallicity varying roughly from $`0.25`$ to $`0.30`$ in \[Fe/H\]. This result is consistent with the conventional view of old, coeval elliptical galaxies where the metallicity scales with the luminosity of the galaxy. This is reflected in scaling relations such as Mg–$`\sigma _0`$. Remarkably, we could show that all other metal line-strength indices also clearly correlate with the central velocity dispersion. In fact all Fe-line–$`\sigma _0`$ relations are consistent with having the same slope.
2. Lenticular Galaxies have luminosity weighted metallicities spanning the whole range of SSP model predictions. Lower luminosity S0s show luminosity weighted ages younger than those of the ellipticals. However, the centres of the bright lenticular galaxies NGC1380 and NGC1381 resemble the properties of ellipticals suggesting that they experienced similar star formation histories. The peculiar S0 galaxy Fornax A (NGC1316), which is the brightest galaxy in the sample, has strong Balmer lines implying a very young luminosity weighted age, yet the metallicity is equal to the most metal rich Es. This is consistent with Fornax A having been involved in a recent gaseous merger. The S0s NGC1380 and NGC1381 follow the index–$`\sigma _0`$ relations of the ellipticals very well. However, the S0s with a young stellar component generally show a large scatter around the scaling relations.
3. Our conclusions are based on several age/metallicity diagnostic diagrams which give consistent results. Furthermore we demonstrate the advantage of using an emission robust age indicator such as H$`\gamma _\mathrm{A}`$ when analysing the stellar populations of extragalactic objects.
4. We have discovered that two of the fainter and very metal poor lenticular galaxies appear to have undergone major star-formation in the last 2 Gyrs (in one case very much more recently). We note that, like Fornax A, most of the young galaxies lie on the periphery of the cluster. This is consistent with the harassment picture where these galaxies are accreted from the field and undergo a morphological transformation with a central star burst.
5. The elliptical galaxies and the S0 NGC1380 exhibit overabundances up to 0.4 dex in Magnesium compared to Fe. There is a trend that the most massive and metal rich galaxies are the most overabundant, whereas the fainter Es approach solar ratios. This trend is inconsistent with the currently available semi-analytical predictions for hierarchical galaxy formation. S0s with young stellar populations are consistent with roughly solar abundance ratios, and may even be slightly underabundant. Remarkably also Fornax A, the brightest galaxy in our sample, shows close to solar abundance ratios which is not what one would expect of an early-type galaxy of its size.
6. Furthermore we note that abundance ratio trends, which are not included in the models, can lead to a change of relative age and metallicity estimates depending on which index combination is used in the analysis. As long as the non-solar abundance ratios are not properly incorporated into the models the estimation of absolute ages of integrated stellar populations remains insecure.
## ACKNOWLEDGEMENTS
HK acknowledges the use of STARLINK computing facilities at the University of Durham and also wishes to thank the Dr. Carl Duisberg Stiftung and University of Durham for generous financial support during the course of this work. Special thanks go to Roger Davies who provided excellent supervision throughout this project. Interesting and very helpful discussions of the following people are acknowledged: Eric Bell, Richard Bower, Taddy Kodama, Reynier Peletier, Jim Rose and Alexandre Vazdekis. Thanks also go to Scott Trager and Guy Worthey for providing the Lick/IDS measurements of the higher order Balmer lines. HK thanks the referee, Inger Jørgensen, for a careful and thorough reading of this paper which helped to improve the final presentation.
## Appendix A Kinematics
In Figure 13 a literature comparison of the central velocity dispersion measurements is presented. The most recent data for Fornax galaxies from Graham et al. (1998) as well as the literature compilation of Smith (1998, Private communication) is in good agreement with our data. The comparison with the literature compilation by McElroy (1995) shows a somewhat larger scatter. In order to establish an average error for the velocity dispersion estimates the mean scatter of the data compared to Smith (1998) and Graham et al. (1998) is evaluated. For the comparison all galaxies with $`\sigma _0<70`$ km s<sup>-1</sup> and also NGC1419 (marked in Figure 13) from the Smith compilation were excluded. The mean scatter is 0.027 in $`\mathrm{log}\sigma _0`$. Subtracting in quadrature a mean error of 0.015 for the literature data gives an error of 0.022 \[in $`\mathrm{log}\sigma _0`$\] for galaxies with $`\sigma _070`$ km s<sup>-1</sup>. For galaxies with $`\sigma _0<70`$ km s<sup>-1</sup> we adopt the rms scatter of the template stars as estimate of the error ($`\mathrm{\Delta }\mathrm{log}\sigma _0=0.074`$).
## Appendix B Lick/IDS calibration
In Figure 14 a comparison between the original Lick/IDS index measurements of stars and selected galaxies in common with our data is shown. The individual diagrams show the scatter (Lick/IDS - AAT) around the mean offset vs the average of Lick/IDS and AAT.
Figure 15 shows the velocity dispersion corrections for each index as derived from broadened stellar and selected galaxy spectra.
## Appendix C Comparison of derived ages, metallicities and \[Mg/Fe\] ratios
Here we present three Figures (16, 17 & 18) comparing our age, metallicity and \[Mg/Fe\] estimates derived with different index combinations. In general we find very consistent results when we use e.g., $`\mathrm{H}\beta `$ instead of $`\mathrm{H}\gamma _\mathrm{A}`$ as an age indicator. Most of the differences, such as the \[Mg/Fe\] ratios of some of the young lenticular galaxies, is caused by the data points lying at the edge of the model predictions or where they have been extrapolated. The extrapolation was necessary because the models reflect roughly solar abundance ratios and cannot account for e.g., the strongest Mg<sub>2</sub> absorption found in our sample.
## Appendix D Final calibrated central indices
### D.1 Table of central index measurements
The final, fully corrected central (2$`\stackrel{}{.}`$$`\times `$3$`\stackrel{}{.}`$85) index measurements and associated errors for the Fornax galaxies and NGC3379 are presented in Table D2. For each galaxy we give our Lick/IDS index measurement in the first row and the 1$`\sigma `$ error in the second row. The second column (first part of Table D2) lists central velocity dispersions in $`\mathrm{log}\sigma `$ units. The last column (part two of Table D2) lists H$`\beta _\mathrm{G}`$ (not Lick/IDS index); for further details see Section D.2.
### D.2 $`\mathrm{H}\beta `$ vs H$`\beta _\mathrm{G}`$ relation
Figure 19 shows a plot of $`\mathrm{H}\beta `$ equivalent width vs H$`\beta _\mathrm{G}`$ equivalent width. The bandpasses of the H$`\beta _\mathrm{G}`$ index are defined by Jørgensen (1997), based on an earlier definition of an $`\mathrm{H}\beta `$ emission index by González (1993), in such a way that the influence of the Fe feature right next to the $`\mathrm{H}\beta `$ absorption feature is minimized. For wavelength definitions see Table 7.
The indices $`\mathrm{H}\beta `$ & H$`\beta _\mathrm{G}`$ show an excellent correlation for the Fornax sample. The solid line in Figure 19 shows the fit to all the data, excluding ESO358-G25 (open triangle) because of its emission contamination, using a method which bisects the ordinary least squares fits made by minimising the X and the Y residuals:
$$\mathrm{H}\beta _\mathrm{G}=(0.862\pm 0.027)\mathrm{H}\beta +(0.568\pm 0.051)$$
(7)
The derived relation (errors from a Jack-Knife analysis) is in very good agreement with the relation found by Jørgensen (1997). |
warning/0001/hep-th0001081.html | ar5iv | text | # hep-th/0001081 BRST Formalism and Zero Locus Reduction
## 1. Introduction
The “BRST” quantization of general gauge theories in the Hamiltonian and Lagrangian formalisms includes the Batalin–Fradkin–Vilkovisky (BFV) and Batalin–Vilkovisky (BV) formalisms. From a geometric standpoint, these quantization formalisms deal with an even or odd QP manifold $`𝒩`$ , i.e., a symplectic or antisymplectic manifold equipped with a compatible odd vector field $`𝑸`$ such that $`𝑸^2=0`$. This condition is ensured by imposing the master equation on the Hamiltonian function of the vector field $`𝑸`$. In the standard physicists’ notation, the respective equations are
(1.1)
$$\{\mathrm{\Omega },\mathrm{\Omega }\}=0\text{and}(S,S)=0,$$
where $`\mathrm{\Omega }`$ (by a widespread abuse of terminology) is the “BRST generator” in the Hamiltonian quantization and $`S`$ is the master action in the Lagrangian quantization.
Under appropriate regularity conditions, the zero locus $`𝒵_𝑸𝒩`$ of $`𝑸=\{\mathrm{\Omega },\}`$ (of $`𝑸=(S,)`$) is an odd Poisson manifold (respectively, a Poisson manifold) , whose geometry captures crucial information about the theory on $`𝒩`$. In this paper, we mainly concentrate on even QP manifolds (which correspond to the BFV quantization and were implicit in ) because they have not been considered before; however, we formulate the general facts about the zero-locus reduction such that they apply to both even and odd QP manifolds. On an even QP manifold, $`𝒵_𝑸`$ carries an antibracket; we then show that the equivalence classes of observables (the cohomology of $`𝑸`$) are in a $`1:1`$ correspondence with characteristic (Casimir) functions of the antibracket on $`𝒵_𝑸`$, and gauge symmetries in the BFV theory on $`𝒩`$ are Hamiltonian vector fields on $`𝒵_𝑸`$.
Moreover, the zero locus $`𝒵_𝑸`$ of the BFV differential on the extended phase space is a proper counterpart of the constraint surface in the following sense. In geometric terms, a first-class constrained system can be specified by its phase space (a symplectic manifold $`𝒩_0`$) and the constraint surface $`\mathrm{\Sigma }`$. On the extended phase space $`𝒩`$ constructed in the BFV quantization, we can consider the dynamical system whose constraint surface, by definition, is $`𝒵_𝑸`$ (in local coordinates on $`𝒩`$, the constraints can be chosen as the components of $`𝑸`$). Then the constrained systems $`(𝒩_0,\mathrm{\Sigma })`$ and $`(𝒩,𝒵_𝑸)`$ are equivalent: the respective algebras of the equivalence classes of observables are naturally isomorphic as Poisson algebras.
Beyond the BRST context, algebras of functions on QP manifolds, which are differential Poisson algebras (associative supercommutative algebras endowed with a bracket operation and a differential that is a derivation of the bracket) can arise from complexes endowed with a super-commutative associative multiplication and a Gerstenhaber-like multiplication (“bracket”); the differential is then interpreted as the $`𝑸`$-structure, and the bracket becomes the P-structure (the Poisson or the BV bracket on the dual (super)manifold). The basic examples are the cohomology complexes of a Lie algebra $`𝔞`$ with coefficients in $`𝔞`$ or $`𝖲𝔞`$ (the exterior and symmetric tensor algebras); the general case involves $`L_{\mathrm{}}`$ algebras .
In this algebraic context, reduction to the zero locus can yield relations between different complexes. In certain cases, the zero-locus reduction can be applied repeatedly; the equation ensuring the existence of a nilpotent vector field on the reduced manifold at the second step of the reduction can be the classical Yang–Baxter equation (CYBE), in which case the reduction leads to the well-known Sklyanin and Berezin–Kirillov brackets.
In addition to the usual QP manifolds, one can consider bi-QP manifolds, which are the geometric counterparts of bicomplexes, and in physical terms, originate in the BRST–anti-BRST ($`Sp(2)`$-symmetric/triplectic) quantization . With two BRST operators represented by two commuting (odd and nilpotent) vector fields, bi-QP manifolds might be called QQP manifolds; interestingly enough, the corresponding zero-locus reduction (to the submanifold on which both vector fields vanish) results in a “PP” manifold, i.e., gives rise to a bi-Hamiltonian structure. A typical example is obtained by starting with a Lie algebra $`𝔞`$ and deriving the second differential from a coalgebra structure. Compatibility between two differentials then implies that $`(𝔞,𝔞^{},𝔞𝔞^{})`$ is a Manin triple . There also exists an alternative construction of a bi-QP manifold from a single Lie algebra structure, which results in non-Abelian triplectic antibrackets on the space of common zeroes of the differentials (and thus, the zero locus reduction leads to a nontrivial relation to the bicomplex used in the extended BRST symmetry).
This paper is organized as follows. In Sec. 2.2, we recall the main points of the zero locus reduction on (odd or even) QP manifolds. Symmetries of QP manifolds are reviewed in Sec. 2.3. In Sec. 3, we turn to a more detailed analysis of even QP manifolds corresponding to the BFV quantization. In Secs. 3.13.2, we recall several facts about the BFV formalism in the form that is suitable for what follows. The results given in 3.4 state the relation between objects in the bulk of the phase space and on the zero locus submanifold. We briefly discuss in Sec. 3.5 how these results can be restated for the BV formalism. In Sec. 4, we consider specific brackets resulting from the zero-locus reduction. In Sec. 5, we study bi-QP manifolds.
## 2. Geometry of QP manifolds and zero locus reduction
Geometric objects underlying the BRST quantization are the $`QP`$ manifolds.
###### 2.1 Definition ().
A QP manifold is a supermanifold $`𝒩`$ equipped with a bracket $`\mathbf{\{},\mathbf{\}}`$ such that
(2.1)
$$\begin{array}{cc}\hfill \mathbf{\{}F,G\mathbf{\}}=& (1)^{(𝗉(F)+\kappa )(𝗉(G)+\kappa )}\mathbf{\{}G,F\mathbf{\}},\hfill \\ \hfill \mathbf{\{}F,GH\mathbf{\}}=& \mathbf{\{}F,G\mathbf{\}}H+(1)^{(𝗉(F)+\kappa ))𝗉(G)}G\mathbf{\{}F,H\mathbf{\}},\hfill \\ \hfill \mathbf{\{}F,\mathbf{\{}G,H\mathbf{\}}\mathbf{\}}=& \mathbf{\{}\mathbf{\{}F,G\mathbf{\}},H\mathbf{\}}+\mathbf{\{}G,\mathbf{\{}F,H\mathbf{\}}\mathbf{\}}(1)^{(𝗉(F)+\kappa )(𝗉(G)+\kappa )},\hfill \end{array}$$
for $`F,G,H_𝒩`$ (smooth functions on $`𝒩`$), and with an odd nilpotent vector field $`𝐐`$, $`𝐐^2=0`$, such that
(2.2)
$$𝑸\mathbf{\{}F,G\mathbf{\}}\mathbf{\{}QF,G\mathbf{\}}(1)^{𝗉(F)+\kappa }\mathbf{\{}F,𝑸G\mathbf{\}}=0,F,G_𝒩$$
(where $`𝗉(())`$ is the Grassmann parity). QP manifolds with a *Poisson bracket ($`\kappa =0`$)* are called even, and those with an *antibracket ($`\kappa =1`$)*, odd.
Odd $`QP`$ manifolds arise in the BV quantization, and even ones in the BFV quantization. Odd QP manifolds were introduced in and were studied in . In most of our definitions, QP manifolds can be either even or odd; in Sec. 3, however, we concentrate on even QP manifolds, which have not been given enough attention previously.
### 2.2. The zero locus of $`𝑸`$
In what follows, $`𝒵_𝑸`$ denotes the zero locus of the odd vector field $`𝑸`$ on a $`QP`$ manifold $`𝒩`$. We assume $`𝒵_𝑸`$ to be a nonempty smooth submanifold and denote by $`_{𝒵_𝑸}_𝒩`$ the ideal of smooth functions vanishing on $`𝒵_𝐐`$.
The odd vector field $`𝑸`$ is called regular if each function $`f_{𝒵_𝑸}`$ can be represented as
(2.3)
$$f=\underset{\alpha }{}f_\alpha 𝑸\mathrm{\Gamma }^\alpha ,$$
with some $`f_\alpha ,\mathrm{\Gamma }^\alpha _𝒩`$ (i.e., if the components of $`𝑸`$ generate $`_{𝒵_𝑸}`$). We say that a submanifold $`𝒩`$ is coisotropic if
(2.4)
$$\mathbf{\{}_{},_{}\mathbf{\}}_{}.$$
###### 2.2.1 Lemma.
If $`𝐐`$ is regular, $`𝒵_𝐐`$ is a coisotropic submanifold of the QP manifold $`𝒩`$.
###### Proof.
Let $`f,g_𝒩`$ vanish on $`𝒵_𝑸`$. Using representation (2.3), the Leibnitz rule, Eq. (2.2), and nilpotency of $`𝑸`$, we see that $`\mathbf{\{}f_\alpha (𝑸\mathrm{\Gamma }^\alpha ),(𝑸\mathrm{\Gamma }^\beta )g_\beta \mathbf{\}}|_{𝒵_𝑸}=(f_\alpha \mathbf{\{}𝑸\mathrm{\Gamma }^\alpha ,𝑸\mathrm{\Gamma }^\beta \mathbf{\}}g_\beta )|_{𝒵_𝑸}=0`$. ∎
In what follows, we assume $`𝒵_𝑸`$ to be coisotropic even in those cases where $`𝑸`$ is not regular.
The algebra $`_{𝒵_𝑸}`$ of smooth functions on $`𝒵_𝑸`$ is the quotient $`_𝒩/_{𝒵_𝑸}`$. We then have
###### 2.2.2 Lemma.
There is a well-defined binary operation given by $`\mathbf{\{},\mathbf{\}}_𝐐:_{𝒵_𝐐}\times _{𝒵_𝐐}_{𝒵_𝐐}`$
(2.5)
$$\mathbf{\{}f,g\mathbf{\}}_𝑸=\mathbf{\{}F,𝑸G\mathbf{\}}|_{𝒵_𝑸},f,g_{𝒵_𝑸},F,G_𝒩,F|_{𝒵_𝑸}=f,G|_{𝒵_𝑸}=g,$$
where $`F`$ and $`G_𝒩`$ are viewed as representatives of functions on $`𝒵_𝐐`$. It makes $`𝒵_𝐐`$ into a Poisson manifold.
The proof is a straightforward generalization of a proof given in . It is obvious that the parity of the induced bracket on $`𝒵_𝑸`$ is opposite to the parity of the $`\mathbf{\{},\mathbf{\}}`$ bracket on $`𝒩`$. An important characteristic of the differential $`𝑸`$ is the homology of the linear operators $`𝑸_p:T_p𝒩T_p𝒩`$, $`p𝒵_𝑸`$, defined as follows. We consider the tangent space $`T_p𝒩`$ as the quotient of the vector fields $`\mathrm{Vect}_𝒩`$ modulo those that vanish at $`p`$. Then
(2.6)
$$𝑸_p(x)=([𝑸,X])|_p,X\mathrm{Vect}_𝒩,x=X_pT_p𝒩.$$
This operation is well-defined once $`𝑸`$ vanishes at $`p`$.
###### 2.2.3 Definition.
A QP manifold $`𝒩`$ is called *proper* if the homology of the linear operator $`𝐐_p:T_p𝒩T_p𝒩`$ is trivial at each point $`p𝒵_𝐐`$.
This definition is equivalent to the one given in (and ), but uses only invariant notions (in local coordinates $`\mathrm{\Gamma }^A`$, we would have $`(𝑸_px)^A=(1)^{𝗉(x)+1}x^B\frac{𝑸^A}{\mathrm{\Gamma }^B}`$.). We now have
###### 2.2.4 Proposition ().
Let $`𝒩`$ be a proper QP manifold with a nondegenerate bracket. Then $`𝒵_𝐐`$ is (anti)symplectic with respect to the induced bracket (2.5).
One can replace $`𝒵_𝑸`$ with a submanifold that still is coisotropic. As a straightforward generalization of 2.2.2, we have
###### 2.2.5 Theorem.
Let $`𝒩`$ be a QP manifold and $`𝒵_𝐐𝒩`$ a coisotropic submanifold of $`𝒩`$. Then $``$ is a Poisson manifold<sup>1</sup><sup>1</sup>1By Poisson manifolds, we mean those with either an even bracket or an antibracket. with the Poisson structure given by
(2.7)
$$\mathbf{\{}f,g\mathbf{\}}_𝑸=\mathbf{\{}F,𝑸G\mathbf{\}}|_{},f,g_{},F,G_𝒩,F|_{}=f,G|_{}=g.$$
###### Proof.
It is easy to see that (2.7) does not depend on the choice of representatives $`F,G_𝒩`$ of $`f,g_{}`$. The Jacobi identity and the Leibnitz rule follow in the same way as for the bracket in Eq. (2.5), see . ∎
### 2.3. Symmetries of QP manifolds
We now recall several basic facts about symmetries of QP structures on a manifold.
###### 2.3.1 Definition.
A vector field $`X`$ on a QP manifold $`𝒩`$ is called a symmetry of $`𝒩`$ if it commutes with $`𝐐`$ and is a Poisson vector field, i.e.,
(2.8)
$$X\mathbf{\{}F,G\mathbf{\}}\mathbf{\{}XF,G\mathbf{\}}(1)^{(𝗉(F)+\kappa )𝗉(X)}\mathbf{\{}F,XG\mathbf{\}}=0,F,G_𝒩.$$
Symmetries of the form $`X=\mathbf{\{}𝐐F,\mathbf{\}}`$ (with $`F_𝒩`$) are called *trivial*.
The Lie algebras of symmetries and trivial symmetries behave in a very regular manner under the restriction to $`𝒵_𝑸`$.
###### 2.3.2 Proposition.
Let $`X`$ be a symmetry of $`𝒩`$. Then $`X`$ restricts to $`𝒵_𝐐`$ and its restriction $`x`$ is a Poisson vector field on $`𝒵_𝐐`$ with respect to the bracket (2.7) on $`𝒵_𝐐`$, namely
(2.9)
$$x\mathbf{\{}F,G\mathbf{\}}_𝑸\mathbf{\{}xF,G\mathbf{\}}_𝑸(1)^{(𝗉(F)+\kappa +1)𝗉(X)}\mathbf{\{}F,xG\mathbf{\}}_𝑸=0,F,G_{𝒵_𝑸}.$$
If in addition $`X=\mathbf{\{}𝐐H,\mathbf{\}}`$ is a trivial symmetry, $`x`$ is a Hamiltonian vector field with respect to the $`\mathbf{\{},\mathbf{\}}_𝐐`$ bracket.
###### Proof.
Any symmetry $`X`$ restricts to $`𝒵_𝑸`$ because $`XF|_{𝒵_𝑸}=0`$ for any $`F`$ vanishing on $`𝒵_𝑸`$. Indeed, every such function can be represented as $`F=F_\alpha 𝑸\mathrm{\Gamma }^\alpha `$ with some functions $`F_\alpha `$ and $`\mathrm{\Gamma }^\alpha `$, provided $`𝑸`$ is regular. Because $`[X,𝑸]=0`$, we have $`XF|_{𝒵_𝑸}=((XF_\alpha )(𝑸\mathrm{\Gamma }^\alpha ))|_{𝒵_𝑸}+(1)^{𝗉(X)(𝗉(F_\alpha )+1)F_\alpha }(𝑸X\mathrm{\Gamma }^\alpha )|_{𝒵_𝑸}=0`$. Equation (2.9) immediately follows from the definition of the zero locus bracket and the definition of symmetries. If in addition $`X=\mathbf{\{}𝑸H,\mathbf{\}}`$ is a trivial symmetry, for any function $`f_{𝒵_𝑸}`$ we have
(2.10)
$$xf=X|_{𝒵_𝑸}f=\mathbf{\{}𝑸H,F\mathbf{\}}|_{𝒵_𝑸}=(1)^{p(H)+\kappa +1}\mathbf{\{}H|_{𝒵_𝑸},f\mathbf{\}}_𝑸,$$
where $`F_𝒩`$ is a lift of $`f`$ (i.e., $`f=F|_{𝒵_𝑸}`$) and $`\kappa `$ is the parity of the $`\mathbf{\{},\mathbf{\}}`$ bracket. Thus, $`x=X|_{𝒵_𝑸}`$ is a Hamiltonian vector field with respect to the bracket $`\mathbf{\{},\mathbf{\}}_𝑸`$. ∎
## 3. Observables, gauge symmetries, and zero locus reduction in BFV and BV quantizations
We now consider the embedding of a constrained system into the BFV extended theory with the BRST charge $`\mathrm{\Omega }`$ and study the “on-shell” gauge symmetries in the two descriptions of the same theory. In the Dirac (“non-extended”) formalism, the on-shell gauge symmetries are those nonvanishing on the constraint surface, and in the BFV extended formalism, these are symmetries nonvanishing on the zero locus $`𝒵_𝑸`$. We show that the former are mapped into the latter such that the equivalence classes of observables in the original theory are mapped into equivalence classes of observables in the BFV theory (the latter can be considered as gauge invariant functions on $`𝒵_𝑸`$). In this sense, the zero locus $`𝒵_𝐐`$ plays the role of a constraint surface in the BFV theory. We concentrate on the BFV case, where we assume the phase space to be finite-dimensional; reformulation of our results for the BV quantization, although straightforward at the formal level, requires some care because the BV configuration space of any realistic model is infinite-dimensional (see 3.5).
### 3.1. A reminder on constrained dynamics
We begin with recalling several basic facts about constrained dynamics in the form that will be suitable in what follows.
#### 3.1.1. Basics of the Dirac constrained dynamics
We consider a first-class constrained Hamiltonian system, defined on a phase space (symplectic manifold) $`𝒩_0`$ with the constraints $`T_\alpha `$ (functions on $`𝒩_0`$) such that
(3.1)
$$\{T_\alpha ,T_\beta \}=U_{\alpha \beta }^\gamma T_\gamma ,$$
where $`\{,\}`$ is the Poisson bracket on $`𝒩_0`$. For simplicity, we assume the first-class constraints $`T_\alpha `$ to be irreducible. Let $`\mathrm{\Sigma }`$ denote the constraint surface $`T_\alpha =0`$. A geometrically invariant way to specify a first-class constrained system is to fix the pair $`(𝒩_0,\mathrm{\Sigma })`$ (a symplectic manifold and a coisotropic submanifold). Different choices for $`T_\alpha `$ then give different generators of the ideal of functions vanishing on $`\mathrm{\Sigma }`$.
By definition, an observable is a function on $`𝒩_0`$ satisfying $`\mathbf{\{}A,T_\alpha \mathbf{\}}|_\mathrm{\Sigma }=0`$. Under the standard regularity conditions, each function vanishing on $`\mathrm{\Sigma }`$ is proportional to the constraints, and therefore,
(3.2)
$$\{A,T_\alpha \}=A_\alpha ^\beta T_\beta $$
for some functions $`A_\beta ^\alpha `$. Observables vanishing on $`\mathrm{\Sigma }`$ are called trivial. Two observables are called equivalent if they differ by a trivial observable. The space of equivalence classes of observables is a Poisson algebra, i.e., is closed under multiplication and under the Poisson bracket (these operations are well-defined on the equivalence classes via representatives). This algebra can be conveniently thought of as a subalgebra in the algebra of functions on $`\mathrm{\Sigma }`$.
Infinitesimal gauge transformations, or gauge symmetries, are the Hamiltonian vector fields $`X_0=\{\varphi _0,\}`$, where $`\varphi _0=\varphi _0^\alpha T_\alpha `$ is a trivial observable. Gauge symmetries form a Lie algebra with respect to the commutator. For any observable $`A`$ and a gauge symmetry $`X_0=\{\varphi _0,\}`$, we have
(3.3)
$$X_0A=\{\varphi _0,A\}=\{\varphi _0^\alpha T_\alpha ,A\}=\varphi _0^\alpha \{T_\alpha ,A\}+T_\alpha \{\varphi _0^\alpha ,A\},$$
which vanishes on $`\mathrm{\Sigma }`$ because $`A`$ is an observable. Therefore, gauge symmetries preserve equivalence classes of observables.
By the on-shell gauge symmetries, we mean the equivalence classes of gauge symmetries modulo those vanishing on the constraint surface $`\mathrm{\Sigma }`$. On-shell gauge symmetries can also be viewed as a subalgebra in the algebra of vector fields on $`\mathrm{\Sigma }`$. Equivalence classes of observables (viewed as functions on $`\mathrm{\Sigma }`$) are then represented by functions annihilated by on-shell gauge symmetries.
#### 3.1.2. Basics of the BFV/BRST approach
In the BFV quantization, the extended phase space $`𝒩`$ is an even QP manifold whose $`𝑸`$-structure is given by $`𝑸=\{\mathrm{\Omega },\}`$, where $`\mathrm{\Omega }`$ is a function on $`𝒩`$ (called the BRST charge) satisfying $`\{\mathrm{\Omega },\mathrm{\Omega }\}=0`$. In applications, the BFV extended phase space is usually equipped with an additional structure, the ghost charge $`G_𝒩`$. Functions with a definite ghost number are eigenfunctions of the ghost number operator
(3.4)
$$g=\{G,\},$$
corresponding to integer eigenvalues. The BRST charge is required to have the ghost number $`1`$,
(3.5)
$$\{G,\mathrm{\Omega }\}=\mathrm{\Omega }.$$
We now consider a QP manifold $`𝒩`$ that is not necessarily constructed via the BFV prescription; however, we refer to the objects on $`𝒩`$ as BFV ones because the applications in what follows will be to the case where $`𝒩`$ does result from the BFV construction. This also helps to distinguish between observables and symmetries on the QP manifold and those in the initial theory (Sec. 3.1.1), with ‘BFV’ used to refer to the former.
A BFV observable $`A`$ is a function on the QP manifold $`𝒩`$ satisfying
(3.6)
$$𝑸A=\{\mathrm{\Omega },A\}=0,\mathrm{gh}(A)=0.$$
The $`𝑸`$-exact BFV observables are called trivial. Two BFV observables $`A`$ and $`\stackrel{~}{A}`$ are equivalent if $`A\stackrel{~}{A}=𝑸B`$ for some function $`B`$; the equivalence classes of observables are then the cohomology of $`𝑸`$ in the ghost number zero. The algebra of BFV observables is a Poisson algebra (multiplication and the Poisson bracket can be defined via representatives).
A vector field $`X`$ is called a BFV gauge symmetry if $`X=\{𝑸H,\}`$ for some function $`H`$ with $`\mathrm{gh}(X)=\mathrm{gh}(𝑸H)=0`$ (these are trivial symmetries (see 2.3.1) of the corresponding QP manifold). In other words, BFV gauge symmetries are the Hamiltonian vector fields generated by trivial BFV observables. If $`A`$ is an observable and $`X=\{𝑸H,\}`$ a BFV gauge symmetry, we see that $`XA=\{𝑸H,A\}=𝑸\{H,A\}`$ is a trivial observable, i.e., BFV gauge symmetries preserve the equivalence classes of BFV observables.
### 3.2. From Dirac to the BFV formulation of a constrained system
Formal similarities between the Dirac and BFV formalisms are summarized in Table 1.
We now make contact between 3.1.1 and 3.1.2 by taking the extended phase space $`𝒩`$ and the BRST charge $`\mathrm{\Omega }`$ to be those arising in the BFV formalism from a given first-class constrained system $`(𝒩_0,\mathrm{\Sigma })`$. As before, the constraints $`T_\alpha _{𝒩_0}`$ are taken to be irreducible; to construct the BFV formalism, one then introduces ghosts $`c^\alpha `$, with $`\mathrm{gh}(c^\alpha )=1`$, $`𝗉(c^\alpha )=𝗉(T_\alpha )+1`$ and their conjugate momenta $`𝒫_\alpha `$,
(3.7)
$$\{c^\alpha ,𝒫_\beta \}=\delta _\beta ^\alpha ,$$
with $`\mathrm{gh}(𝒫_\alpha )=1`$, $`𝗉(𝒫_\alpha )=𝗉(T_\alpha )+1`$. The extended phase space $`𝒩`$ is the direct product of $`𝒩_0`$ with the superspace spanned by $`c^\alpha `$ and $`𝒫_\alpha `$.<sup>2</sup><sup>2</sup>2When the constraints are defined locally, the extended phase space is a vector bundle over the original phase space, as, for example, in . The Poisson bracket on $`𝒩`$ is the product Poisson bracket of that on $`𝒩_0`$ and (3.7).
One introduces the ghost charge (where we assume the constraints to be bosonic to avoid extra sign factors)
(3.8)
$$G=c^\alpha 𝒫_\alpha ,\{G,c^\alpha \}=c^\alpha ,\{G,𝒫_\alpha \}=𝒫_\alpha .$$
The BRST charge $`\mathrm{\Omega }`$ is an odd function defined by the condition that it has the ghost number 1 and satisfies
(3.9)
$$\{\mathrm{\Omega },\mathrm{\Omega }\}=0$$
with the boundary condition
(3.10)
$$\mathrm{\Omega }=c^\alpha T_\alpha +\mathrm{},$$
where $`\mathrm{}`$ means higher-order terms in the ghost momenta. It is well known that under standard assumptions, the BRST charge $`\mathrm{\Omega }`$ exists for any constrained system. Up to the first order in $`𝒫_\alpha `$, one has
(3.11)
$$\mathrm{\Omega }=c^\alpha T_\alpha \frac{1}{2}𝒫_\gamma U_{\alpha \beta }^\gamma c^\alpha c^\beta +\text{},$$
where the structure functions are those from (3.1).
As regards observables, the following statement is well known (see also ).
###### 3.2.1 Proposition.
The algebra of the equivalence classes of observables on $`𝒩_0`$ and the algebra of the equivalence classes of BFV observables (the cohomology of $`𝐐`$ in the ghost number zero) on the extended phase space $`𝒩`$ are isomorphic as Poisson algebras.
This means that if $`A_0_{𝒩_0}`$ is an observable of the constrained system on $`𝒩_0`$, there exists a BFV observable $`A_𝒩`$ with $`\mathrm{gh}(A)=0`$ such that
(3.12)
$$A|_{𝒩_0}=A_0.$$
Moreover, two BFV observables corresponding to the same observable $`A_0`$ differ by a trivial BFV observable. If in addition $`A_0`$ is a trivial observable, it follows that $`A=\{\mathrm{\Omega },B\}`$. The Poisson bracket on $`𝒩`$ induces a bracket on the cohomology of $`𝑸`$, and one has
(3.13)
$$\{A,B\}|_{𝒩_0}=\{A_0,B_0\}.$$
The isomorphism between the BRST cohomology in the ghost number zero and the algebra of equivalence classes of observables of the constrained system on $`𝒩_0`$ is given by the restriction of representatives to the initial constrained surface $`\mathrm{\Sigma }𝒩_0𝒩`$ (recall that equivalence classes of observables are gauge invariant functions on $`\mathrm{\Sigma }`$).
It also follows from 3.2.1 that because gauge symmetries of the initial system $`(𝒩_0,\mathrm{\Sigma })`$ are generated by trivial observables, each gauge symmetry can be lifted to a BFV gauge symmetry.
### 3.3. Zero locus $`𝒵_𝑸`$ in the BFV theory, the general case
We now consider an even QP-manifold $`𝒩`$ that is not necessarily constructed by the BFV procedure for a constrained system. We assume $`𝒩`$ to be symplectic, and the odd nilpotent vector field $`𝑸`$ to be regular in the sense of 2.3. The zero locus $`𝒵_𝑸`$ is thus a coisotropic submanifold of $`𝒩`$. Because each trivial BFV observable $`A=𝑸B`$ vanishes on $`𝒵_𝑸`$, each cohomology class uniquely determines a function on $`𝒵_𝑸`$. Thus, there is a mapping
(3.14)
$$H_𝑸^0_{𝒵_𝑸}$$
from the space of inequivalent observables to functions on $`𝒵_𝑸`$.
In what follows, we say that a statement holds locally if it is true in every sufficiently small neighbourhood. Mapping (3.14) is locally an embedding in view of the following proposition.
###### 3.3.1 Proposition.
Let $`𝐐=\{\mathrm{\Omega },\}`$ be regular in the sense of 2.2. Locally, each BFV observable vanishing on $`𝒵_𝐐`$ is a trivial BFV observable.
###### Proof.
Let $`A`$ be an observable vanishing on $`𝒵_𝑸`$, i.e., $`𝑸A=0`$$`A|_{𝒵_𝑸}=0`$. We must show that $`A=𝑸X`$ in a sufficiently small neighbourhood $`U`$ of any point $`p𝒵_𝑸`$. It is well known that locally there exists a coordinate system $`p_i,q^j,p_\alpha ,q^\beta ,c^\alpha ,𝒫_\beta `$ on $`𝒩`$ such that
(3.15)
$$\begin{array}{cc}& \mathrm{\Omega }=p_ic^i,\hfill \\ \hfill \{q^i,p_j\}=\delta _j^i,& \{q^\alpha ,p_\beta \}=\delta _\beta ^\alpha ,\{c^\alpha ,𝒫_\beta \}=\delta _\beta ^\alpha .\hfill \end{array}$$
Since the function $`A`$ vanishes on $`𝒵_𝑸`$, it can be represented as
(3.16)
$$A=A^\alpha p_\alpha +A_\alpha c^\alpha .$$
Now the odd vector field $`𝑸`$ becomes
(3.17)
$$𝑸=c^\alpha \frac{}{q^\alpha }+p_\alpha \frac{}{𝒫_\alpha },$$
and can be considered as the exterior differential under the identification $`c^\alpha =dq^\alpha `$, $`p_\alpha =d𝒫_\alpha `$, while $`A`$ becomes a 1-form. The assertion immediately follows from the super analogue of the Poincaré lemma in $`U`$. ∎
#### 3.3.2. Embedding $`H_𝑸^0`$ into functions on $`𝒵_𝑸`$.
As before, $`𝒵_𝑸`$ is the zero locus submanifold of $`𝑸=\{\mathrm{\Omega },\}`$ We recall from 2.3.2 that each BFV gauge symmetry $`X`$ can be restricted to $`𝒵_𝑸`$ and the restriction $`x=X|_{𝒵_𝑸}`$ is a $`\{,\}_𝑸`$-Hamiltonian vector field on $`𝒵_𝑸`$. The image of BFV gauge symmetries under the restriction to $`𝒵_𝑸`$ is called the algebra of the on-shell BFV symmetries. The functions on $`𝒵_𝑸`$ that are annihilated by the on-shell BFV symmetries are then the characteristic functions of the $`\{,\}_𝑸`$ antibracket on $`𝒵_𝑸`$.<sup>3</sup><sup>3</sup>3We recall that a function $`f`$ on the (odd) Poisson manifold $`𝒩`$ is said to be a characteristic (Casimir) function of an (odd) Poisson bracket $`\{,\}`$ if $`\{f,H\}=0`$ for any function $`H`$.
Because BFV observables are annihilated by BFV gauge symmetries, the restriction to $`𝒵_𝑸`$ maps BFV observables into characteristic functions of $`\{,\}_𝑸`$. For the equivalence classes of BFV observables (the cohomology of $`𝑸`$), this mapping is certainly an embedding locally. It is also an isomorphism in the important case of a BFV QP manifold considered in 3.4. Locally, we choose flat coordinates in some neighborhood $`U`$ of a point $`p𝒵_𝑸`$ and use explicit form (3.15) of the BRST charge $`\mathrm{\Omega }`$ and the Poisson bracket to arrive at
###### 3.3.3 Theorem.
Locally, the equivalence classes of BFV observables (the cohomology of $`𝐐`$) are in a $`1:1`$ correspondence with characteristic functions of the $`\{,\}_𝐐`$ antibracket on $`𝒵_𝐐`$.
We note that in one direction, this statement holds in general (i.e., not only locally) because for any observable $`A`$, we have
(3.18)
$$\{f,A|_{𝒵_𝑸}\}_𝑸=\{F,𝑸A\}|_{𝒵_𝑸}=0,$$
where $`F_𝒩`$ is the lift of a function $`f_{𝒵_𝑸}`$ and $`A|_{𝒵_𝑸}`$ is the image of $`A`$ under (3.14). Thus, $`A|_{𝒵_𝑸}`$ is a characteristic function of the antibracket $`\{,\}_𝑸`$ on $`𝒵_𝑸`$.
The “$`𝒵_𝑸`$-based” view on the BFV formalism developed here can be expressed as follows. Any even $`QP`$ manifold $`𝒩`$ gives rise to the constrained system $`(𝒩,𝒵_𝐐)`$, i.e., a constrained system whose phase space is $`𝒩`$ and the constrained surface is $`𝒵_𝐐`$. We recall from 3.1.1 that gauge transformations and the algebra of observables can be reconstructed if a first-class constrained system is specified in geometric terms, via its phase space (a symplectic manifold) and the constraint surface (a coisotropic submanifold). We now take this pair to be $`(𝒩,𝒵_𝑸)`$ (with $`𝒵_𝑸`$ being coisotropic in view of 2.2.1). In local coordinates, the constraints are the components of $`𝑸`$; in a neighborhood $`U𝒩`$, the following statement is obvious in the special coordinates in which $`\mathrm{\Omega }`$ and $`\{,\}_𝑸`$ are given by (3.15).
###### 3.3.4 Theorem.
On a QP manifold $`𝒩`$, the constrained system $`(𝒩,𝒵_𝐐)`$ is locally equivalent to the BFV theory on the extended phase space $`𝒩`$ with the BRST charge $`\mathrm{\Omega }`$ (i.e., the respective algebras of equivalence classes of observables are isomorphic as Poisson algebras).
In a more physical language, the equivalence can be reformulated by saying that the two constrained dynamics are equivalent.
The above considerations show that BFV observables are related to $`𝒵_𝑸`$ in the same way as observables in the initial theory (Sec. 3.1.1) are related to the constraint surface $`\mathrm{\Sigma }`$. This allows us to interpret $`𝒵_𝑸`$ as the extended constraint surface. In the general case, this correspondence takes place at the local level only.
### 3.4. Zero locus $`𝒵_𝑸`$ in the BFV formulation of a constrained system
We now concentrate on the important case where the QP manifold under consideration is a BFV extended phase space obtained by the BFV procedure from a given constrained system $`(𝒩_0,\mathrm{\Sigma })`$.
###### 3.4.1 Proposition.
The initial constraint surface $`\mathrm{\Sigma }𝒩_0`$ is a submanifold of the zero locus $`𝒵_𝐐𝒩`$ of the BRST differential $`𝐐=\{\mathrm{\Omega },\}`$.
###### Proof.
We restrict ourselves to an irreducible theory with constraints $`T_\alpha `$ (although the statement also is true for reducible constraints); the structure of the BRST charge is then given by (3.10). Considered as a submanifold in $`𝒩`$, the initial phase space $`𝒩_0`$ is determined by the equations $`c^\alpha =0`$ and $`𝒫_\alpha =0`$. It follows from (3.10) and from $`\mathrm{gh}(\mathrm{\Omega })=1`$ that the zero locus $`𝒵_𝑸`$ is determined by the equations
(3.19)
$$T_\alpha +\mathrm{}=0,\mathrm{}=0,$$
where $`\mathrm{}`$ denotes terms vanishing on $`𝒩_0`$. Then the intersection $`𝒵_𝑸𝒩_0`$ (considered as a submanifold in $`𝒩_0`$) is determined by the equations $`T_\alpha =0`$, and therefore, coincides with the initial constraint surface $`\mathrm{\Sigma }`$. Thus, $`\mathrm{\Sigma }`$ is a submanifold in $`𝒵_𝑸`$. ∎
The zero locus can be described somewhat more explicitly if we recall that in the BFV formalism, functions on the extended phase space are formal power series in the ghost variables $`c^\alpha `$ and $`𝒫_\alpha `$. This means that $`𝒵_𝑸`$ is actually determined by the equations
(3.20)
$$T_\alpha =0,c^\alpha =0.$$
This, in its turn, gives an explicit construction of the antibracket $`\{,\}_𝑸`$ on $`𝒵_𝑸`$. Let $`y^i`$ be local coordinates on $`\mathrm{\Sigma }`$. Then $`y^i`$ and $`𝒫_\alpha `$ can be considered as local coordinates on $`𝒵_𝑸`$. Evaluating (2.5), we now obtain
(3.21)
$$\{y^i,y^j\}_𝑸=0,\{𝒫_\alpha ,y^i\}_𝑸=R_\alpha ^i(y),\{𝒫_\alpha ,𝒫_\beta \}_𝑸=U_{\alpha \beta }^\gamma (y)𝒫_\gamma ,$$
where $`R_\alpha ^i(y)=\{T_\alpha ,y^i\}|_\mathrm{\Sigma }`$ and $`U_{\alpha \beta }^\gamma (y)=U_{\alpha \beta }^\gamma |_\mathrm{\Sigma }`$ with $`U_{\alpha \beta }^\gamma `$ from (3.11).
Using the explicit form (3.21) of the antibracket on $`𝒵_𝑸`$, it is easy to describe its characteristic functions in terms of the initial constraint surface $`\mathrm{\Sigma }`$. The following statement is obvious for irreducible constraints $`T_\alpha `$ and can be easily generalized to reducible constraints.
###### 3.4.2 Proposition.
Characteristic functions of the antibracket $`\{,\}_𝐐`$ are in a $`1:1`$ correspondence with gauge invariant functions on $`\mathrm{\Sigma }`$.
On a QP manifold constructed in accordance with the BFV prescription, the relation between the BRST cohomology and the geometry of the extended constrained surface $`𝒵_𝑸`$ can be made more precise than in the previous section. In particular, the respective counterparts of statements 3.3.1, 3.3.3, and 3.3.4 hold globally. We first see that (3.14) is an embedding.
###### 3.4.3 Proposition.
On a QP manifold $`𝒩`$ constructed in the BFV formalism, each BFV observable that vanishes on $`𝒵_𝐐𝒩`$ is a trivial BFV observable.
###### Proof.
Let $`A`$ be a BFV observable and $`A|_{𝒵_𝑸}=0`$. According to 3.4.1, $`\mathrm{\Sigma }𝒵_𝑸`$. Then $`A|_{𝒵_𝑸}=0`$ implies $`A|_\mathrm{\Sigma }=0`$ (a trivial observable). By 3.2.1, $`A`$ is a trivial BFV observable. ∎
We now consider the QP manifold constructed in the BFV formalism. Combining 3.4.3 with the argument given after 3.3.3 proves the next theorem in one direction; the other direction follows because each characteristic function on $`𝒵_𝑸`$ can be lifted to a BFV observable on $`𝒩`$, see 3.2.1 and 3.4.2.
###### 3.4.4 Theorem.
Equivalence classes of BFV observables (the cohomology of $`𝐐`$ with the ghost number zero) on the BFV QP manifold are in a $`1:1`$ correspondence with characteristic functions of the zero locus antibracket on $`𝒵_𝐐`$.
As in 3.3.3, we now consider the extended phase space of the BFV formulation as the phase space of a “new” constrained system determined by the constraint surface $`𝒵_𝑸`$. With $`𝒩`$ in its turn obtained from a constrained dynamical system $`(𝒩_0,\mathrm{\Sigma })`$ in accordance with the BFV formalism, we have a global version of 3.3.4.
###### 3.4.5 Theorem.
Let $`𝒩`$ be a QP manifold constructed in the BFV formalism. The constrained system determined by the pair $`(𝒩,𝒵_𝐐)`$ is equivalent to the BFV theory on $`𝒩`$ (i.e., the respective algebras of equivalence classes of observables are isomorphic as Poisson algebras).
Combining this with 3.2.1, we obtain a remarkable relation between the constrained systems specified by the respective pairs $`(𝒩_0,\mathrm{\Sigma })`$ and $`(𝒩,𝒵_𝑸)`$:
###### 3.4.6 Corollary.
The constrained systems $`(𝒩_0,\mathrm{\Sigma })`$ and $`(𝒩,𝒵_𝐐)`$ are equivalent (the respective algebras of inequivalent observables are isomorphic as Poisson algebras).
We also note a difference between the initial and the extended constraint surfaces $`\mathrm{\Sigma }`$ and $`𝒵_𝑸`$ in that $`\mathrm{\Sigma }`$ carries an action of the gauge generators $`\{T_i,\}`$, while $`𝒵_𝑸`$ is equipped with the zero locus antibracket. This is not unnatural, because the on-shell gauge symmetries are Hamiltonian vector fields with respect to the zero-locus antibracket, while inequivalent observables are (identified with) the characteristic functions of the zero-locus antibracket.<sup>4</sup><sup>4</sup>4This applies at the classical level. The notion of the initial and the extended constraint surfaces is essentially classical and has no obvious counterparts at the quantum level. At the quantum level, restrictions to the constraint surface should be understood as restriction to some quotient of the full Hilbert space of the quantum system. We do not discuss this very interesting subject here, and refer instead to , where a related problem was considered. We thank I.A. Batalin for an illuminating discussion of this point.
Finally, we note that there is a slightly different point of view on the interpretation of BFV observables in terms of the geometry of $`𝒵_𝑸`$. Namely, to each (odd) Poisson structure, one can associate the coboundary operator (differential) acting on antisymmetric tensor fields, with the action being the adjoint action of the Poisson bivector with respect to the Schouten-Nijenhuis bracket. Inequivalent observables are then the zero-degree cohomology of this differential on $`𝒵_𝑸`$ (tensors of zero degree are functions).
### 3.5. Observables, gauge symmetries, and zero locus reduction in the BV quantization
The above can be reformulated for odd QP manifolds/BV quantization. In the BV formulation, the zero locus of $`𝑸=(S,)`$, where $`S`$ is the master action, is the stationary surface of $`S`$ (provided the BV antibracket $`(,)`$ is nondegenerate). The BV observables are the cohomology of $`𝑸`$ in the ghost number zero. The BV gauge symmetries are the vector fields of the form
(3.22)
$$X=(𝑸B,),$$
and, thus, are Hamiltonian vector fields generated by trivial observables. Whenever the master action $`S`$ is constructed via the BV prescription starting from a given initial action $`S_0`$, the zero locus of $`𝑸=(S,)`$ is a certain extension of the stationary surface of the initial action $`S_0`$.
At the formal level, all the statements considered in the BFV scheme have their counterparts in the BV formalism. We do not restate here the contents of 3.13.4 for the odd case and refer instead to .<sup>5</sup><sup>5</sup>5To avoid misunderstanding, we note that we have changed our point of view on how the BFV/BV gauge symmetries should be defined: the BV gauge symmetries were called the “trivial gauge symmetries of the master action” in ; translating the results proved in into the present conventions, therefore, requires some care with the “obsolete” definitions. We only point out one important difference. Unlike the Hamiltonian picture, the Lagrangian one can be considered in the scope of a finite dimensional analogue only formally. The finite dimensional configuration space (the space of field histories) does not correspond to any physically relevant system. Thus all the BV counterparts of the statements of the previous section should be considered with some care. In particular, the BV quantization prescription requires the master action $`S`$ to be a proper solution to the master equation. The condition imposed on the master action to be proper has no counterpart in the Hamiltonian picture. It implies that the corresponding configuration space is a proper QP manifold (which in general is not the case for the BFV phase space). In the finite dimensional case, this in turn implies that all the observables (the cohomology of $`𝑸`$) are trivial (except those of a topological nature). The $`𝑸`$ cohomology becomes nontrivial only when evaluated on space-time local functionals .
## 4. Towers of brackets
In this section, we study the possibility of a “second” zero-locus reduction, i.e., the reduction on a QP manifold which itself is the result of a zero-locus reduction. This leads to several well-known structures, including the classical Yang–Baxter equation.
### 4.1. A “second” zero-locus reduction
On a QP manifold $`𝒩`$ (which can be either even or odd), a coisotropic submanifold $`𝒵_𝑸`$ (for example, a Lagrangian submanifold in $`𝒩`$) is a P-manifold, i.e., is equipped with an (even or odd) Poisson structure (see 2.2.5). One can try to equip $``$ with a compatible $`𝑸`$ structure, thereby making it into a QP manifold. On a general QP manifold $`𝒩`$, there is no canonical structure inducing a $`𝑸`$ operator on $``$. Instead, we can look for a $`𝑸`$ operator on $``$ in the form $`𝑸_{}=\mathbf{\{}H,\mathbf{\}}_𝑸`$, where $`\mathbf{\{},\mathbf{\}}_𝑸`$ is the bracket given by (2.7) and $`H`$ is a solution of the equation
(4.1)
$$\mathbf{\{}H,H\mathbf{\}}_𝑸=0,H_{},𝗉(H)=𝗉(\mathbf{\{},\mathbf{\}}_𝑸)+1.$$
Whenever such an $`H`$ is found, $``$ becomes a QP manifold. With this $`𝑸`$-structure, we can repeat the procedure, thereby producing a sequence of QP manifolds.
This construction can be restated in terms of differential Poisson algebras (the algebras of functions on QP manifolds). Even “more algebraically,” we consider the case where a differential Poisson algebra arises from a complex endowed with a super-commutative associative multiplication and a Gerstenhaber-like multiplication (see the Appendix). To these differential Poisson algebras, we can then apply one or more zero-locus reduction steps, resulting in relations between different complexes.
### 4.2. Examples of the zero locus reduction on an even QP manifold
Let $``$ be a cotangent bundle $`=T^{}𝒳`$. We then write $`(q^a,p_a)`$ for local coordinates on $``$ (which we take to be bosonic to avoid extra sign factors); the Poisson bracket then is $`\{q^a,p_b\}=\delta _b^a`$. We assume a Hamiltonian action of a Lie algebra $`𝔞`$ on $``$. For simplicity, we consider the Hamiltonian action that is the lift of an action on $`𝒳`$ via the vector fields $`X_i=X_i^a\frac{}{q^a}`$, with $`[X_i,X_j]=C_{ij}^kX_k`$. The generators of the Hamiltonian action on $``$ are then given by $`T_i=p_aX_i^a(q)`$. Applying the BFV scheme to the constraints $`T_i`$ gives the BRST generator
(4.2)
$$\mathrm{\Omega }=p_aX_i^a(q)\theta ^i\frac{1}{2}C_{ij}^k\xi _k\theta ^i\theta ^j.$$
We now take the submanifold $`𝒵_𝑸`$ (which is Lagrangian in $``$) determined by $`\theta ^i=0`$ and $`p_a=0`$ and view $`q^a`$ and $`\xi _i`$ as local coordinates on $``$. The antibracket $`(,)\mathbf{\{},\mathbf{\}}_𝑸`$ from 2.2.5 is then given by
(4.3)
$$(\xi _i,\xi _j)=C_{ij}^k\xi _k,(q^a,\xi _i)=X_i^a.$$
Using this antibracket structure on $``$, we consider the equation
(4.4)
$$(H,H)=0$$
for an even function $`H_{}`$. Given a solution $`H`$, we can construct the odd nilpotent vector field $`𝑸=(H,)`$ that makes $``$ into a QP manifold.
We consider solutions to (4.4) of the form
(4.5)
$$H_{\mathrm{YB}}=\frac{1}{2}r^{ij}\xi _i\xi _j,$$
where $`r`$ is a skew-symmetric matrix with entries from $`_𝒳`$. Explicitly, Eq. (4.4) is the following generalization of the CYBE:
(4.6)
$$r^{\mathrm{}[i}C_\mathrm{}m^kr^{j]m}+X_{\mathrm{}}^ar^{\mathrm{}[i}\frac{}{q^a}r^{jk]}=0.$$
We now proceed with the next stage of the zero locus reduction. The zero locus of the “Yang–Baxter differential” $`𝑸_{\mathrm{YB}}=(H_{\mathrm{YB}},)`$ is determined by $`r^{ij}\xi _j=0`$. We choose a smaller submanifold $`𝒳𝒵_{𝑸_{\mathrm{YB}}}`$ determined by $`\xi _i=0`$. Whenever (4.4) is satisfied, $`\{,\}=(,𝑸_{\mathrm{YB}})`$ is a Poisson bracket on $`𝒳`$. Explicitly, the Poisson brackets are given by
(4.7)
$$\{q^a,q^b\}=X_i^ar^{ij}X_j^b.$$
#### 4.2.1. The classical Yang–Baxter equation
Antibracket (4.3) considered on $`q`$-independent functions coincides with the Schouten bracket on $`𝔞`$ viewed as the Grassmann algebra generated by $`\xi _i`$. In the case where $`r^{ij}`$ is a constant matrix, (4.6) becomes the CYBE
(4.8)
$$r^{j[i}C_{jl}^kr^{m]l}=0.$$
For each $`r^{ij}`$ satisfying (4.8), the corresponding differential $`𝑸_{\mathrm{YB}}`$ (considered on $`𝔞`$) is nothing but the cohomology differential of the Lie algebra complex with trivial coefficients (see Appendix A), for the Lie algebra defined on $`𝔞^{}`$ by the structure constants $`F_k^{ij}=r^{il}C_{lk}^jr^{jl}C_{lk}^i`$.
#### 4.2.2. The Sklyanin bracket
With $`𝒳`$ taken to be the Lie group corresponding to the Lie algebra $`𝔞`$, we have two natural ways to define the action of $`𝔞`$ on $`𝒳`$, by the left- and right-invariant vector fields $`L_i`$ and $`R_i`$. Proceeding along the steps described in the previous paragraphs with $`X_i^a`$ taken to be $`L_i^a`$ or $`R_i^a`$, we arrive at two Poisson brackets on $`𝒳`$,
(4.9)
$$\{q^a,q^b\}_{\text{right}}=L_i^ar^{ij}L_j^b\text{and}\{q^a,q^b\}_{\text{left}}=R_i^ar^{ij}R_j^b,$$
which are compatible in view of $`[R_i,L_j]=0`$. The Poisson bracket
(4.10)
$$\{q^a,q^b\}_{\text{Sklyanin}}=\{q^a,q^b\}_{\text{right}}\{q^a,q^b\}_{\text{left}}$$
makes the Lie group $`𝒳`$ into a Poisson–Lie group.
### 4.3. Zero locus reduction on an odd QP manifold
To reformulate the above for an odd QP manifold, we construct the BV scheme starting with a manifold $`𝒳`$ with an $`𝔞`$ action. The $`\xi _i`$ variables are then even, and because of the symmetry properties, the “tower of reductions” is shorter than for odd $`\xi _i`$. We then introduce antifields $`q_a^{}`$, ghosts $`\theta ^i`$, and their antifields $`\xi _i`$, with $`(\theta ^i,\xi _j)=\delta _j^i`$ (where restored the traditional notation for the antibracket). The differential
(4.11)
$$𝑸=(S,),S=q_a^{}X_i^a\theta ^i\frac{1}{2}\xi _kC_{ij}^k\theta ^i\theta ^j$$
corresponds to the quantization of a theory with the vanishing classical action.
We choose a Lagrangian subspace $`𝒵_𝑸`$ determined by $`\theta ^i=0`$ and $`q_a^{}=0`$. In accordance with Sec. 2, the zero locus reduction induces a Poisson bracket $`\{,\}_𝑸`$ on $``$ with the nonvanishing components
(4.12)
$$\{\xi _i,\xi _j\}_𝑸=C_{ij}^k\xi _k,\{q^a,\xi _i\}_𝑸=X_i^a.$$
Unless $`𝒳`$ is a supermanifold, $``$ is a purely even manifold, and therefore, the new generating equation with respect to the $`\{,\}_𝑸`$-bracket has only the trivial solution. The tower of brackets is thus terminated.
We now recall that the even variables $`\xi _i`$ generate the algebra of functions on $`𝔞^{}`$. Restricting ourselves to functions that are independent of the coordinates on $`𝒳`$, we see that (4.12) becomes the Berezin–Kirillov bracket on $`𝔞^{}`$,
(4.13)
$$\{f,g\}=f\frac{\stackrel{}{}}{\xi _i}\xi _kC_{ij}^k\frac{}{\xi _j}g.$$
#### 4.3.1. Linear and nonlinear brackets
The bracket in (4.13) is “linear” in the sense of its explicit dependence on $`\xi _i`$. For a Lie algebra $`𝔞`$, one can construct “nonlinear” brackets $`\frac{\stackrel{}{}}{\xi _i}\mathrm{\Omega }_{ij}\frac{}{\xi _j}`$ on $`𝔞^{}`$, where the expansion of $`\mathrm{\Omega }_{ij}`$ in $`\xi _i`$ starts with $`\xi _kC_{ij}^k`$. For a given bracket of this form, a natural problem is whether it can be transformed into the Berezin–Kirillov bracket by a change of coordinates. With the help of the zero-locus reduction, this is solved as follows. The Poisson bracket is represented as the zero-locus reduction of the canonical antibracket on a QP manifold with $`𝑸`$ determined by the Hamiltonian $`H=\mathrm{\Omega }_{ij}(\xi )\theta ^i\theta ^j`$. The Jacobi identity for the Poisson bracket is rewritten as the master equation for $`H`$, and moreover, the terms containing higher powers of $`\xi _i`$ are closed with respect to the differential $`𝑸_0=\{H_0,\}`$, where $`H_0=\xi _kC_{ij}^k\theta ^i\theta ^j`$ is the “linear” part of the Hamiltonian. We thus have proved the fact known from other considerations (and in a more powerful analytic version)
###### 4.3.2 Corollary.
Let $`\mathrm{\Omega }_{ij}(\xi )=\xi _kC_{ij}^k+\xi _k\xi _lC_{ij}^{kl}+\mathrm{}`$ be the matrix of a Poisson bracket on $`𝔞^{}`$, where $`C_{ij}^k`$ are the structure constants of a Lie algebra $`𝔞`$. Then $`\mathrm{\Omega }_{ij}(\xi )`$ can be reduced to the form $`\xi _kC_{ij}^k`$ by a change of variables $`\xi _if_i(\xi )`$ if the second cohomology group $`H^2(𝔞,𝖲𝔞)`$ of $`𝔞`$ with coefficients in $`𝖲𝔞`$ is trivial.
Similar considerations in the BFV case lead to similar statements for the nonlinear antibracket.
## 5. Bi-QP manifolds
Up to now, we have studied QP manifolds whose differential corresponds to a single solution of the corresponding “master” equation. We now consider bi-QP manifolds.
### 5.1. A BFV-like formulation of the bialgebra complex
In the previous section, we associated an even QP manifold with a vector space $`𝔞`$ and a smooth manifold $`=T^{}𝒳`$. Namely, a Lie algebra structure on $`𝔞`$ and the vector fields $`X_i`$ (giving an $`𝔞`$-module structure on $`_𝒳`$) can be read off from a solution of the generating equation
(5.1)
$$\{\mathrm{\Omega },\mathrm{\Omega }\}=0$$
with the ansatz (4.2). The algebra of functions on the thus constructed QP manifold is $`𝔄=𝖧𝗈𝗆(𝔞,𝔞)_{T^{}𝒳}`$; we interpret $`𝖧𝗈𝗆(𝔞,𝔞)`$ as the algebra generated by the odd variables $`\theta ^i`$ and $`\xi _j`$. The basic Poisson bracket relations are
(5.2)
$$\{\theta ^i,\xi _j\}=\delta _j^i,\{q^a,p_b\}=\delta _b^a,$$
where $`q,p`$ are the standard local coordinates on the cotangent bundle $`=T^{}𝒳`$. We have the solution
(5.3)
$$C=p_aX_i^a(q)\theta ^i\frac{1}{2}\xi _kC_{ij}^k\theta ^i\theta ^j.$$
At the same time, every solution of (5.1) of the form
(5.4)
$$F=p_aX^{ia}\xi _i\frac{1}{2}\theta ^kF_k^{ij}\xi _i\xi _j$$
determines a coalgebra structure on the vector space $`𝔞`$, or equivalently, a Lie algebra structure on $`𝔞^{}`$, and makes $`_𝒳`$ into an $`𝔞^{}`$-module, with the vector fields $`X^i=R^{ia}\frac{}{q^a}\mathrm{Vect}_𝒳`$ representing the action of the basis elements of $`𝔞^{}`$. Then $`𝔄`$ is equipped with Poisson bracket (5.2) and the differentials
(5.5)
$$\begin{array}{cc}\hfill d_C=& \{C,\}\hfill \\ \hfill =& \frac{1}{2}C_{ij}^k\theta ^i\theta ^j\frac{}{\theta ^k}\xi _kC_{ij}^k\theta ^i\frac{}{\xi _j}p_aX_i^a\frac{}{\xi _i}+\theta ^iX_i^a\frac{}{q^a}\theta ^ip_aX_{i,b}^a\frac{}{p_b},\hfill \\ \hfill d_F=& \{F,\}\hfill \\ \hfill =& \frac{1}{2}F_k^{ij}\xi _i\xi _j\frac{}{\xi _k}\theta ^kF_k^{ij}\xi _i\frac{}{\theta ^j}p_aX^{ia}\frac{}{\theta ^i}+\xi _iX^{ia}\frac{}{q^a}\xi _ip_aX_{,b}^{ia}\frac{}{p_b}.\hfill \end{array}$$
We next impose the condition that the differentials be compatible, i.e.,
(5.6)
$$[d_C,d_F]=0\{C,F\}=0.$$
###### 5.1.1 Proposition.
Condition (5.6) implies that $`(𝔞,𝔞^{},𝔞𝔞^{})`$ is a Manin triple , with the Lie bracket on $`𝔞𝔞^{}`$ given by
(5.7)
$$[e_i,e_j]=C_{ij}^ke_k,[e^i,e^j]=F_k^{ij}e^k,[e_i,e^j]=C_{ik}^je^k+F_i^{jk}e_k.$$
where $`e_i`$ and $`e^i`$ are dual bases in $`𝔞`$ and $`𝔞^{}`$ respectively. Equivalently, $`𝔞`$ is a Lie bialgebra. Moreover, $`_𝒳`$ is a module over the Lie algebra $`𝔞𝔞^{}`$.
The proof is straightforward.
That $`_𝒳`$ is a module over $`𝔞𝔞^{}`$ means that under the mapping $`e_iX_i`$, $`e^iX^i`$, the following commutation relations between vector fields are satisfied:
(5.8)
$$[X_i,X_j]=C_{ij}^kX_k,[X^i,X^j]=F_k^{ij}X^k,[X_i,X^j]=C_{ik}^jX^k+F_i^{jk}X_k.$$
It also follows from (5.6) that
(5.9)
$$X^{ia}X_i^b+X^{ib}X_i^a=0.$$
#### 5.1.2. Zero locus reduction on a bi-QP manifold
We next consider the submanifolds of the zero loci, $`_C𝒵_C`$ and $`_F𝒵_F`$ defined by $`(\theta ^i=0,p_a=0)`$ and $`(\xi _i=0,p_a=0)`$, respectively. Since $`_C`$ and $`_F`$ are coisotropic, we can apply Theorem 2.2.5. We thus have the respective antibrackets
(5.10)
$$\begin{array}{cc}\hfill \{\xi _i,\xi _j\}_C=& C_{ij}^k\xi _k,\{\xi _i,q^a\}_C=X_i^a,\hfill \\ \hfill \{\theta ^i,\theta ^j\}_F=& F_k^{ij}\theta ^k,\{\theta ^i,q^a\}_F=X^{ia},\hfill \end{array}$$
on $`_C`$ and $`_F`$.
###### 5.1.3 Proposition.
The differential $`d_C`$ induces a well-defined operator (vector field) $`\overline{d}_C=d_C|__F:__F__F`$ and the differential $`d_F`$ induces an operator $`\overline{d}_F=d_F|__C:__C__C`$. Thus, $`__F`$ ($`__C`$) is an odd differential Poisson algebra and $`_F`$ (respectively, $`_C`$) is an odd QP manifold.
Thus, the manifolds $`_C`$ and $`_F`$ are equipped with $`𝑸`$ structures. We now proceed to the next step of the zero locus reduction.
Recall that the submanifold $`𝒳=_C_F`$ is determined by the equations $`p_a=\xi _i=\theta ^j=0`$. It is easy to see that $`𝒳`$ is a coisotropic submanifold of $`_C`$ and also a coisotropic submanifold of $`_F`$. On $`𝒳`$, we then have the Poisson bracket
(5.11)
$$\{,\}_𝒳=\{,\overline{d}_F\}_C=\{,\overline{d}_C\}_F$$
or in the coordinate form,
(5.12)
$$\{q^a,q^b\}_𝒳=X_i^aX^{ib}.$$
It follows from (5.9) that bracket (5.12) is skew-symmetric; the Jacobi identity follows from the compatibility of $`d_C`$ and $`d_F`$.
#### 5.1.4. Coboundary bialgebras
Up to this point, the situation was symmetric with respect to $`\theta ^i\xi _i`$, but now we try to solve Eq. (5.6) for $`F`$. Namely, suppose that $`F`$ is a coboundary
(5.13)
$$F=d_Cr=\{C,r\},$$
where $`r=r^{ij}\xi _i\xi _j`$ and $`r^{ij}`$ is taken to be a constant matrix. Then the condition $`d_F^2=0`$ yields
(5.14)
$$\{C,\{r,\{C,r\}\}\}=d_C\{r,d_Cr\}=0.$$
This is the generalized CYBE. An even stronger condition
(5.15)
$$\{r,d_Cr\}=\{r,r\}_C=0$$
leads to the CYBE (see (4.4)).
### 5.2. Two differentials from a Lie algebra action
We now look at the bicomplex setting from a somewhat different point of view. Rather than associating a second differential with a coalgebra structure, we construct a pair of differentials for a single Lie algebra. This subject attracts one’s attention because of its possibly deep relation to the extended BRST symmetry . We now show that the bicomplex generalization of the zero locus reduction method induces the non-Abelian triplectic antibrackets on the space of common zeroes of the differentials.<sup>6</sup><sup>6</sup>6The non-Abelian triplectic antibrackets were introduced in , see also , as the structure underlying a possible generalization of the well known Lagrangian version of the extended BRST quantization.
#### 5.2.1. Left and right $`𝔞`$ actions
We consider the left and the right actions of $`𝔞`$ on $`𝒳`$. To illustrate the idea, we restrict ourselves to the case where $`𝒳=𝒢`$ is the Lie group corresponding to the Lie algebra $`𝔞`$. Let the basis elements $`e_i`$ of $`𝔞`$ act on $`𝒢`$ via the left invariant vector fields $`L_i`$ (which correspond to the right action) and via the right invariant vector fields $`R_i`$ (which correspond to the left action). Obviously, $`[L_i,R_j]=0`$. Let $`q^a`$ and $`p_a`$ be the standard coordinates on $`T^{}𝒢`$. Unlike in the case considered above, we introduce the doubled set of variables $`\xi _i^1`$, $`\xi _j^2`$, $`\theta _1^k`$, and $`\theta _2^l`$, $`i,j,k,l=1,\mathrm{},dim𝔞`$, with the basic Poisson brackets
(5.16)
$$\{q^a,p_b\}=\delta _b^a,\{\theta _1^i,\xi _j^1\}=\delta _j^i,\{\theta _2^i,\xi _j^2\}=\delta _j^i.$$
The functions
(5.17)
$$\begin{array}{cc}\hfill \mathrm{\Omega }^1& =p_aR_i^a\theta _1^i\frac{1}{2}\xi _k^1C_{ij}^k\theta _1^i\theta _1^j,\hfill \\ \hfill \mathrm{\Omega }^2& =p_aL_i^a\theta _2^i\frac{1}{2}\xi _k^2C_{ij}^k\theta _2^i\theta _2^j\hfill \end{array}$$
satisfy $`\{\mathrm{\Omega }^\alpha ,\mathrm{\Omega }^\beta \}=0`$ for $`\alpha ,\beta =1,2`$, as follows immediately from the commutativity of the left- and right-invariant vector fields. These generating functions give rise to the anticommuting differentials $`𝑸^a=\{\mathrm{\Omega }^a,\}`$, thereby providing $`_{\mathrm{ext}}`$ with a bicomplex structure.
#### 5.2.2. Zero locus reduction in $`_{\mathrm{ext}}`$ and nonabelian triplectic antibrackets
We now apply the zero locus reduction along the lines of Sec. 2. We identify the zero locus $`𝒵_{𝑸^1}`$ (respectively, $`𝒵_{𝑸^2}`$) of the differential $`𝑸^1`$ (of $`𝑸^2`$) determined by the equations $`\theta _1^i=0`$ and $`p_a=0`$ (respectively, $`\theta _2^i=0`$ and $`p_a=0`$). The intersection $`=𝒵_{𝑸^1}𝒵_{𝑸^2}`$ is then endowed with a pair of compatible antibrackets. Identifying $`_{}`$ (functions on the intersection) with functions of $`q^a`$, $`\xi _i^1`$, and $`\xi _j^2`$, we have
(5.18)
$$\begin{array}{c}\hfill \{\xi _i^1,q^a\}_{𝑸^1}=R_i^a,\{\xi _i^1,\xi _j^1\}_{𝑸^1}=C_{ij}^k\xi _k^1,\\ \hfill \{\xi _i^2,q^a\}_{𝑸^2}=L_i^a,\{\xi _i^2,\xi _j^2\}_{𝑸^2}=C_{ij}^k\xi _k^2,\end{array}$$
with all the other brackets vanishing. These are precisely the non-Abelian triplectic antibrackets from .
## 6. Conclusions
Our results give a geometric interpretation to a number of structures involved in the BFV/BV formalism; the interpretation of the BRST cohomology in terms of the constraint surface geometry can thus be extended in terms of geometry of a “more invariant” object—the zero locus $`𝒵_𝑸`$ that plays the role of the extended constraint surface. Although this is presently limited to the ghost number zero, it would be interesting to extend this interpretation to other ghost numbers. Another interesting application of the zero locus reduction consists in interpreting $`𝒵_𝑸`$ with the induced Poisson bracket in the BV formulation of a pure-gauge model as an extended phase space and the extended Poisson bracket in the BFV formulation of the same model . As noted above, the zero locus reduction applies to finite-dimensional models; it would be interesting to extend it to local field theory, for example in the jet language formulation of the BRST formalism .
#### Acknowledgments
We are grateful to I. A. Batalin, P. H. Damgaard, O. M. Khudaverdyan, and I. V. Tyutin for illuminating discussions. This work was supported in part by the Russian Federation President Grant 99-15-96037. The work of AMS and MAG was supported in part by the RFBR Grant 99-01-00980, and the work of MAG was also supported by the INTAS YSF-98-156. MAG is grateful to P. H. Damgaard for kind hospitality at the Niels Bohr Institute, where a part of this paper was written.
## Appendix A Lie algebra cohomology and the (anti)bracket
Let $`𝔞`$ denote a Lie algebra of dimension $`N`$ and $`𝔐`$ denote an $`𝔞`$-module. We denote by
(A.1)
$$𝔞=\underset{n=0}{\overset{N}{}}^n𝔞$$
the exterior algebra of the vector space $`𝔞`$ and by $`𝖲𝔞`$ the symmetric tensor algebra.
The cohomology complex of $`𝔞`$ with coefficients in the module $`𝔐`$ is
(A.2)
$$𝖢^{}(𝔞,𝔐)=\{𝖧𝗈𝗆(𝔞,𝔐),d\}.$$
Decomposition (A.1) induces the grading $`𝖢^{}(𝔞,𝔐)=_{n=0}^NC^n(𝔞,𝔐)`$, where $`C^n(𝔞,𝔐)=𝖧𝗈𝗆(^n𝔞,𝔐)`$. The differential $`d`$ has the degree $`1`$ and acts as $`d:C^n(𝔞,𝔐)C^{n+1}(𝔞,𝔐)`$ via
(A.3)
$$\begin{array}{c}(da)(g_1,\mathrm{},g_{n+1})=\underset{1i<jn+1}{}(1)^{i+j1}a([g_i,g_j],g_1,\mathrm{},\widehat{g}_i,\mathrm{},\widehat{g}_j,\mathrm{},g_{n+1})+\hfill \\ \hfill +\underset{1in+1}{}(1)^ig_ia(g_1,\mathrm{},\widehat{g}_i,\mathrm{},g_{n+1}),aC^n(𝔞,𝔐).\end{array}$$
We also use the simplified notation $`C^n=C^n(𝔞,𝔐)`$.
We can identify the cohomology complex $`𝖢^{}(𝔞,𝔐)`$ with $`𝔞^{}𝔐`$ as follows.<sup>7</sup><sup>7</sup>7We here assume that the algebras are finite dimensional or graded, $`𝔞=_i𝔞_i`$, with finite dimensional homogeneous spaces $`𝔞_i`$, and $`𝔞^{}`$ is by definition $`𝔞^{}=_i𝔞_i^{}`$, where $`𝔞_i^{}`$ are finite dimensional spaces dual to $`𝔞_i`$. Let $`e_i`$ be a basis in $`𝔞`$, with $`[e_i,e_j]=C_{ij}^ke_k`$. Let also $`\theta ^i`$ be the basis of $`𝔞^{}`$ dual to $`e_i`$. The Grassmann algebra generated by $`\theta ^i`$ is then identified with $`𝔞^{}`$. To every cochain $`x𝖧𝗈𝗆(^n𝔞,𝔐)`$, we associate the element (with the summations implied)
(A.4)
$$\overline{x}=\frac{1}{n!}x(e_{i_1},\mathrm{},e_{i_n})\theta ^{i_1}\mathrm{}\theta ^{i_n}𝔞^{}𝔐.$$
The differential $`d`$ then acts on $`𝔞^{}𝔐`$ as the differential operator
(A.5)
$$d=\frac{1}{2}C_{ij}^k\theta ^i\theta ^j\frac{}{\theta ^k}\theta ^iX_i,$$
where $`X_i:𝔐𝔐`$ is the action of $`e_i𝔞`$ on $`𝔐`$.
We next specialize to the coefficients in $`𝔞`$ (viewed as the adjoint representation $`𝔞`$-module). The complex is then endowed with the Gerstenhaber bracket ,\[26, and references therein\]
$$\mathbf{\{},\mathbf{\}}:C^nC^mC^{n+m1}$$
given by
(A.6)
$$\mathbf{\{}x,y\mathbf{\}}=xy(1)^{(m+1)(n+1)}yx,xC^n,yC^m$$
where
(A.7)
$$(xy)(a_1,\mathrm{},a_{n+m1})=\frac{1}{m!(n1)!}\underset{\sigma P_{n+m1}}{}(1)^\sigma x(a_{\sigma (1)},\mathrm{},a_{\sigma (n1)},y(a_{\sigma (n)},\mathrm{},a_{\sigma (n+m1)}))$$
This makes $`𝖧𝗈𝗆(𝔞,𝔞)`$ into a graded differential Lie algebra.
Let $`\xi _i`$ denote the basis of $`𝔞`$ viewed as an $`𝔞`$-module (equivalently, coordinates on $`𝔞^{}`$). For each cochain $`xC^n`$, we then expand $`\overline{x}`$ from (A.4) as
(A.8)
$$\overline{x}=\frac{1}{n!}\xi _jx^j(e_{i_1},\mathrm{},e_{i_n})\theta ^{i_1}\mathrm{}\theta ^{i_n}$$
and rewrite the Gerstenhaber bracket as
(A.9)
$$\mathbf{\{}\overline{x},\overline{y}\mathbf{\}}=\overline{x}\overline{y}(1)^{(k+1)(l+1)}\overline{y}\overline{x},\overline{x}\overline{y}=\overline{x}\frac{\stackrel{}{}}{\theta ^i}\frac{}{\xi _i}\overline{y},xC^k,yC^l,$$
where $`\frac{\stackrel{}{}}{\theta ^i}`$ is the (right) derivative in the Grassmann algebra and the $`\frac{}{\xi _i}`$ operation is defined on the elements of form (A.8) as the contraction with the element $`\xi _i^{}`$ of the dual basis in $`𝔞^{}`$. The differential then becomes
(A.10)
$$d=\mathbf{\{}\frac{1}{2}C_{ij}^k\xi _k\theta ^i\theta ^j,\mathbf{\}}=\frac{1}{2}C_{ij}^k\theta ^i\theta ^j\frac{}{\theta ^k}\xi _k\theta ^iC_{ij}^k\frac{}{\xi _j}.$$
On the elements $`\underset{¯}{a}`$ as in (A.8), the second term represents the adjoint action (in accordance with the above choice $`𝔐=𝔞`$). Equation (A.9) suggests the interpretation of a Poisson/Batalin–Vilkovisky bracket. As it stands, however, (A.9) can be neither of these, since no associative supercommutative multiplication has been defined on the cochains.<sup>8</sup><sup>8</sup>8Superficially, the bracket in (A.9) has the grade $`1`$ since it maps as $`C^m\times C^nC^{m+n1}`$, however the gradings of all the terms in the complex can be shifted by $`1`$, after which the bracket becomes a grade-$`0`$ operation. On the other hand, an associative graded commutative multiplication defined on the complex would fix the grading, and (A.9) would become either the Batalin–Vilkovisky or the Poisson bracket. There are two remarkable possibilities to embed $`𝖢^{}(𝔞,𝔐)=𝖢^{}(𝔞,𝔞)`$ into a complex endowed with a multiplication: the complex
(A.11)
$$𝖢^{}(𝔞,𝖲𝔞)=𝔞^{}𝖲𝔞$$
corresponding to the BV quantization, or the complex
(A.12)
$$𝖢^{}(𝔞,𝔞)=𝔞^{}𝔞$$
corresponding to the BFV quantization. Geometrically, these two possibilities correspond to even and odd QP manifolds (see Definition 2.1).
Choosing $`𝔐=𝖲𝔞`$, we have the complex $`_{m,n}𝖧𝗈𝗆(^m𝔞,𝖲^n𝔞)`$, which can be viewed as the associative supercommutative algebra generated by the variables $`\theta ^i`$ and $`\xi _j`$ satisfying $`\xi _i\xi _j\xi _j\xi _i=0`$, $`\theta ^i\theta ^j+\theta ^j\theta ^i=0`$, and $`\theta ^i\xi _j\xi _j\theta ^i=0`$.<sup>9</sup><sup>9</sup>9These relations between $`\theta `$ and $`\xi `$ variables correspond to the case (tacitly implied in most of our formulae) where $`𝔞`$ is a Lie algebra, not a superalgebra; then the Grassmann parities are simply $`𝗉(\xi _i)=0`$ and $`𝗉(\theta ^i)=1`$. However, if $`𝔞`$ is a Lie superalgebra, let $`𝗉(e_i)=\epsilon _i`$ be the Grassmann parities of its generators. Then $`𝗉(\xi _i)=\epsilon _i`$ and $`𝗉(\theta ^i)=\epsilon _i+1`$, and therefore, $`\xi _i\xi _j(1)^{\epsilon _i\epsilon _j}\xi _j\xi _i=0`$, $`\theta ^i\theta ^j(1)^{(\epsilon _i+1)(\epsilon _j+1)}\theta ^j\theta ^i=0`$, and $`\xi _i\theta ^j(1)^{\epsilon _i(\epsilon _j+1)}\theta ^j\xi _i=0`$. It then follows that (A.9) can be extended to an odd bracket on this complex. The differential extends to $`𝖧𝗈𝗆(𝔞,𝖲𝔞)`$ by the same formula $`d=\mathbf{\{}C_0,\mathbf{\}}`$, $`C_0=\frac{1}{2}C_{ij}^k\xi _k\theta ^i\theta ^j`$. The complex is endowed with the grading known as the ghost number in the BV quantization or as the Weyl complex grading in homology theory: for a cochain $`x𝖧𝗈𝗆(^m𝔞,𝖲^n𝔞)`$, one has $`\mathrm{gh}(x)=m2n`$.
On the other hand, taking the coefficients to be the exterior algebra $`𝔞`$, we can extend (A.9) to an even bracket. With $`𝔞`$ identified with the algebra generated by $`\xi _i`$ viewed as anticommuting variables (with obvious modifications in the case where $`𝔞`$ is a Lie superalgebra, see footnote 9), the bracket becomes the Poisson bracket on the space $`𝔞^{}𝔞`$ (which is identified with functions of $`\theta ^i`$ and $`\xi _j`$; we also assume that $`\xi _i\theta ^j+\theta ^j\xi _i=0`$ in addition to $`\xi _i\xi _j+\xi _j\xi _i=0`$). The ghost number grading on this complex taken from the BFV quantization is $`\mathrm{gh}(x)=mn`$ for an element $`x𝖧𝗈𝗆(^m𝔞,^n𝔞)`$.
The coefficients can be further extended (cf. ) by $`𝔐=_{}`$, the algebra of smooth functions on a manifold $``$ such that $`𝔞`$ acts on $`_{}`$ by derivations (vector fields on $``$). We write $`X_i`$ for the image of the basis elements of $`𝔞`$ in $`\mathrm{Vect}_{}`$. In accordance with the BRST paradigm, one wishes the vector fields representing the action of $`𝔞`$ on $``$ to be Hamiltonian with respect to a bracket structure. For even $`\xi _i`$, this can be achieved by replacing $``$ with the odd cotangent bundle $`\mathrm{\Pi }T^{}`$ and, thus, the algebra $`_{}`$ with the algebra $`_{\mathrm{\Pi }T^{}}`$ of smooth functions on the odd cotangent bundle. Then each vector field $`V=V^a\frac{}{q^a}`$ on $``$ is generated by the canonical antibracket structure on $`\mathrm{\Pi }T^{}`$; the action of the basis elements $`X_i=X_i^a\frac{}{q^a}`$ on functions is given by the antibracket
(A.13)
$$X_iF=\mathbf{\{}X_i^aq_a^{},F\mathbf{\}},F_{},$$
with $`q_a^{}`$ being the standard coordinates on the fibers of $`\mathrm{\Pi }T^{}`$ (and the standard antibracket given by $`\mathbf{\{}q^a,q_b^{}\mathbf{\}}=\delta _b^a`$).
For odd $`\xi _i`$, similarly, we can consider the functions $`_T^{}`$ on the cotangent bundle, which allows the action of $`𝔞`$ to be implemented by the bracket on $`_T^{}`$ (the same formula (A.13) for the bracket, where now $`q_a^{}`$ are the canonical coordinates on the fibers of $`T^{}`$).
We note, however, that the differential
(A.14)
$$d=\mathbf{\{}\frac{1}{2}C_{ij}^k\xi _k\theta ^i\theta ^j,\mathbf{\}}+\theta ^i\mathbf{\{}X_i^aq_a^{},\mathbf{\}}$$
in either of the complexes
(A.15) $`𝖢_{\mathrm{odd}}^{}(𝔞,)=`$ $`𝖢^{}(𝔞,𝖲𝔞)_{\mathrm{\Pi }T^{}},`$
(A.16) $`𝖢_{\mathrm{even}}^{}(𝔞,)=`$ $`𝖢^{}(𝔞,𝔞)_T^{}`$
is not compatible with the bracket. Remarkably, the compatibility can be achieved by changing the differentials such that (A.15) and (A.16) become the well-known BV and BFV complexes used in the Lagrangian and Hamiltonian quantization of gauge theories. The term to be added to the differential is the Koszul differential involving precisely the same “auxiliary” variables $`\xi _i`$ that were originally introduced to rewrite the Gerstenhaber bracket in the “geometric” form.
To conclude, we note that we have given a homological interpretation of the structures appearing in the BRST quantization in the example of a Lie algebra structure (i.e., in the case where the constraints or gauge generators form a Lie algebra). In the most general setting, the BRST charge and the master action in the BFV and BV cases, respectively, can be considered as the generating functions for the $`L_{\mathrm{}}`$ algebras (see also ). From this general standpoint, the Lie algebra structure appears as a particular case. |
warning/0001/hep-ph0001274.html | ar5iv | text | # Electroweak Bubble Wall Friction: Analytic Results
## I Introduction
Electroweak baryogenesis is the name for the production of the baryon number asymmetry of the universe at the electroweak epoch. It is possible in extensions of the standard model where the three famous Sakharov conditions can be met:
1. Baryon number violation is efficient,
2. The discrete symmetries C and CP are violated, and
3. There is a departure from equilibrium, coincident with baryon number violation turning off.
All three of these criteria potentially exist in the standard model. Baryon number violation is not only present, but efficient . The expansion of the universe provides departure from equilibrium. And of course, C and CP are known not to be true symmetries.
Whether or not electroweak baryogenesis can explain the size of the observed baryon number of the universe, which is about
$$\frac{\mathrm{baryons}}{\mathrm{entropy}}(2\mathrm{\hspace{0.25em}7})\times 10^{11},$$
(1)
is a more detailed question. In the minimal standard model, both the available CP violation , and the departure from equilibrium , appear to be grossly insufficient. However the question is quite open in extensions of the standard model.
Baryogenesis at a first order electroweak phase transition is a complicated process, and requires understanding several things. First, we must be able to compute the strength of the phase transition. The tools for doing this are now well developed . Next, we must be able to compute the efficiency of baryon number violation; here the tools are also well developed . Finally one must be able to compute the microscopic dynamics of baryon number carrying excitations in the presence of an electroweak phase interface, henceforth called the “bubble wall” . This problem is not well under control, but the papers quoted show that the results strongly depend on another factor, which is the velocity of propagation of the bubble wall. This paper will discuss the computation of this bubble wall velocity, which is important for baryogenesis, and is also interesting as an example where such a dynamical quantity can be computed from first principles.
A fairly substantial literature already exists on the electroweak bubble wall velocity. Since a bubble wall liberates latent heat as it propagates, hydrodynamic considerations are potentially important. Hydrodynamic considerations appear in . The conclusion is that, if the friction on the bubble wall is small, then the hydrodynamics are important; but if the bubble friction is large, so the wall velocity is small, then all that matters is the general rise in temperature from the bubbles in aggregate as the transition proceeds.
There is also a literature on the friction the bubble wall feels . The paper of Khlebnikov shows how the friction is related to the self-energy of the zero mode of the bubble wall. That is, the friction arises from the back-reaction on the wall of the disturbance from equilibrium of excitations, induced by the motion of the wall. Most of the papers quoted study this back reaction by treating the excitations with kinetic theory. The exception is , where it is argued that infrared SU(2) gauge field and Higgs field excitations are most important, and that it is more appropriate to treat them as classical fields, which can be done nonperturbatively on the lattice.
Both the kinetic descriptions, and classical nonperturbative treatment, missed one important piece of physics, however, which as we will see leads them to be incorrect parametrically. That is, they miss the physics of screening and Landau damping, which dominates the dynamics of infrared gauge fields. The importance of this physics has been pointed out by Arnold, Son, and Yaffe, , in the context of determining the baryon number violation rate, and has been further discussed in . The central result is that the SU(2) gauge field $`A`$, instead of evolving under (classical) equations of motion of the form (in temporal gauge, ignoring nonlinearities)
$$\frac{d^2A}{dt^2}=(k^2+m^2)A,$$
(2)
with $`m^2=(g^2/4)\varphi ^2`$ the mass squared induced by a Higgs condensate, instead undergoes overdamped evolution,
$$\frac{\pi m_\mathrm{D}^2}{4k}\frac{dA}{dt}=(k^2+m^2)A+\mathrm{noise},$$
(3)
with $`m_\mathrm{D}^2`$ the Debye mass squared, $`m_\mathrm{D}^2=(11/6)g^2T^2`$ in the standard model. In this letter we will see what consequences this has for the kinetic description of the bubble wall friction.
## II Ingredients
For concreteness we will work here with the minimal standard model, even though baryogenesis in that model is ruled out. Since only the gauge fields are overdamped, and since only the SU(2) and U(1) fields have large interactions with the Higgs condensate, the extension to models with more fields, such as the minimal supersymmetric standard model, should be straightforward. From here on we will also neglect the U(1) field, which is the same as taking $`g^{}g`$. This is probably a reasonable approximation. Going beyond it would increase the bookkeeping but would not make our treatment substantially more difficult.
This work will be strictly analytic and strictly based on parametric expansions in formally small quantities. In order to obtain an electroweak phase transition of a strength which can be analyzed perturbatively, we will take
$$g^3\lambda g^2/\mathrm{log}(1/g),$$
(4)
where $`\lambda `$ is the Higgs self-coupling in the same normalization as used in . As we see in a moment, the gauge field condensate is parametrically $`\varphi _0g^3T/\lambda `$, so the induced gauge field masses are $`g\varphi _0/2g^4T/\lambda `$. The longitudinal gauge fields have a Debye mass $`gTg^4T/\lambda `$ and can be neglected ; the one loop effective potential is then approximately
$$V_{1\mathrm{loop}}(\varphi )=\frac{m^2(T)}{2}\varphi ^2\frac{g^3T}{16\pi }\varphi ^3+\frac{\lambda }{4}\varphi ^4.$$
(5)
A broken minimum exists if $`dV/d\varphi =0`$ at some nonzero $`\varphi _0`$. The broken minimum is degenerate with the minimum at $`\varphi =0`$ at
$$m^2(T_{\mathrm{eq}})=\frac{g^6T^2}{128\pi ^2\lambda },$$
(6)
and the value of $`\varphi _0`$ is
$$\varphi _0(T_{\mathrm{eq}})=\frac{g^3T}{8\pi \lambda }.$$
(7)
The $`\varphi `$ profile of the electroweak bubble wall at $`T_{\mathrm{eq}}`$ is given at leading order by
$$\varphi (z)=\frac{\varphi _0}{2}\left[1+\mathrm{tanh}\frac{z}{L}\right],L=\frac{2}{m(T_{\mathrm{eq}})}.$$
(8)
Here $`z`$ is a space coordinate orthogonal to the bubble wall.
The gauge field mass times the wall width,
$$Lm_W=L\frac{g\varphi }{2}=\frac{16\pi \sqrt{2\lambda }}{g^3T}\frac{g^4T}{16\pi \lambda }\frac{\varphi }{\varphi _0}=\frac{\varphi }{\varphi _0}\sqrt{\frac{2g^2}{\lambda }}1,$$
(9)
is large, and it is therefore possible to treat gauge field excitations, in the broken phase and inside the bubble wall, with kinetic theory. The kinetic theory description, which amounts to taking the Higgs field background as approximately homogeneous and expanding in its gradients, breaks down at the small $`z`$ (symmetric phase) edge of the bubble wall, when $`(\varphi /\varphi _0)\sqrt{\lambda /g^2}`$. In fact, since $`\lambda =(g^2/8)`$ is the condition for $`m_H=m_W`$, the kinetic description works fairly well at remarkably large Higgs masses.
All of the equations above are valid at leading order in $`\lambda /g^2`$ or $`g^3/\lambda `$. when the former breaks down, higher loop corrections and corrections from Higgs loops (neglected here) become important. When the latter breaks down, the longitudinal gauge fields become important. We also need $`g^3\lambda `$ to ensure that the relevant gauge fields will be overdamped.
## III Friction
The friction on an electroweak bubble wall is defined as the excess pressure on the wall (directed towards the broken phase), over the equilibrium value;
$$\mathrm{friction}=PP_{\mathrm{eq}}=P+V(\varphi =0)V(\varphi _0)P+\mathrm{\Delta }V.$$
(10)
We expect the friction to depend on the bubble wall velocity, with the condition $`P=0`$ determining the steady state bubble wall velocity $`v_w`$, which is what we want to know. We can define a linear response friction coefficient $`\eta `$ as a limit
$$\eta \underset{\mathrm{\Delta }V0}{lim}\frac{\mathrm{\Delta }V}{v_w}=\frac{\mathrm{steady}\mathrm{state}\mathrm{friction}}{v_w}.$$
(11)
This is what we want to determine.
The friction on the bubble wall depends on the departure from equilibrium of the plasma excitations inside the bubble wall. In a kinetic theory description, the excitations are described by population functions $`f(k,x)`$. We write them as an equilibrium part $`f_0`$,
$$f_0=\frac{1}{\mathrm{exp}(E/T)\pm 1},E=\sqrt{k^2+m^2(x)},$$
(12)
with $`+`$ for fermions and $``$ for bosons, which is the case we will care about, plus a departure from equilibrium $`\delta f`$. The friction a bubble wall feels is, in the kinetic description <sup>*</sup><sup>*</sup>*A derivation of sorts can be found in , but the expression is implicit in the earlier references as well.,
$$\mathrm{friction}=_{\mathrm{}}^{\mathrm{}}𝑑z\underset{\mathrm{DOF}}{}\frac{d^3k}{(2\pi )^3}\frac{d\varphi }{dz}\frac{dm^2}{d\varphi }\frac{dE}{dm^2}\delta f,$$
(13)
where the sum is over degrees of freedom which get a mass from the Higgs field. This expression has a clear intuitive meaning; it is the sum over excess particles of $`(dE/dz)`$, the force the wall exerts on them. It remains to determine $`\delta f`$ and evaluate the integral.
The friction will be dominated by the gauge boson contribution. The contribution from Higgs bosons is smaller because their $`dm^2/d\varphi `$ is smaller by $`\lambda /g^2`$, and because their evolution is not overdamped. The fermionic contributions, such as that from the top quark, are smaller because Fermi-Dirac statistics lack the infrared divergence of Bose-Einstein statistics, so their $`\delta f`$ has much weaker infrared behavior. The parametric argument appears in , and a reasonable estimate of friction from top quarks appears in ; it proves numerically smaller than what we find below. The case where there is a light scalar top is more difficult and we do not consider it; for the MSSM the friction we find should be viewed as a lower bound rather than a tight estimate.
As discussed above, the gauge fields undergo overdamped evolution given by Eq. (3). Since $`fA^2`$, the equation for $`f`$ is
$$\frac{\pi m_\mathrm{D}^2}{8k}\frac{df}{dt}=E^2f+\mathrm{noise},$$
(14)
where the noise is of the right size to ensure that, for $`m^2`$ time independent, $`f`$ will approach $`f_0`$; so averaging over the noise,
$$(k^2+m^2)f+\mathrm{noise}E^2\delta f.$$
(15)
Also, $`df/dt=df_0/dt+d(\delta f)/dt`$. At small $`v_w`$, the limit we are interested in, $`\delta ff_0`$, and $`d(\delta f)/dt`$ may be dropped. Further,
$$\frac{df_0}{dt}=\frac{d\varphi }{dt}\frac{dm^2}{d\varphi }\frac{dE}{dm^2}\frac{df_0}{dE}=v_w\frac{d\varphi }{dz}\frac{dm^2}{d\varphi }\frac{1}{2ET}f_0(1+f_0).$$
(16)
Therefore, the departure from equilibrium of a gauge boson degree of freedom is
$$\delta f=\frac{\pi m_\mathrm{D}^2v_w}{16kE^3T}f_0(1+f_0)\frac{d\varphi }{dz}\frac{dm_W^2}{d\varphi }.$$
(17)
Note that transport plays no role in setting $`\delta f`$; this is because the gauge fields are overdamped. Substituting this into Eq. (13), and noting that there are 6 species of transverse $`W`$ bosons (3 flavors times 2 spins), gives
$$\mathrm{friction}=\frac{6\pi v_wm_\mathrm{D}^2}{8}_{\mathrm{}}^{\mathrm{}}𝑑z\left(\frac{d\varphi }{dz}\frac{dm_W^2}{d\varphi }\right)^2\frac{d^3k}{(2\pi )^3}\frac{f_0(1+f_0)}{4kE^4T}.$$
(18)
Since $`f_0`$ is monotonically decreasing, the momentum integral is infrared dominated, cut off by the nonvanishing $`W`$ boson mass. Therefore it is appropriate to make the approximation, valid in the infrared, that $`f_01+f_0T/E`$, and evaluate the integral;
$$\frac{d^3k}{(2\pi )^3}\frac{f_0(1+f_0)}{4kE^4T}_0^{\mathrm{}}\frac{k^2dk}{2\pi ^2}\frac{T}{4kE^6}=\frac{T}{16\pi ^2}_{m_W^2}^{\mathrm{}}\frac{d(E^2)}{E^6}=\frac{T}{32\pi ^2m_W^4}.$$
(19)
The friction, using $`m_W=g\varphi /2`$, is then
$$\mathrm{friction}=v_w\frac{3m_\mathrm{D}^2T}{32\pi }_{\mathrm{}}^{\mathrm{}}𝑑z\left(\frac{d\varphi }{dz}\right)^2\frac{1}{\varphi ^2}.$$
(20)
Writing $`(d\varphi /dz)dz=d\varphi `$, and using Eq. (8) to write
$$\frac{d\varphi }{dz}=\frac{2\varphi (\varphi _0\varphi )}{L\varphi _0},$$
(21)
Eq. (20) now gives
$$\eta =\frac{3m_\mathrm{D}^2T}{16\pi L}\times \left(_0^{\varphi _0}\frac{(\varphi _0\varphi )d\varphi }{\varphi _0\varphi }=_0^1\frac{(1x)dx}{x}\right).$$
(22)
There is a log divergence arising from the symmetric phase side of the bubble wall. The log will be cut off where the first approximation used to derive Eq. (20) breaks down. The perturbative expansion breaks down when $`m_Wg^2T`$, or $`\varphi gT`$, which is at $`(\varphi /\varphi _0)=(\lambda /g^2)`$. The kinetic theory description breaks down at $`(\varphi /\varphi _0)=\sqrt{\lambda /g^2}`$, see Eq. (9); this occurs first.
Since the degrees of freedom which dominate the friction are those with $`km`$, when $`m`$ drops below $`1/L`$, it is no longer appropriate to treat the particles as seeing a slowly varying wall. Such degrees of freedom see a wall which is sharper than their wavelength can resolve. For those degrees of freedom with $`kL1`$, we find the friction scales as $`1/L`$. This must go over to an $`L`$ independent value for wavelengths which cannot resolve the thickness of the wall, which means that their contribution is less than the kinetic theory estimate. Hence the log is cut off at $`(\varphi /\varphi _0)\sqrt{\lambda /g^2}`$, the contribution from very infrared degrees of freedom is subdominant. Hence the friction we determine is
$$\eta =\frac{3m_\mathrm{D}^2T}{16\pi L}\left(\mathrm{log}(m_WL)+O(1)\right)=\frac{3}{16\sqrt{2}}\frac{m_D^2}{g^2T^2}\frac{g}{\sqrt{\lambda }}\alpha _w^2T^4\left(\mathrm{log}\frac{g}{\sqrt{\lambda }}+O(1)\right).$$
(23)
The first expression makes no assumptions about the effective potential or wall thickness and should be valid in extensions as well as the standard model.
Now we comment on the parametric form of $`\eta `$. Taking $`\lambda g^2`$ and neglecting logs, the friction coefficient is $`\eta \alpha _w^2T^4`$. We can guess this on dimensional grounds by noting that
$$\left[\eta \right]=\left[\frac{\mathrm{pressure}}{\mathrm{velocity}}\right]=\left[\frac{\mathrm{energy}\times \mathrm{time}}{\mathrm{length}^4}\right],$$
(24)
and that $`\eta `$ arises from infrared gauge field physics. Such physics has a natural energy scale $`T`$, a natural length scale $`1/\alpha _wT`$, and a natural time scale $`1/(\alpha _w^2T)`$ ; so on dimensional grounds we should have anticipated $`\eta \alpha _w^2T^4`$. This is to be contrasted with the pressure driving the bubble wall, which by the same parametric estimates must be
$$\left[P\right]=\left[\frac{\mathrm{energy}}{\mathrm{length}^3}\right],P\alpha _w^3T^4.$$
(25)
Hence, the bubble wall velocity is parametrically $`v_w\alpha _w`$.
## IV Bödeker’s Effective Theory
The reason we considered the case $`\lambda g^2/\mathrm{log}(1/g)`$, rather than $`\lambda g^2`$, is because in the parametric regime
$$\frac{g^2}{\mathrm{log}(1/g)}\lambda g^2$$
(26)
The gauge bosons with $`mg\varphi _0/2`$ do not obey Eq. (3); instead Bödeker’s effective theory is applicable ;
$$\sigma \frac{dA}{dt}=E^2A+\mathrm{noise},\sigma =\frac{2\pi m_\mathrm{D}^2}{3g^2T\mathrm{log}(1/g)},$$
(27)
up to corrections suppressed by $`\mathrm{log}(1/g)`$. In this case the derivation proceeds analogously, but the behavior is slightly less infrared dominated, and no log occurs in the $`\varphi `$ integral. The final expression is
$$\eta =\frac{1}{256\sqrt{2}\mathrm{log}(1/g)}\left(\frac{g}{\sqrt{\lambda }}\right)^3\left(\frac{m_D^2}{g^2T^2}\right)\alpha _w^2T^4.$$
(28)
This result is in parametric agreement with what we found above. It is only applicable in an extremely narrow parametric range, and receives corrections suppressed only by $`(\lambda /g^2)`$ (from the loop expansion) and $`g^2\mathrm{log}(1/g)/\lambda `$ (from the breakdown of Bödeker’s effective theory).
## V Conclusion: Value of the Wall Velocity
Since we only have the friction at leading log, which means with at least a factor of $`2`$ error, we will do with a fairly crude estimate of the pressure which drives the bubble. In references the nucleation temperature is estimated as occurring at $`m^2(T_{\mathrm{nuc}})=0.8m^2(T_{\mathrm{eq}})`$. For this value, using Eqs. (5) and (6), the pressure driving the bubble wall is approximately $`P=.94\times (g^6/1024\pi \lambda ^3)\alpha _w^3T^4`$, and the wall velocity is about
$$v_w=\frac{P}{\eta }.0012\left(\frac{g}{\sqrt{\lambda }}\right)^5\frac{\alpha _w}{\mathrm{log}(g/\sqrt{\lambda })},$$
(29)
which, for $`m_Hm_W`$ or $`\lambda =g^2/8`$, and estimating $`\mathrm{log}(g/\sqrt{\lambda })1`$, is $`v_w\alpha _w/4<0.01`$. For $`\lambda /g^2.04`$, as required to give a sufficiently strong phase transition, $`v_w0.1`$. The numerical value for the bubble wall velocity is expected to be small. The numerical estimate from the result determined using Bödeker’s effective theory is similar; $`v_w0.2`$ for $`\lambda /g^2=.04`$. Neither estimate is very reliable because the constant under the log has not been determined; but the conclusion $`v_w1`$ is clear.
The friction we find is larger, and the value of $`v_w`$ smaller, than in previous literature. In particular we find a much larger friction than the numerical results of indicate. This is partly because there we did not include hard thermal loop effects for the gauge fields, but overly aggressive data fitting may have contributed. It would be interesting to make a new numerical analysis using the techniques developed in .
We should also briefly comment on how the result may change in extensions to the standard model and beyond leading order in $`v_w`$. Beyond first order in $`v_w`$, a few effects become important. We may not be able to neglect $`\delta f`$ next to $`f_0`$ in getting Eq. (17); $`\delta f`$, and the friction, will be larger. Also, the frictive pressure will change the bubble wall shape. Since most of the friction is on the symmetric phase side of the wall, the wall will become narrower, see , which also increases the friction. Finally, in extensions to the standard model (which are the only viable candidates for baryogenesis because they can provide both a strong phase transition and the heavy Higgs boson required by experiment), the bubble wall is typically thinner, because the Higgs mass is larger. This increases the friction, even before considering new contributions from extra light bosons such as a light scalar top. For these reasons we anticipate that electroweak bubble wall velocities are quite generally much less than 1. |
warning/0001/hep-th0001114.html | ar5iv | text | # 1. Introduction
## 1. Introduction
One of the remarkable features of pure three dimensional gravity is its reformulation as a Chern-Simons theory of gauge connections . The precise relationship between the two theories is well understood when the spacetime has no boundaries. However, since the discovery of the BTZ (Bañados-Teitelboim-Zanelli) black hole solution of three dimensional gravity , there has been a new interest in this subject. This is partly in an attempt to provide a statistical mechanical interpretation of the black hole entropy . The main idea behind this approach is, roughly, to treat the horizon of the black hole as a two dimensional boundary of spacetime. This argument is based on some quantum mechanical considerations of black holes. As a consequence, one is forced to deal with a theory of gravity in the presence of boundaries. Imposing then appropriate boundary conditions would lead to physical observables that live on the boundary. After quantization, these extra degrees of freedom (which are absent if the manifold has no boundaries) would correctly account for the black hole entropy.
At this point one might wonder about the precise nature of the two dimensional theory that gives rise to these observables. It has been known for a long time that a Chern-Simons theory, defined on a compact surface, is intimately related to a two dimensional Wess-Zumino-Witten-Novikov (WZWN) model . This connection is in fact made at the quantum level and links the physical Hilbert space of the Chern-Simons theory to the conformal blocks of the WZWN model. If the three dimensional manifold is with boundaries, then the equivalence between the two theories relies on the breaking of the gauge symmetry of the Chern-Simons theory at the boundary . However, in the treatement of three dimensional gravity as a Chern-Simons theory with boundaries, one is forced to preserve gauge invariance as this is indeed what replaces diffeomorphism invariance in gravity .
The usual way to proceed in maintaining gauge invariance is as follows: One imposes some boundary conditions on the gauge fields. This requires then the introduction of a corresponding boundary action in order to have a well-defined variational problem. However, the combined action in the bulk and on the boundary is not gauge invariant. A remedy for this problem consists in introducing group-valued fields $`g`$ living on the boundary. The resulting theory, for very particular boundary conditions, is a Chern-Simons theory coupled to a WZWN model . Another point of view in arriving to this conclusion is presented in . The quantisation of the WZWN theory is then believed to account for the entropy of the black hole. It should be emphasised that different boundary conditions yield different conformal field theories on the boundary . A BRST treatment leads also to the same conclusions , where different gauge fixings give different theories.
On the other hand, the conventional method in showing the equivalence of Chern-Simons gauge theory and the WZWN model follows a different path . In this procedure one does not introduce (by hand) new dynamical degrees of freedom on the boundary, namely the WZWN fields. In fact, one starts from the Chern-Simons action
$$I_{\mathrm{SC}}\left(A\right)=\frac{k}{4\pi }_{}\mathrm{d}^3yϵ^{\mu \nu \rho }\mathrm{tr}\left(A_\mu _\nu A_\rho +\frac{2}{3}A_\mu A_\nu A_\rho \right)$$
(1)
We assume here that $`=𝐑\times 𝚺`$, where $`𝚺`$ is a disc whose boundary we denote $`𝚺`$ and whose radial and angular coordinates are $`r`$ and $`\theta `$, respectively. The time coordinate $`\tau `$ parametrises $`𝐑`$ and $`ϵ^{\tau r\theta }=1`$. It is clear that gauge invariance, which holds only at the level of the partition function if $``$ has no boundaries, is broken by the presence of the boundary $`=𝐑\times 𝚺`$. The only gauge invariance left is the one for which the gauge parameters reduce to the identity at the boundary. Expanding the action we find
$$I_{\mathrm{SC}}\left(A\right)=\frac{k}{4\pi }_{}\mathrm{d}^3y\mathrm{tr}\left(2A_\tau F_{r\theta }A_r_\tau A_\theta +A_\theta _\tau A_r\right)+\frac{k}{4\pi }_{}\mathrm{d}^2x\mathrm{tr}\left(A_\theta A_\tau \right),$$
(2)
where $`F_{r\theta }=_rA_\theta _\theta A_r+[A_r,A_\theta ]`$.
If one ignores the boundary action, then $`A_\tau `$ is a non-dynamical field which forces, in the bulk, the constraint $`F_{r\theta }=0`$. The solution to this zero curvature condition is $`A_i=L^1_iL`$, $`(i=r,\theta )`$, for some group element $`L`$. Upon substitution, the Chern-Simons action reduces to a WZWN model. It is not clear though how to justify the fact that the boundary term is not taken into account. The usual given explanation resides in choosing a gauge for which $`A_\tau =0`$. However, if one treats $`A_\tau `$ in the same manner everywhere then one has the constraint $`A_\theta =0`$ on the boundary. Solving simultaneously the bulk and the boundary constraints would put conditions on $`L`$ at the boundary. Furthermore, one might choose to impose, by hand, some other boundary conditions. This in turn introduces a boundary action involving the gauge fields $`A_\mu `$. Here also different boundary conditions would yield different theories, not necessarily of the WZWN type. In conclusion, the precise connection between Chern-Simons theory, on a manifold with boundaries, and the WZWN model is far from being transparent.
We notice that in all the above mentioned methods, the starting point is the Chern-Simons theory. The WZWN model is obtained either by a direct coupling or by a particular parametrisation of the gauge fields. The aim of this paper is to clarify, at the level of the action and in a classical manner, the nature of the relationship between the WZWN model and Chern-Simons theory. Our starting point is the WZWN model itself. We find that the WZWN theory is dual to a topological BF theory coupled to a Chern-Simons theory in the bulk. The boundary action is unique for this equivalence to hold and no boundary conditions are imposed by hand. This is the same guiding principle for possible boundary terms as that presented in . Our approach makes use of the techniques of non-Abelian T-duality transformations in non-linear sigma models . Here one relies on the existence of isometries and their gauging. In the case of the WZWN model there are two types of isometries. The first consists of the isometries for which the gauge fields live entirely on the boundary and are at most quadratic in the action. The second corresponds to those isometries for which the gauge fields live on the boundary as well as on the bulk. It is this last category of isometries which is used in this study. It has the advantage of allowing for a first order formulation of the WZWN model in terms of gauge fields and Lagrange multipliers where the original fields do not appear anymore. Finally, we specialise to the case of the Lie algebras $`SO(2,1)`$ and $`SO(2,1)\times SO(2,1)`$ and explore their relation to three dimensional gravity. We show that, in general, the two theories do not coincide. This is due to the presence of the BF theory action which is usually ignored in the literature when comparing three dimensional gravity to the WZWN model.
## 2. First order formulation of the WZWN model
The action for the WZWN model defined on the group manifold $`_𝒢`$, based on the Lie algebra $`𝒢`$, is given by
$`S_{\mathrm{WZWN}}\left(g\right)`$ $`=`$ $`{\displaystyle \frac{k}{8\pi }}{\displaystyle _{}}\mathrm{d}^2x\sqrt{|\gamma |}\gamma ^{\mu \nu }\mathrm{tr}\left(g^1_\mu g\right)\left(g^1_\nu g\right)+\mathrm{\Gamma }\left(g\right)`$
$`\mathrm{\Gamma }\left(g\right)`$ $`=`$ $`{\displaystyle \frac{k}{12\pi }}{\displaystyle _{}}\mathrm{d}^3yϵ^{\mu \nu \rho }\mathrm{tr}\left(g^1_\mu g\right)\left(g^1_\nu g\right)\left(g^1_\rho g\right)`$ (3)
where $`g_𝒢`$ and $``$ is a three dimensional ball whose boundary is the two dimensional surface $``$. The metric on this two dimensional worldsheet is denoted by $`\gamma _{\mu \nu }`$. The remarkable thing about this action is that a variation of the type $`gg+\delta g`$ leads to a change in the action, $`\delta S_{\mathrm{WZWN}}`$, which is an integral over the boundary $``$ only. The equations of motion are obtained without a need to impose any boundary conditions on the field $`g`$ or the variation $`\delta g`$. This property will serve as a guiding principle for our analyses when gauge fields are included. We shall adopt the philosophy of not imposing any boundary conditions on the fields but rather let the equations of motion determine the behaviour of the fields everywhere. This point will be explained in details below.
As it is well-known, the WZWN action has the global symmetry
$$gLgR,$$
(4)
where $`L`$ and $`R`$ are two constant (more precisely chiral) group elements. Our first step in constructing the dual theory is to gauge this symmetry. We, therefore, introduce two Lie algebra-valued gauge functions $`A_\mu `$ and $`\stackrel{~}{A}_\mu `$ transforming as
$`A_\mu `$ $``$ $`LA_\mu L^1_\mu LL^1`$
$`\stackrel{~}{A}_\mu `$ $``$ $`R^1\stackrel{~}{A}_\mu R+R^1_\mu R.`$ (5)
Since the usual minimal coupling of the gauge fields does not lead to an invariant theory, the gauged WZWN action is found by applying Noether’s method . The final result is
$`S_{\mathrm{gauge}}`$ $`=`$ $`S_{\mathrm{WZWN}}\left(g\right)+I_{\mathrm{SC}}\left(A\right)I_{\mathrm{SC}}\left(\stackrel{~}{A}\right)`$ (6)
$`+`$ $`{\displaystyle \frac{k}{4\pi }}{\displaystyle _{}}\mathrm{d}^2x\mathrm{tr}\left[P_+^{\mu \nu }\left(_\mu gg^1A_\nu \right)P_{}^{\mu \nu }\left(g^1_\mu g\stackrel{~}{A}_\nu \right)P_{}^{\mu \nu }\left(A_\mu g\stackrel{~}{A}_\nu g^1\right)\right]`$
$`+`$ $`{\displaystyle \frac{k}{8\pi }}{\displaystyle _{}}\mathrm{d}^2x\sqrt{|\gamma |}\gamma ^{\mu \nu }\mathrm{tr}\left(A_\mu A_\nu +\stackrel{~}{A}_\mu \stackrel{~}{A}_\nu \right).`$
We have defined, for convenience, the two quantities $`P_\pm ^{\mu \nu }=\sqrt{|\gamma |}\gamma ^{\mu \nu }\pm ϵ^{\mu \nu }`$. Notice also the natural appearance of the Chern-Simons actions corresponding to the gauge fields $`A_\mu `$ and $`\stackrel{~}{A}_\mu `$.
There are only two kinds of subgroups of the transformations (4) for which the gauge fields live entirely on the boundary $``$. This corresponds to the situation when the combination $`\left[I_{\mathrm{CS}}\left(A\right)I_{\mathrm{SC}}\left(\stackrel{~}{A}\right)\right]`$ vanishes. The first category are the diagonal subgroups where $`R=L^1`$ and $`A_\mu =\stackrel{~}{A}_\mu `$ with $`L`$ and $`R`$ being Abelian or non-Abelian group elements. The second kind are the axial subgroups for which $`R=L`$ and $`A_\mu =\stackrel{~}{A}_\mu `$, where both $`R`$ and $`L`$ are Abelian group elements. The sort of gauging we are interested in corresponds to taking the gauge functions $`L`$ and $`R`$ to be two independent, Abelian or non-Abelian, group elements. The Chern-Simons parts of the gauged WZWN action (6) are then present.
The next step towards the dual theory consists in casting the WZWN action in a first order formulation. We begin from the following gauge invariant action
$$S_{\mathrm{total}}=S_{\mathrm{gauge}}\frac{k}{4\pi }_{}\mathrm{d}^3yϵ^{\mu \nu \rho }\mathrm{tr}\left(B_\mu F_{\nu \rho }\right)\frac{k}{4\pi }_{}\mathrm{d}^3yϵ^{\mu \nu \rho }\mathrm{tr}\left(\stackrel{~}{B}_\mu \stackrel{~}{F}_{\nu \rho }\right).$$
(7)
Here $`F_{\mu \nu }`$ and $`\stackrel{~}{F}_{\mu \nu }`$ are the two gauge curvatures correponding, respectively, to $`A_\mu `$ and $`\stackrel{~}{A}_\mu `$. The Lie algebra-valued fields $`B_\mu `$ and $`\stackrel{~}{B}_\mu `$ are two Lagrange multipliers transforming as $`B_\mu LB_\mu L^1`$ and $`\stackrel{~}{B}_\mu R^1\stackrel{~}{B}_\mu R`$. The equations of motion of these Lagrange multipliers (or their integration out in a path integral formulation) lead to the constraints $`F_{\mu \nu }=\stackrel{~}{F}_{\mu \nu }=0`$. The solutions to these two equations are, up to gauge transformations, given by $`A_\mu =h^1_\mu h`$ and $`\stackrel{~}{A}_\mu =\stackrel{~}{h}^1_\mu \stackrel{~}{h}`$ for two group elements $`h`$ and $`\stackrel{~}{h}`$. Substituting for $`A_\mu `$ and $`\stackrel{~}{A}_\mu `$ in (7) we find that $`S_{\mathrm{total}}=S_{\mathrm{WZWN}}\left(hg\stackrel{~}{h}^1\right)`$. Therefore, by a change of variables such that $`g^{}=hg\stackrel{~}{h}^1`$ (or equivalently by fixing a gauge such that $`h=\stackrel{~}{h}=1`$) one recovers the original WZWN model. We conclude that the WZWN model in (3) is equivalent to the theory described by the action in (7).
The gauge invariance of the action $`S_{\mathrm{total}}`$ allows one to choose a gauge such that $`g=1`$. Notice that this gauge choice is not possible for diagonal or axial subgroups. Substituting for $`g`$, we obtain the sought first order action
$`S_{\mathrm{first}}`$ $`=`$ $`I_{\mathrm{SC}}\left(A\right)I_{\mathrm{SC}}\left(\stackrel{~}{A}\right){\displaystyle \frac{k}{4\pi }}{\displaystyle _{}}\mathrm{d}^3yϵ^{\mu \nu \rho }\mathrm{tr}\left(B_\mu F_{\nu \rho }\right){\displaystyle \frac{k}{4\pi }}{\displaystyle _{}}\mathrm{d}^3yϵ^{\mu \nu \rho }\mathrm{tr}\left(\stackrel{~}{B}_\mu \stackrel{~}{F}_{\nu \rho }\right)`$ (8)
$``$ $`{\displaystyle \frac{k}{4\pi }}{\displaystyle _{}}\mathrm{d}^2xP_{}^{\mu \nu }\mathrm{tr}\left(A_\mu \stackrel{~}{A}_\nu \right)+{\displaystyle \frac{k}{8\pi }}{\displaystyle _{}}\mathrm{d}^2x\sqrt{|\gamma |}\gamma ^{\mu \nu }\mathrm{tr}\left(A_\mu A_\nu +\stackrel{~}{A}_\mu \stackrel{~}{A}_\nu \right).`$
We should mention that the integration over the Lagrange multipliers would always lead to the WZWN model. This follows from the above explanation upon setting $`g=1`$, where we obtain $`S_{\mathrm{first}}=S_{\mathrm{WZWN}}\left(g^{}\right)`$ with $`g^{}=h\stackrel{~}{h}^1`$. At this point, it is important to emphasise the crucial rôle played by the Lagrange multiplier terms in relating the first order action (8) to the WZWN theory. A further practical manipulation consists in making a field redefinition such that $`Q_\mu =B_\mu +\frac{1}{2}A_\mu `$ and $`\stackrel{~}{Q}_\mu =\stackrel{~}{B}_\mu \frac{1}{2}\stackrel{~}{A}_\mu `$. The first order action takes then the form
$`S_{\mathrm{first}}`$ $`=`$ $`{\displaystyle \frac{k}{4\pi }}{\displaystyle _{}}\mathrm{d}^3yϵ^{\mu \nu \rho }\mathrm{tr}\left[Q_\mu F_{\nu \rho }+\stackrel{~}{Q}_\mu \stackrel{~}{F}_{\nu \rho }{\displaystyle \frac{1}{3}}A_\mu A_\nu A_\rho +{\displaystyle \frac{1}{3}}\stackrel{~}{A}_\mu \stackrel{~}{A}_\nu \stackrel{~}{A}_\rho \right]`$ (9)
$``$ $`{\displaystyle \frac{k}{4\pi }}{\displaystyle _{}}\mathrm{d}^2xP_{}^{\mu \nu }\mathrm{tr}\left(A_\mu \stackrel{~}{A}_\nu \right)+{\displaystyle \frac{k}{8\pi }}{\displaystyle _{}}\mathrm{d}^2x\sqrt{|\gamma |}\gamma ^{\mu \nu }\mathrm{tr}\left(A_\mu A_\nu +\stackrel{~}{A}_\mu \stackrel{~}{A}_\nu \right).`$
In this reformulation the equivalence to the WZWN model is much more transparent and the distinction from pure Chern-Simons theory is evident.
Our first order action (8) has a rich structure in terms of symmetries. Indeed, the Lagrange multiplier terms present in this action correspond to what is commonly known as “$`BF`$ theories”. They are topological theories which have been widely studied (see for a review). A typical characteristic of these theories is their invariance under the finite transformations
$`B_\mu B_\mu +𝒟_\mu \alpha =B_\mu +_\mu \alpha +[A_\mu ,\alpha ]`$
$`\stackrel{~}{B}_\mu \stackrel{~}{B}_\mu +\stackrel{~}{𝒟}_\mu \stackrel{~}{\alpha }=\stackrel{~}{B}_\mu +_\mu \stackrel{~}{\alpha }+[\stackrel{~}{A}_\mu ,\stackrel{~}{\alpha }]`$ (10)
for some two arbitrary local functions $`\alpha `$ and $`\stackrel{~}{\alpha }`$ evaluated in the Lie algebra $`𝒢`$. In our case, this invariance does not hold due to the presence of the boundary. There are, however, two possible ways to restore this invariance. The simplest would be to demande that $`\alpha `$ and $`\stackrel{~}{\alpha }`$ vanish when evaluated at the boundary $``$. The second is to put no restrictions on $`\alpha `$ and $`\stackrel{~}{\alpha }`$ and to supply the first order action (8) with the additional boundary term
$$S_{\mathrm{add}}=\frac{k}{4\pi }_{}\mathrm{d}^2xϵ^{\mu \nu }\mathrm{tr}\left(\lambda F_{\mu \nu }\right)\frac{k}{4\pi }_{}\mathrm{d}^2xϵ^{\mu \nu }\mathrm{tr}\left(\stackrel{~}{\lambda }\stackrel{~}{F}_{\mu \nu }\right).$$
(11)
We associate then to the new fields $`\lambda `$ and $`\stackrel{~}{\lambda }`$ the transformations $`\lambda \lambda \alpha `$ and $`\stackrel{~}{\lambda }\stackrel{~}{\lambda }\stackrel{~}{\alpha }`$. The equations of motion corresponding to $`\lambda `$ and $`\stackrel{~}{\lambda }`$ force the field strength $`F_{\mu \nu }`$ and $`\stackrel{~}{F}_{\mu \nu }`$ to vanish on the boundary $``$. However, these constraints are already imposed by the Lagrange multipliers $`B_\mu `$ and $`\stackrel{~}{B}_\mu `$. In this sense, $`\lambda `$ and $`\stackrel{~}{\lambda }`$ are redundant fields and can be set to zero at anytime. In other words, a gauge fixing of the transformations (10) for which $`\lambda =\stackrel{~}{\lambda }=0`$ can be made. We will nevertheless keep this additional term in our analyses for later use. Of course we could have also chosen to impose appropriate boundary conditions on the fields themselves. However, this is not in the spirit of our procedure as mentioned above.
The Lagrange multiplier terms in the first order action are seperately invariant under the gauge transformations (5). However, the Chern-Simons parts together with the boundary action are not. It is though possible to recover this invariance by putting restrictions on the gauge parameters as given by
$$L|_{}=R|_{}=0$$
(12)
and demande, for quantum invariance, that the topological charges $`\mathrm{\Gamma }\left(L\right)`$ and $`\mathrm{\Gamma }\left(R\right)`$ are integer-valued.
As a matter of fact, not all of the gauge symmetry (5) is exhausted by the gauge fixing conditon $`g=1`$. Indeed, when both $`A_\mu `$ and $`\stackrel{~}{A}`$ are present, the first order action is still invariant everywhere under the gauge transformations
$$A_\mu HA_\mu H^1_\mu HH^1,\stackrel{~}{A}_\mu H\stackrel{~}{A}_\mu H^1_\mu HH^1$$
(13)
with no restrictions on the group element $`H`$ at the boundary. Accordingly, the Lagrange multipliers transform as $`B_\mu HB_\mu H^1`$ and $`\stackrel{~}{B}_\mu H\stackrel{~}{B}_\mu H^1`$. Since the Lagrange multiplier terms are each gauge invariant under this last symmetry, we have therefore a way of making two Chern-Simons actions (more precisely their difference) gauge invariant. Namely, by coupling them at the boundary in the unique manner as in the first order action (8). This presents an alternative to the usual procedure employed, namely by introducing new dynamical fields, in preserving gauge invariance in Chern-Simons theory.
We return now to our main stream to explore duality. The idea of duality is not to use the equations of motion of the Lagrange multipliers (as these would always lead to the original WZWN theory) but use instead those of the gauge fields. In this context, we distinguish two different situations: If the gauge fields belong to the diagonal or axial subgroups then they appear at most in a quadratic form in the action $`S_{\mathrm{first}}`$. Hence, they can be completely eliminated from the action through their equations of motion (which is equivalent to performing the Gaussian integration over these fields in the path integral). The resulting theory (the dual theory) is another non-linear sigma model but with a different target space metric and a different torsion from those encountred in the original WZWN model. This case has been the subject of many investigations and will not concern us here. The second situation, which will be our main interest, deals with the gauging of subgroups for which the presence of the Chern-Simons action is required. The gauge fields are no longer quadratic in the gauged action. Therefore, they cannot be integrated out from the first order action and the dual theory is certainly not a non-linear sigma model. Let us nevertheless stay at the classical level and examin the equations of motion of the gauge fields $`A_\mu `$ and $`\stackrel{~}{A}_\mu `$.
The Lie algebras $`𝒢`$ has generators $`T_a`$ and is specified by the structure constants $`f_{bc}^a`$ and the trace $`\eta _{ab}=\mathrm{tr}\left(T_aT_b\right)`$, where $`\eta _{ab}f_{cd}^b+\eta _{cb}f_{ad}^b=0`$. We do not assume anything on the ivertibility of the invariant bilinear form $`\eta _{ab}`$. The fields of the first order action are decomposed as $`A_\mu =A_\mu ^aT_a`$, $`\stackrel{~}{A}_\mu =\stackrel{~}{A}_\mu ^aT_a`$, $`B_\mu =B_\mu ^aT_a`$ and $`\stackrel{~}{B}_\mu =\stackrel{~}{B}_\mu ^aT_a`$. We start by calculating the equations of motion in the bulk. The variation of the action (8) with respect to the gauge fields yields
$$ϵ^{\mu \nu \rho }\eta _{ab}(F_{\nu \rho }^b+2𝒟_\nu B_\rho ^b)=0,ϵ^{\mu \nu \rho }\eta _{ab}(\stackrel{~}{F}_{\nu \rho }^b2\stackrel{~}{𝒟}_\nu \stackrel{~}{B}_\rho ^b)=0.$$
(14)
Acting with the covariant derivatives $`𝒟_\mu `$ and $`\stackrel{~}{𝒟}_\mu `$ on these equations and using the Bianchi identities, leads to the following consistency relations
$$ϵ^{\mu \nu \rho }\eta _{ab}f_{cd}^bF_{\mu \nu }^cB_\rho ^d=0,ϵ^{\mu \nu \rho }\eta _{ab}f_{cd}^b\stackrel{~}{F}_{\mu \nu }^c\stackrel{~}{B}_\rho ^d=0.$$
(15)
Among the possible solutions to these consistency equations, two are of particular interest to us. The first is provided by taking $`F_{\mu \nu }^a=\stackrel{~}{F}_{\mu \nu }^a=0`$ (we assume for this discussion that $`\eta _{ab}`$ is invertible). These conditions are as if one is using the equations of motion of the Lagrange multiplier fields $`B_\mu `$ and $`\stackrel{~}{B}_\mu `$. Solving these constraints by taking $`A_\mu `$ and $`\stackrel{~}{A}_\mu `$ to be pure gauge fields would eventually yield the original WZWN theory. Furthermore, the vanishing of the gauge curvatures (and hence the disappearance of the Lagrange multiplier terms from the action) are the usual equations of motion of pure Chern-Simons theory. This explains why pure Chern-Simons theory (without the BF terms) could be equivalent, in some particular situations, to the WZWN model.
The second solution is given by
$$B_\mu ^a=G_{\mu \nu }ϵ^{\nu \alpha \beta }F_{\alpha \beta }^a,\stackrel{~}{B}_\mu ^a=G_{\mu \nu }ϵ^{\nu \alpha \beta }\stackrel{~}{F}_{\alpha \beta }^a,$$
(16)
where $`G_{\mu \nu }`$ is the metric on the three dimensional manifold $``$. If we substitute these expressions for $`B_\mu `$ and $`\stackrel{~}{B}_\mu `$ in the equations of motion (14), we obtain
$`\eta _{ab}\left(ϵ^{\mu \nu \rho }F_{\nu \rho }^b+2\sqrt{|G|}G^{\mu \nu }G^{\alpha \beta }𝒟_\alpha F_{\nu \beta }^b\right)=0`$
$`\eta _{ab}\left(ϵ^{\mu \nu \rho }\stackrel{~}{F}_{\nu \rho }^b2\sqrt{|G|}G^{\mu \nu }G^{\alpha \beta }\stackrel{~}{𝒟}_\alpha \stackrel{~}{F}_{\nu \beta }^b\right)=0.`$ (17)
These are the equations of motion in the bulk of a three dimensional Yang-Mills theory in the presence of Chern-Simons terms. This observation might explain the origin of the boundary Kac-Moody algebra found in a Yang-Mills-Chern-Simons gauge theory . Therefore, different dual theories to the WZWN model are obtained depending on which solution one chooses for the consistency equations.
In order for the above variational procedure to be well-defined, one needs to specify the boundary conditions of the problem. We choose not to impose by hand boundary conditions on the fields. We rather let the equations of motion play their full rôle and determine for us the boundary conditions. This is achieved by treating the boundary variations in the same manner as those of the bulk. In other words, we demand that $`\delta A_\mu `$ and $`\delta \stackrel{~}{A}_\mu `$ are arbitrary both on the bulk and on the boundary. When varying with respect to $`A_\mu ^a`$ and $`\stackrel{~}{A}_\mu ^a`$, one obtains the boundary terms
$`{\displaystyle \frac{k}{4\pi }}{\displaystyle _{}}\mathrm{d}^2x\eta _{ab}\left[P_{}^{\rho \mu }\left(\stackrel{~}{A}_\mu ^aA_\mu ^a\right)+2ϵ^{\rho \mu }\left(B_\mu ^a+𝒟_\mu \lambda ^a\right)\right]\delta A_\rho ^b`$
$`{\displaystyle \frac{k}{4\pi }}{\displaystyle _{}}\mathrm{d}^2x\eta _{ab}\left[P_+^{\rho \mu }\left(A_\mu ^a\stackrel{~}{A}_\mu ^a\right)+2ϵ^{\rho \mu }\left(\stackrel{~}{B}_\mu ^a+\stackrel{~}{𝒟}\stackrel{~}{\lambda }^a\right)\right]\delta \stackrel{~}{A}_\rho ^b,`$ (18)
where we have included the contribution due to the action (11) with $`\lambda =\lambda ^aT_a`$ and $`\stackrel{~}{\lambda }=\stackrel{~}{\lambda }^aT_a`$. If we choose not to impose the vanishing of $`\delta A_\mu ^a`$ and $`\delta \stackrel{~}{A}_\mu ^a`$ at the boundary, then the equations of motion on the boundary are
$`\eta _{ab}\left[P_{}^{\rho \mu }\left(\stackrel{~}{A}_\mu ^aA_\mu ^a\right)+2ϵ^{\rho \mu }\left(B_\mu ^a+𝒟_\mu \lambda ^a\right)\right]=0`$
$`\eta _{ab}\left[P_+^{\rho \mu }\left(A_\mu ^a\stackrel{~}{A}_\mu ^a\right)+2ϵ^{\rho \mu }\left(\stackrel{~}{B}_\mu ^a+\stackrel{~}{𝒟}_\mu \stackrel{~}{\lambda ^a}\right)\right]=0.`$ (19)
This determines for us, in a natural way, the behaviour of the gauge fields at the boundary. Of course, these boundary conditions must be compatible with the bulk equations of motion (14). Since, the boundary action was found in a unique fashion, namely through demanding gauge invariance of the WZWN action, we suspect that the two sets of equations are always compatible. This has been checked at least for the case of the BTZ black hole in the absence of BF terms .
## 3. Comparaison with three dimensional gravity
We consider now some special Lie algebras which are relevant to the study of three dimensional gravity. The first of these cases consists in taking the left gauging, $`gLg`$, of the WZWN model. This amounts to setting $`\stackrel{~}{A}=0`$ in the first order action (8). We choose a Lie algebra whose generators $`T_i`$ and $`J_i`$ satisfy
$`[T_i,T_j]=f_{ij}^kT_k,[T_i,J_j]=f_{ij}^kJ_k,[J_i,J_j]=0`$
$`\mathrm{tr}\left(T_iT_j\right)=0,\mathrm{tr}\left(J_iJ_j\right)=0,\mathrm{tr}\left(T_iJ_j\right)=\eta _{ij}.`$ (20)
We will work with the modified first order action in (9) together with the additional boundary action (11) and expand the different fields there according to
$$A_\mu =\omega _\mu ^iT_i+e_\mu ^iJ_i,Q_\mu =\theta _\mu ^iT_i+v_\mu ^iJ_i,\lambda _\mu =\chi ^iT_i+t^iJ_i.$$
(21)
The action for this particular Lie algebra takes then the form
$`S_{\mathrm{first}}^{(1)}`$ $`=`$ $`{\displaystyle \frac{k}{4\pi }}{\displaystyle _{}}\mathrm{d}^3yϵ^{\mu \nu \rho }\eta _{ij}\left[v_\mu ^iF_{\nu \rho }^j+2\theta _\mu ^i𝒟_\nu e_\rho ^j{\displaystyle \frac{1}{2}}f_{kl}^i\omega _\mu ^k\omega _\nu ^le_\rho ^j\right]`$ (22)
$`+`$ $`{\displaystyle \frac{k}{4\pi }}{\displaystyle _{}}\mathrm{d}^2x\eta _{ij}\left[\sqrt{|\gamma |}\gamma ^{\mu \nu }\omega _\mu ^ie_\nu ^jϵ^{\mu \nu }\left(t^iF_{\mu \nu }^j+2\chi ^i𝒟_\mu e_\nu ^j\right)\right],`$
where $`F_{\mu \nu }^i=_\mu \omega _\nu ^i_\nu \omega _\mu ^i+f_{jk}^i\omega _\mu ^j\omega _\nu ^k`$ and $`𝒟_\mu e_\nu ^i=_\mu e_\nu ^i+f_{jk}^i\omega _\mu ^je_\nu ^k`$. We remark that the gauge field component $`e_\mu ^i`$ appears linearly in this action. It has, therefore, the function of a Lagrange multiplier which imposes the bulk contraint $`ϵ^{\mu \nu \rho }\eta _{ij}\left(𝒟_\nu \theta _\rho ^j\frac{1}{4}f_{kl}^j\omega _\nu ^k\omega _\rho ^l\right)=0`$. If a formal solution of the form $`\omega _\mu ^i=O_\mu ^i\left(\theta \right)`$ exists then the action (22) is effectively a BF theory with just the first term present in the bulk. The true fields are, as expected, the original Lagrange multipliers $`\theta _\mu ^i`$ and $`v_\mu ^i`$. This statement will be of use to us when dealing with three dimensional gravity.
The other example we consider is obtained for left and right gauging, $`gLgR`$, of the WZWN model. Here both gauge fields $`A_\mu `$ and $`\stackrel{~}{A}_\mu `$ are kept in the first order action (8). We assume that the Lie algebra $`𝒢`$ consists of two identical copies (left and right) spanned by the generators $`T_i`$ and $`\stackrel{~}{T}_i`$. We have chosen to label the two copies with the same indices. The structure constants are denoted $`f_{jk}^i`$ for both copies and the trace is such that $`\mathrm{tr}\left(T_iT_j\right)=\mathrm{tr}\left(T_i\stackrel{~}{T}_j\right)=\mathrm{tr}\left(\stackrel{~}{T}_i\stackrel{~}{T}_j\right)=\eta _{ij}`$. The two independent gauge fields are decomposed according to
$$A_\mu =\left(\omega _\mu ^i+\alpha e_\mu ^i\right)T_i,\stackrel{~}{A}_\mu =\left(\omega _\mu ^i\alpha e_\mu ^i\right)\stackrel{~}{T}_i.$$
(23)
Similarly, the redefined Lagrange multipliers in (9) are written in the form
$$Q_\mu =(\theta _\mu ^i+\frac{1}{\alpha }v_\mu ^i)T_i,\stackrel{~}{Q}_\mu =(\theta _\mu ^i\frac{1}{\alpha }v_\mu ^i)\stackrel{~}{T}_i.$$
(24)
The action $`S_{\mathrm{first}}`$ for this kind of Lie algebra is given by
$`S_{\mathrm{first}}^{(2)}`$ $`=`$ $`{\displaystyle \frac{k}{2\pi }}{\displaystyle _{}}\mathrm{d}^3yϵ^{\mu \nu \rho }\eta _{ij}[\theta _\mu ^i(F_{\nu \rho }^j+\alpha ^2f_{kl}^je_\nu ^ke_\rho ^l)+2v_\mu ^i𝒟_\nu e_\rho ^j`$ (25)
$``$ $`{\displaystyle \frac{1}{2}}\alpha f_{kl}^i(\omega _\mu ^k\omega _\nu ^le_\rho ^j+{\displaystyle \frac{1}{3}}\alpha ^2e_\mu ^ke_\nu ^le_\rho ^j)]+{\displaystyle \frac{k}{2\pi }}{\displaystyle }_{}\mathrm{d}^2x\eta _{ij}\{\alpha ^2\sqrt{|\gamma |}\gamma ^{\mu \nu }\eta _{ij}e_\mu ^ie_\nu ^j`$
$``$ $`ϵ^{\mu \nu }[\alpha \omega _\mu ^ie_\nu ^j+\chi ^i(F_{\mu \nu }^j+\alpha ^2f_{kl}^je_\mu ^ke_\nu ^l)+2t^i𝒟_\mu e_\nu ^j]\},`$
where we have also expanded the additional fields in (11) as $`\lambda =\left(\chi ^i+\frac{1}{\alpha }t^i\right)T_i`$ and $`\stackrel{~}{\lambda }=\left(\chi ^i\frac{1}{\alpha }t^i\right)\stackrel{~}{T}_i`$. Since the above action is cubic in $`e_\mu ^i`$, it is not merely an effective BF theory.
Another reason behind the choice of these particular Lie algebras resides in their close connection with three dimensional gravity. Before entering into the details of the precise relashionship, let us briefly recall the main features of three dimensional gravity in the Palatini formalism (see for a review). The Einstein-Hilbert action in three dimensions,
$$S^{\mathrm{EH}}=\frac{1}{16\pi \kappa }_{}\mathrm{d}^3y\sqrt{|G|}\left(R2\mathrm{\Lambda }\right)+_{}\mathrm{d}^2x\left(G\right),$$
(26)
can be cast in a first order formalism as
$$S_{\mathrm{Palatini}}^{\mathrm{EH}}=\frac{1}{16\pi \kappa }_{}ϵ_{IJK}\left[F^{IJ}\left(\mathrm{\Omega }\right)E^K\frac{\mathrm{\Lambda }}{3}E^IE^JE^K\right]+_{}(\mathrm{\Omega },E).$$
(27)
The fundamental variables are a one-form connection $`\mathrm{\Omega }`$ and a one-form triad field $`E`$. The exact nature of the boundary term is not relevant to us here and we refer the reader to for more details. The indices $`I,J,K,\mathrm{}=0,1,2`$ label an internal space whose flat metric we denote $`h_{IJ}`$ and $`ϵ_{012}=1`$. The spacetime metric is as usual given by $`G_{\mu \nu }=h_{IJ}E_\mu ^IE_\nu ^J`$. We take $`h_{IJ}`$ to be of Lorentzian signature. The curvature two-form is $`F^{IJ}\left(\mathrm{\Omega }\right)=\mathrm{d}\mathrm{\Omega }^{IJ}+\mathrm{\Omega }_K^I\mathrm{\Omega }^{KJ}`$. It is convenient to identify the internal space indices with those of a three dimensional Lie algebra. This is $`SO(2,1)`$ for Lorentzian gravity and $`SO(3)`$ for a Euclidean spacetime. The internal metric $`h_{IJ}`$ is then taken to be proportional to the Killing-Cartan metric of this Lie algebra while $`ϵ_{JK}^I=h^{IL}ϵ_{LJK}`$ are its structure constants. Furthermore, for the connection to take value in the Lie algebra, we introduce the new connection (labelled with one index) through the redefinition $`\mathrm{\Omega }_I^J=ϵ_{IK}^J\mathrm{\Omega }^K`$. The one-form connection is such that $`\mathrm{\Omega }^{IJ}=\mathrm{\Omega }^{JI}`$ and is, therefore, metric-preserving (it satisfies the metricity condition).
We return now to our two examples in (22) and (25) and try to find their connection to three dimensional gravity. We specialise to the case when $`\eta _{ij}`$ and $`f_{jk}^i`$ describe an $`SO(2,1)`$ Lie algebra and are proportional to $`h_{IJ}`$ and $`ϵ_{JK}^I`$ respectively. We start by examining the first example in (22). As noticed above, the action $`S_{\mathrm{first}}^{(1)}`$ is effectively a BF theory in the bulk and can therefore be identified with $`S_{\mathrm{Palatini}}^{\mathrm{EH}}`$ with zero cosmological constant. The boundary Lagrangian $`(\mathrm{\Omega },E)`$ is unique and corresponds to that of $`S_{\mathrm{first}}^{(1)}`$. The rôle of the triads $`E_\mu ^I`$ is played by the field $`v_\mu ^i`$ and the three dimensional metric is then given by $`G_{\mu \nu }=\eta _{ij}v_\mu ^iv_\nu ^j`$. On the other hand, the connection $`\mathrm{\Omega }_\mu ^I`$ is related to $`O_\mu ^i\left(\theta \right)`$; the solution to the constraints imposed by $`e_\mu ^i`$ on the fields of the action (22). There is, however, a crucial difference between the two actions if one interprets $`\omega _\mu ^i`$ (and not $`O_\mu ^i\left(\theta \right)`$) as the spin connection of three dimensional gravity: The variation of (27) with respect to $`\mathrm{\Omega }_\mu ^I`$ implies a torsion-free condition whereas a variation of (22) with respect to $`\omega _\mu ^i`$ does not. This is due to the presence of the last two terms in the bulk part of the action (22). In addition, there are more fields in (22) than in (27) and the two theories coincide only if all the Lagrange multipliers in (8) are set to zero. This amounts to setting $`\theta _\mu ^i=\frac{1}{2}\omega _\mu ^i`$ and $`v_\mu ^i=\frac{1}{2}e_\mu ^i`$ in the action (22).
The situation with the action $`S_{\mathrm{first}}^{(2)}`$ is much more complex. If the Lagrange multipliers were absent then our action in (25) reduces to $`S_{\mathrm{palatini}}^{\mathrm{EH}}`$ in the presence of a non vanishing cosmological constant. This can be seen by setting $`\theta _\mu ^i=\frac{1}{2}\alpha e_\mu ^i`$ and $`v_\mu ^i=\frac{1}{2}\alpha \omega _\mu ^i`$ together with performing an integration by parts in the action (25). In general, however, the two actions are different. This is mainly due to the fact that the field $`e_\mu ^i`$ is cubic in the action (25) and cannot be integrated out. This makes it difficult to determine the fields that play the rôle of the triads and the connection of three dimensional gravity with a non vanishing cosmological constant.
## 4. Conclusions
The WZWN model is a conformal field theory that has been used in the past to reproduce various two dimensional theories like Toda field theories, black holes and others. In this paper we have enlarged this liste by connecting this model to a combination of three dimensional topological BF and Chern-Simons gauge theories defined on a manifold with boundaries. In order for this connection to hold, the boundary action accompanying these topological theories is unique. We arrive to this result by a direct application of non Abelian duality on the WZWN non-linear sigma model. One of our motivations in this work is to provide, at the classical level and at the level of the Lagrangian, a precise relashionship between the WZWN model and three dimensional gravity. This is a special case in our study an two examples are supplied. In the case of the Lie algebra $`SO(2,1)`$, we find that one obtains three dimensional gravity without a cosmological constant. However, the triads and the connection of gravity are not simply given by the components of the Chern-Simons gauge field. The other example is based on the Lie algebra $`SO(2,1)\times SO(2,1)`$ and yields three dimensional gravity with a cosmological constant only if the BF theory is not taken into account.
As we have stressed before, the integration over the Lagrange multipliers at any time would yield the original WZWN model. Notice that in both of our examples, the equations of motion of the Lagrange multipliers are simply Einstein’s equations (with and without cosmological constant) and the torsion-free condition, if one interprest $`e_\mu ^i`$ and $`\omega _\mu ^i`$, respectively, as the triads and connection of tree dimensional gravity. Therefore, it is only on-shell (that is, when Einstein’s equations and the torsion-free condition are satisfied) that three-dimensional gravity, with our unique boundary action, is equivalent to the WZWN theory. This statement is, of course, equivalent to ignoring the BF contribution.
One of the interesting problems would be of course to quantise the resulting BF and Cherns-Simons theories taking into account the boundary terms. This might provide a map between the observables of the two dual theories, namely the WZWN model and the topological theory. In the absence of boundaries, a quantum treatment of a combined BF and Chern-Simons has been presented in . There also, the question of which fields might play the rôle of the triads and connection of gravity is raised. We should also mention that we are still at the level the first order action. The next step in the completion of the dualisation procedure is to perform the integration over the gauge field. These appear in a cubic form in the Lagrangian and no known methods allow their complete integration from the path integral. However, a perturbative analyses is always possible. Notice though that there are situations, like the example in (22), where the gauge components appear at most in a quadratic form and without any derivatives acting on them. Integrating them out from the path integral is, in principle, feasible. It might be instructive to start by the quantisation of such models.
Aknowledgements: I would like to thank J. Balog and P. Forgács for discussions. |
warning/0001/math0001124.html | ar5iv | text | # An inequality for the norm of a polynomial factor
## 1. Introduction
Let $`p(z)`$ be a monic polynomial of degree $`n`$, with complex coefficients. Suppose that $`p(z)`$ has a monic factor $`q(z)`$, so that
$$p(z)=q(z)r(z),$$
where $`r(z)`$ is also a monic polynomial. Define the uniform (sup) norm on a compact set $`E`$ in the complex plane $``$ by
(1.1)
$$f_E:=\underset{zE}{sup}|f(z)|.$$
We study the inequalities of the following form
(1.2)
$$q_EC^np_E,\mathrm{deg}p=n,$$
where the main problem is to find the best (the smallest) constant $`C_E`$, such that (1.2) is valid for any monic polynomial $`p(z)`$ and any monic factor $`q(z)`$.
In the case $`E=\overline{D}`$, where $`D:=\{z:|z|<1\},`$ the inequality (1.2) was considered in a series of papers by Mignotte , Granville and Glesser , who obtained a number of improvements on the upper bound for $`C_{\overline{D}}`$. D. W. Boyd made the final step here, by proving that
(1.3)
$$q_{\overline{D}}\beta ^np_{\overline{D}},$$
with
(1.4)
$$\beta :=\mathrm{exp}\left(\frac{1}{\pi }_0^{2\pi /3}\mathrm{log}\left(2\mathrm{cos}\frac{t}{2}\right)𝑑t\right).$$
The constant $`\beta =C_{\overline{D}}`$ is asymptotically sharp, as $`n\mathrm{},`$ and it can also be expressed in a different way, using Mahler’s measure. This problem is of importance in designing algorithms for factoring polynomials with integer coefficients over integers. We refer to and for more information on the connection with symbolic computations.
A further development related to (1.2) for $`E=[a,a],a>0,`$ was suggested by P. B. Borwein in (see Theorems 2 and 5 there or see Section 5.3 in ). In particular, Borwein proved that if $`\mathrm{deg}q=m`$ then
(1.5)
$$|q(a)|p_{[a,a]}a^{mn}2^{n1}\underset{k=1}{\overset{m}{}}\left(1+\mathrm{cos}\frac{2k1}{2n}\pi \right),$$
where the bound is attained for a monic Chebyshev polynomial of degree $`n`$ on $`[a,a]`$ and a factor $`q`$. He also showed that, for $`E=[2,2]`$, the constant in the above inequality satisfies
$`\underset{n\mathrm{}}{lim\; sup}\left(2^{m1}{\displaystyle \underset{k=1}{\overset{m}{}}}\left(1+\mathrm{cos}{\displaystyle \frac{2k1}{2n}}\pi \right)\right)^{1/n}`$
$``$ $`\underset{n\mathrm{}}{lim}\left(2^{[2n/3]1}{\displaystyle \underset{k=1}{\overset{[2n/3]}{}}}\left(1+\mathrm{cos}{\displaystyle \frac{2k1}{2n}}\pi \right)\right)^{1/n}`$
$`=`$ $`\mathrm{exp}\left({\displaystyle _0^{2/3}}\mathrm{log}\left(2+2\mathrm{cos}\pi x\right)𝑑x\right)=1.9081\mathrm{},`$
which hints that
(1.6)
$$C_{[2,2]}=\mathrm{exp}\left(_0^{2/3}\mathrm{log}\left(2+2\mathrm{cos}\pi x\right)𝑑x\right)=1.9081\mathrm{}.$$
We find the asymptotically best constant $`C_E`$ in (1.2) for a rather arbitrary compact set $`E`$. The general result is then applied to the cases of a disk and a line segment, so that we recover (1.3)-(1.4) and confirm (1.6).
## 2. Results
Our solution of the above problem is based on certain ideas from the logarithmic potential theory (cf. or ). Let $`\mathrm{cap}(E)`$ be the logarithmic capacity of a compact set $`E`$. For $`E`$ with $`\mathrm{cap}(E)>0`$, denote the equilibrium measure of $`E`$ (in the sense of the logarithmic potential theory) by $`\mu _E`$. We remark that $`\mu _E`$ is a positive unit Borel measure supported on $`E`$, $`\mathrm{supp}\mu _EE`$ (see \[13, p. 55\]).
###### Theorem 2.1.
Let $`E`$ be a compact set, $`\mathrm{cap}(E)>0`$. Then the best constant $`C_E`$ in (1.2) is given by
(2.1)
$$C_E=\frac{\underset{uE}{\mathrm{max}}\mathrm{exp}\left({\displaystyle _{|zu|1}}\mathrm{log}|zu|d\mu _E(z)\right)}{\mathrm{cap}(E)}.$$
Furthermore, if $`E`$ is regular then
(2.2)
$$C_E=\underset{uE}{\mathrm{max}}\mathrm{exp}\left(_{|zu|1}\mathrm{log}|zu|d\mu _E(z)\right).$$
The above notion of regularity is to be understood in the sense of the exterior Dirichlet problem (cf. \[13, p. 7\]). Note that the condition $`\mathrm{cap}(E)>0`$ is usually satisfied for all applications, as it only fails for very thin sets (see \[13, pp. 63-66\]), e.g., finite sets in the plane. But if $`E`$ consists of finitely many points then the inequality (1.2) cannot be true for a polynomial $`p(z)`$ with zeros at every point of $`E`$ and for its linear factors $`q(z)`$. On the other hand, Theorem 2.1 is applicable to any compact set with a connected component consisting of more than one point (cf. \[13, p. 56\]).
One can readily see from (1.2) or (2.1) that the best constant $`C_E`$ is invariant under the rigid motions of the set $`E`$ in the plane. Therefore we consider applications of Theorem 2.1 to the family of disks $`D_r:=\{z:|z|<r\}`$, which are centered at the origin, and to the family of segments $`[a,a],a>0.`$
###### Corollary 2.2.
Let $`D_r`$ be a disk of radius $`r`$. Then the best constant $`C_{\overline{D}_r}`$, for $`E=\overline{D_r}`$, is given by
(2.3)
$$C_{\overline{D}_r}=\{\begin{array}{c}\frac{1}{r},0<r1/2,\hfill \\ \\ \frac{1}{r}\mathrm{exp}\left(\frac{1}{\pi }_0^{\pi 2\mathrm{arcsin}\frac{1}{2r}}\mathrm{log}\left(2r\mathrm{cos}\frac{x}{2}\right)𝑑x\right),r>1/2.\hfill \end{array}$$
Note that (1.3)-(1.4) immediately follow from (2.3) for $`r=1.`$ The graph of $`C_{\overline{D}_r}`$, as a function of $`r`$, is in Figure 1.
###### Corollary 2.3.
If $`E=[a,a],a>0`$, then
(2.4)
$$C_{[a,a]}=\{\begin{array}{c}\frac{2}{a},0<a1/2,\hfill \\ \\ \frac{2}{a}\mathrm{exp}\left(_{1a}^a\frac{\mathrm{log}(t+a)}{\pi \sqrt{a^2t^2}}𝑑t\right),a>1/2.\hfill \end{array}$$
Observe that (2.4), with $`a=2`$, implies (1.6) by the change of variable $`t=2\mathrm{cos}\pi x.`$ We include the graph of $`C_{[a,a]}`$, as a function of $`a`$, in Figure 2.
We now state two general consequences of Theorem 2.1. They explain some interesting features of $`C_E`$, which the reader may have noticed in Corollaries 2.2 and 2.3. Let
$$\mathrm{diam}(E):=\underset{z,\zeta E}{\mathrm{max}}|z\zeta |$$
be the Euclidean diameter of $`E`$.
###### Corollary 2.4.
Suppose that $`\mathrm{cap}(E)>0.`$ If $`\mathrm{diam}(E)1`$ then
(2.5)
$$C_E=\frac{1}{\mathrm{cap}(E)}.$$
It is well known that cap$`(D_r)=r`$ and cap$`([a,a])=a/2`$ (see \[12, p. 135\]), which clarifies the first lines of (2.3) and (2.4) by (2.5).
The next Corollary shows how the constant $`C_E`$ behaves under dilations of the set $`E`$. Let $`\alpha E`$ be the dilation of $`E`$ with a factor $`\alpha >0.`$
###### Corollary 2.5.
If $`E`$ is regular then
(2.6)
$$\underset{\alpha +\mathrm{}}{lim}C_{\alpha E}=1.$$
Thus Figures 1 and 2 clearly illustrate (2.6).
We remark that one can deduce inequalities of the type (1.2), for various $`L_p`$ norms, from Theorem 2.1, by using relations between $`L_p`$ and $`L_{\mathrm{}}`$ norms of polynomials on $`E`$ (see, e.g., ).
## 3. Proofs
###### Proof of Theorem 2.1.
The proof of this result is based on the ideas of and . For $`u`$, consider a function
$$\rho _u(z):=\mathrm{max}(|zu|,1),z.$$
One can immediately see that $`\mathrm{log}\rho _u(z)`$ is a subharmonic function in $`z`$, which has the following integral representation (see \[12, p. 29\]):
(3.1)
$$\mathrm{log}\rho _u(z)=\mathrm{log}|zt|d\lambda _u(t),z,$$
where $`d\lambda _u(u+e^{i\theta })=d\theta /(2\pi )`$ is the normalized angular measure on $`|tu|=1`$.
Let $`uE`$ be such that
$$q_E=|q(u)|.$$
If $`z_k,k=1,\mathrm{},m,`$ are the zeros of $`q(z)`$, counted according to multiplicities, then
(3.2) $`\mathrm{log}q_E`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{m}{}}}\mathrm{log}|uz_k|{\displaystyle \underset{k=1}{\overset{m}{}}}\mathrm{log}\rho _u(z_k)`$
$`=`$ $`{\displaystyle \underset{k=1}{\overset{m}{}}}{\displaystyle \mathrm{log}|z_kt|d\lambda _u(t)}={\displaystyle \mathrm{log}|p(t)|d\lambda _u(t)},`$
by (3.1).
We use the well known Bernstein-Walsh lemma about the growth of a polynomial outside of the set $`E`$ (see \[12, p. 156\], for example):
Let $`E`$ be a compact set, $`\mathrm{cap}(E)>0`$, with the unbounded component of $`\overline{}E`$ denoted by $`\mathrm{\Omega }`$. Then, for any polynomial $`p(z)`$ of degree $`n`$, we have
(3.3)
$$|p(z)|p_Ee^{ng_\mathrm{\Omega }(z,\mathrm{})},z,$$
where $`g_\mathrm{\Omega }(z,\mathrm{})`$ is the Green function of $`\mathrm{\Omega }`$, with pole at $`\mathrm{}`$. The following representation for $`g_\mathrm{\Omega }(z,\mathrm{})`$ is found in Theorem III.37 of \[13, p. 82\]).
(3.4)
$$g_\mathrm{\Omega }(z,\mathrm{})=\mathrm{log}\frac{1}{\mathrm{cap}(E)}+\mathrm{log}|zt|d\mu _E(t),z.$$
It follows from (3.1)-(3.4) and Fubini’s theorem that
$`{\displaystyle \frac{1}{n}}\mathrm{log}{\displaystyle \frac{q_E}{p_E}}`$ $``$ $`{\displaystyle \mathrm{log}\frac{|p(t)|^{1/n}}{p_E^{1/n}}d\lambda _u(t)}{\displaystyle g_\mathrm{\Omega }(t,\mathrm{})𝑑\lambda _u(t)}`$
$`=`$ $`\mathrm{log}{\displaystyle \frac{1}{\mathrm{cap}(E)}}+{\displaystyle \mathrm{log}|zt|d\lambda _u(t)𝑑\mu _E(z)}`$
$`=`$ $`\mathrm{log}{\displaystyle \frac{1}{\mathrm{cap}(E)}}+{\displaystyle \mathrm{log}\rho _u(z)𝑑\mu _E(z)}.`$
Using the definition of $`\rho _u(z)`$, we obtain from the above estimate that
$`q_E`$ $``$ $`\left({\displaystyle \frac{\underset{uE}{\mathrm{max}}\mathrm{exp}\left({\displaystyle \mathrm{log}\rho _u(z)𝑑\mu _E(z)}\right)}{\mathrm{cap}(E)}}\right)^np_E`$
$`=`$ $`\left({\displaystyle \frac{\underset{uE}{\mathrm{max}}\mathrm{exp}\left({\displaystyle _{|zu|1}}\mathrm{log}|zu|d\mu _E(z)\right)}{\mathrm{cap}(E)}}\right)^np_E.`$
Hence
(3.5)
$$C_E\frac{\underset{uE}{\mathrm{max}}\mathrm{exp}\left({\displaystyle _{|zu|1}}\mathrm{log}|zu|d\mu _E(z)\right)}{\mathrm{cap}(E)}.$$
In order to prove the inequality opposite to (3.5), we consider the $`n`$-th Fekete points $`\{a_{k,n}\}_{k=1}^n`$ for the set $`E`$ (cf. \[12, p. 152\]). Let
$$p_n(z):=\underset{k=1}{\overset{n}{}}(za_{k,n})$$
be the Fekete polynomial of degree $`n`$. Define the normalized counting measures on the Fekete points by
$$\tau _n:=\frac{1}{n}\underset{k=1}{\overset{n}{}}\delta _{a_{k,n}},n.$$
It is known that (see Theorems 5.5.4 and 5.5.2 in \[12, pp. 153-155\])
(3.6)
$$\underset{n\mathrm{}}{lim}p_n_E^{1/n}=\mathrm{cap}(E).$$
Furthermore, we have the following weak\* convergence of counting measures (cf. \[12, p. 159\]):
(3.7)
$$\tau _n\stackrel{}{}\mu _E,\text{ as }n\mathrm{}.$$
Let $`uE`$ be a point, where the maximum on the right hand side of (3.5) is attained. Define the factor $`q_n(z)`$ for $`p_n(z)`$, with zeros being the $`n`$-th Fekete points satisfying $`|a_{k,n}u|1`$. Then we have by (3.7) that
$`\underset{n\mathrm{}}{lim}q_n_E^{1/n}`$ $``$ $`\underset{n\mathrm{}}{lim}|q_n(u)|^{1/n}=\underset{n\mathrm{}}{lim}\mathrm{exp}\left({\displaystyle \frac{1}{n}}{\displaystyle \underset{|a_{k,n}u|1}{}}\mathrm{log}|ua_{k,n}|\right)`$
$`=`$ $`\mathrm{exp}\left(\underset{n\mathrm{}}{lim}{\displaystyle _{|zu|1}}\mathrm{log}|uz|d\tau _n(z)\right)`$
$`=`$ $`\mathrm{exp}\left({\displaystyle _{|zu|1}}\mathrm{log}|uz|d\mu _E(z)\right).`$
Combining the above inequality with (3.6) and the definition of $`C_E`$, we obtain that
$$C_E\underset{n\mathrm{}}{lim}\frac{q_n_E^{1/n}}{p_n_E^{1/n}}\frac{\mathrm{exp}\left({\displaystyle _{|zu|1}}\mathrm{log}|zu|d\mu _E(z)\right)}{\mathrm{cap}(E)}.$$
This shows that (2.1) holds true. Moreover, if $`uE`$ is a regular point for $`\mathrm{\Omega }`$, then we obtain by Theorem III.36 of \[13, p. 82\]) and (3.4) that
$$\mathrm{log}\frac{1}{\mathrm{cap}(E)}+\mathrm{log}|ut|d\mu _E(t)=g_\mathrm{\Omega }(u,\mathrm{})=0.$$
Hence
$$\mathrm{log}\frac{1}{\mathrm{cap}(E)}+_{|zu|1}\mathrm{log}|ut|d\mu _E(t)=_{|zu|1}\mathrm{log}|ut|d\mu _E(t),$$
which implies (2.2) by (2.1). ∎
###### Proof of Corollary 2.2.
It is well known \[13, p. 84\] that cap$`(\overline{D_r})=r`$ and $`d\mu _{\overline{D_r}}(re^{i\theta })=d\theta /(2\pi ),`$ where $`d\theta `$ is the angular measure on $`D_r`$. If $`r(0,1/2]`$ then the numerator of (2.1) is equal to 1, so that
$$C_{\overline{D_r}}=\frac{1}{r},0<r1/2.$$
Assume that $`r>1/2.`$ We set $`z=re^{i\theta }`$ and let $`u_0=re^{i\theta _0}`$ be a point where the maximum in (2.1) is attained. On writing
$$|zu_0|=2r\left|\mathrm{sin}\frac{\theta \theta _0}{2}\right|,$$
we obtain that
$`C_{\overline{D_r}}`$ $`=`$ $`{\displaystyle \frac{1}{r}}\mathrm{exp}\left({\displaystyle \frac{1}{2\pi }}{\displaystyle _{\theta _0+2\mathrm{arcsin}\frac{1}{2r}}^{2\pi +\theta _02\mathrm{arcsin}\frac{1}{2r}}}\mathrm{log}\left|2r\mathrm{sin}{\displaystyle \frac{\theta \theta _0}{2}}\right|d\theta \right)`$
$`=`$ $`{\displaystyle \frac{1}{r}}\mathrm{exp}\left({\displaystyle \frac{1}{2\pi }}{\displaystyle _{2\mathrm{arcsin}\frac{1}{2r}\pi }^{\pi 2\mathrm{arcsin}\frac{1}{2r}}}\mathrm{log}\left(2r\mathrm{cos}{\displaystyle \frac{x}{2}}\right)𝑑x\right)`$
$`=`$ $`{\displaystyle \frac{1}{r}}\mathrm{exp}\left({\displaystyle \frac{1}{\pi }}{\displaystyle _0^{\pi 2\mathrm{arcsin}\frac{1}{2r}}}\mathrm{log}\left(2r\mathrm{cos}{\displaystyle \frac{x}{2}}\right)𝑑x\right),`$
by the change of variable $`\theta \theta _0=\pi x.`$
###### Proof of Corollary 2.3.
Recall that cap$`([a,a])=a/2`$ (see \[13, p. 84\]) and
$$d\mu _{[a,a]}(t)=\frac{dt}{\pi \sqrt{a^2t^2}},t[a,a].$$
It follows from (2.1) that
(3.8)
$$C_{[a,a]}=\frac{2}{a}\mathrm{exp}\left(\underset{u[a,a]}{\mathrm{max}}_{[a,a](u1,u+1)}\frac{\mathrm{log}|tu|}{\pi \sqrt{a^2t^2}}𝑑t\right).$$
If $`a(0,1/2]`$ then the integral in (3.8) obviously vanishes, so that $`C_{[a,a]}=2/a`$. For $`a>1/2`$, let
(3.9)
$$f(u):=_{[a,a](u1,u+1)}\frac{\mathrm{log}|tu|}{\pi \sqrt{a^2t^2}}𝑑t.$$
One can easily see from (3.9) that
$$f^{}(u)=_{u+1}^a\frac{dt}{\pi (ut)\sqrt{a^2t^2}}<0,u[a,1a],$$
and
$$f^{}(u)=_a^{u1}\frac{dt}{\pi (ut)\sqrt{a^2t^2}}>0,u[a1,a].$$
However, if $`u(1a,a1)`$ then
$$f^{}(u)=_{u+1}^a\frac{dt}{\pi (ut)\sqrt{a^2t^2}}+_a^{u1}\frac{dt}{\pi (ut)\sqrt{a^2t^2}}.$$
It is not difficult to verify directly that
$$\frac{dt}{\pi (ut)\sqrt{a^2t^2}}=\frac{1}{\pi \sqrt{a^2u^2}}\mathrm{log}\left|\frac{a^2ut+\sqrt{a^2t^2}\sqrt{a^2u^2}}{tu}\right|+C,$$
which implies that
$$f^{}(u)=\frac{1}{\pi \sqrt{a^2u^2}}\mathrm{log}\left(\frac{a^2u^2+u+\sqrt{a^2(u1)^2}\sqrt{a^2u^2}}{a^2u^2u+\sqrt{a^2(u+1)^2}\sqrt{a^2u^2}}\right),$$
for $`u(1a,a1).`$ Hence
$$f^{}(u)<0,u(1a,0),\text{ and }f^{}(u)>0,u(0,a1).$$
Collecting all facts, we obtain that the maximum for $`f(u)`$ on $`[a,a]`$ is attained at the endpoints $`u=a`$ and $`u=a`$, and it is equal to
$$\underset{u[a,a]}{\mathrm{max}}f(u)=_{1a}^a\frac{\mathrm{log}(t+a)}{\pi \sqrt{a^2t^2}}𝑑t.$$
Thus (2.3) follows from (3.8) and the above equation. ∎
###### Proof of Corollary 2.4.
Note that the numerator of (2.1) is equal to 1, because $`|zu|1,zE,uE`$. Thus (2.5) follows immediately. ∎
###### Proof of Corollary 2.5.
Observe that $`C_E1`$ for any $`E`$, so that $`C_{\alpha E}1`$. Since $`E`$ is regular, we use the representation for $`C_E`$ in (2.2). Let $`T:E\alpha E`$ be the dilation mapping. Then $`|TzTu|=\alpha |zu|,z,uE,`$ and $`d\mu _{\alpha E}(Tz)=d\mu _E(z)`$. This gives that
$`C_{\alpha E}`$ $`=`$ $`\underset{Tu(\alpha E)}{\mathrm{max}}\mathrm{exp}\left({\displaystyle _{|TzTu|1}}\mathrm{log}|TzTu|d\mu _{\alpha E}(Tz)\right)`$
$`=`$ $`\underset{uE}{\mathrm{max}}\mathrm{exp}\left({\displaystyle _{|zu|1/\alpha }}\mathrm{log}(\alpha |zu|)𝑑\mu _E(z)\right)`$
$`=`$ $`\underset{uE}{\mathrm{max}}\mathrm{exp}\left(\mu _E(\overline{D_{1/\alpha }(u)})\mathrm{log}\alpha {\displaystyle _{|zu|1/\alpha }}\mathrm{log}|zu|d\mu _E(z)\right)`$
$`<`$ $`\underset{uE}{\mathrm{max}}\mathrm{exp}\left({\displaystyle _{|zu|1/\alpha }}\mathrm{log}|zu|d\mu _E(z)\right),`$
where $`\alpha 1`$. Using the absolute continuity of the integral, we have that
$$\underset{\alpha +\mathrm{}}{lim}_{|zu|1/\alpha }\mathrm{log}|zu|d\mu _E(z)=0,$$
which implies (2.6). ∎ |
warning/0001/astro-ph0001314.html | ar5iv | text | # Theory of “Jitter” Radiation from Small-Scale Random Magnetic Fields and Prompt Emission from Gamma-Ray Burst Shocks
## 1. Introduction
The conventional paradigm of the generation of radiation by relativistic electrons in magnetic fields is totally based on the theory of synchrotron radiation. We demonstrate that this theory is invalid if a magnetic field is tangled on very short spatial scales and we develop a quantitative theory of radiation in this case. Apparently, the required short-scale fields may naturally be present in astrophysical shocks. Here we focus on radiation from gamma-ray burst (GRB) shocks, for which a detailed theory of the formation of magnetic fields has recently been elaborated.
The relativistic blast-wave model of cosmological GRBs (see, e.g., a review by Piran 1999) explains fairly well many observational features of this phenomenon, such as the rapid variability of $`\gamma `$-ray flux, the prompt optical flash, the light curves and spectra of delayed afterglows, etc.. This model interprets the prompt $`\gamma `$-ray flash as synchrotron radiation emitted by Fermi-accelerated electrons in internal shocks propagating in the ejecta (Rees & Mészáros 1994) and then Lorentz-boosted to the $`\gamma `$-ray band. The afterglows are explained in a similar way, as the emission from an external shock (Mészáros & Rees 1993) propagating into the interstellar medium. To achieve the observed very high luminosities, the magnetic field in the GRB shocks must be of nearly equipartition strength, $`ϵ_B=B^2/8\pi e_t1`$, where $`e_t`$ is the thermal energy density of the shocked material.
Until very recently, the equipartition assumption was completely unjustified. Medvedev & Loeb (1999) have shown that the relativistic two-stream instability is capable of producing magnetic fields with $`ϵ_B10^110^5`$ in both internal and external shocks. Observations of afterglow spectra and light curves yield values of $`ϵ_B`$ from $`10^110^2`$ for GRB 970508 (Wijers & Galama 1998; Frail et al. 1999; Granot et al. 1999) to $`10^5`$ for GRB 990123 and GRB 971214 (Galama et al. 1999). Recent detection of polarization of the optical afterglow of GRB 990510 (Covino et al. 1999; Wijers et al. 1999) indicates that the geometry of the magnetic field in the shock is consistent with the predictions of Medvedev & Loeb (1999) for collimated outflows (Ghisellini & Lazzati 1999; Gruzinov 1999; Sari 1999).
The magnetic field produced in GRB shocks randomly fluctuates on a very small scale of roughly the relativistic skin depth, which is $`\lambda _B10^2`$ cm in internal shocks, for instance. On the other hand, the emitting ultrarelativistic electrons have much larger Larmor radii. Therefore, the electron trajectories are not helical, as they would be in a homogeneous field. Thus, the theory of synchrotron radiation derived for homogeneous fields is not applicable and the spectrum of the emergent radiation is different. Such a situation has never been considered in the astrophysical literature.
In this paper we investigate the effect of small-scale magnetic fields on the properties of radiation. We primarily focus on internal shocks, for concreteness. We show that there are two regimes, depending on the ratio of particle’s deflection angle and the relativistic beaming angle. Which regime is realized depends on the magnetic field properties, $`B`$ and $`\lambda _B`$, but is independent of the particle’s energy. When deflections are large compared to beaming, synchrotron radiation is emitted. Otherwise, when particle’s deflections are small, a new type of radiation — jitter radiation — is produced. A quantitative analytical theory of this radiation is developed in this paper. For the power-law distributed electrons with a cutoff, $`N(\gamma )\gamma ^p`$ for $`\gamma \gamma _{\mathrm{min}}`$, where $`\gamma `$ is the particle’s Lorentz factor and $`p`$ is the index, the emergent spectrum has the following properties. First, the spectral power peaks at the so-called jitter frequency, $`\omega _{jm}=\omega _j(\gamma _{\mathrm{min}})`$, which, unlike the synchrotron frequency, is independent of the magnetic field strength but, instead, depends on the particle density in the shock, $`\omega _{jm}\sqrt{n}`$. Second, at low frequencies, $`\omega \omega _{jm}`$, the spectral power scales as $`P(\omega )\omega ^1`$, in contrast to the synchrotron spectrum, for which $`P(\omega )\omega ^{1/3}`$. The high-frequency portion is, however, determined by the electrons, $`P(\omega )\omega ^{(p1)/2}`$. Third, the total (i.e., frequency integrated) powers emitted in the jitter and synchrotron regimes are identical.
Since large-scale fields (e.g., due to a magnetized progenitor, for instance) may also be present in the shocked material, we construct a composite, two-component jitter+synchrotron (JS) spectral model of the prompt $`\gamma `$-ray emission. We then compare the predictions of this model with presently available data collected (mostly) by the Burst And Transient Source Experiment (BATSE) on the Compton Gamma-Ray Observatory (CGRO). It turns out that the JS model is able to naturally explain some properties of the GRB spectra which are inconsistent with a simple synchrotron shock model. First, almost a half of all BATSE bursts violate the so-called synchrotron “line of death” prediction, i.e., their low-frequency spectral indices are greater than the maximal admissible value of $`1/3`$ (Preece et al. 1998). The spectra of these bursts are, however, well consistent with the steeper $`\omega ^1`$-law of jitter radiation. Second, the sharp change of the spectral index at a peak frequency seen in some bursts is also consistent with our model. Third, the JS model theoretically supports the fact that some GRBs have two spectral sub-components (Pendleton et al. 1994). Fourth, there is an indication that spectra of some GRBs exhibit spectral features — “GRB lines” see Briggs 1999 for a review and some BATSE candidates. It should be noted, however, that the results from Ginga and BATSE are somewhat controversial. We demonstrate that a line-like spectral feature may be associated with a weak jitter component in a synchrotron-dominated spectrum. If a future analysis shall confirm that “GRB lines” are real features and not instrumental (or other) artifacts, then they provide, together with the synchrotron component, precise information about properties of cosmological fireballs just a few hundred seconds after the explosion. We illustrate this on the example of GRB 910503 which has been observed by all four instruments on the CGRO and which exhibits a second spectral peak at roughly twice the synchrotron peak frequency. A simple fit of a spectral shape readily yields the value of the magnetic field, $`ϵ_B4\times 10^4`$ in the shock, which is in agreement with results of a completely different analysis by Chiang & Dermer (1999). It should be noted that a reliable identification/detection of the jitter spectral features will provide a direct evidence that the magnetic fields in GRBs are generated by the two-stream instability, since this is the only presently known mechanism capable of producing small-scale, large-amplitude fields in shocks.
At last, we emphasize that the advantage of our JS model is that it was not specially designed in order to explain peculiarities of the GRB spectra, but solely to study the physical effect of small-scale magnetic fields. The phenomenon considered in this paper is quite general and, clearly, relevant not only to the emission from internal GRB shocks. A similar mechanism of emission is expected to operate in external GRB shocks and, possibly, in more conventional supernova shocks and blazar jets.
The rest of this paper is organized as follows. A qualitative consideration of the radiation mechanisms in a nonuniform magnetic field is presented in §2. In §3 the structure and properties of magnetic fields in GRB internal shocks are discussed. In §4 we present a quantitative analytical theory of jitter radiation. A two-component, jitter+synchrotron spectral model of the prompt $`\gamma `$-ray emission is presented in §5. We compare the predictions of the model with recent observational results in §6. Finally, §7 is the conclusion.
## 2. Radiation from Small-Scale Fields: General Consideration
Let’s consider a plasma at rest threaded by a magnetic field. Let’s now consider an ultrarelativistic electron with the Lorentz factor $`\gamma 1`$ moving in a magnetic field. Because of beaming the emerging radiation is concentrated in a narrow cone with the opening angle $`\mathrm{\Delta }\theta 1/\gamma 1`$ in the direction of the particle’s velocity. In a uniform magnetic field the electron moves along a helical trajectory, so that the radiation seen by an observer consists of short pulses repeated every cyclotron period. The synchrotron spectrum, therefore, consists of a large number of cyclotron harmonics, the envelope of which is determined by the inverse duration of the pulses (Rybicki & Lightman 1979). The spectrum is peaked near the critical synchrotron frequency $`\omega _c=\frac{3}{2}\gamma ^2eB_{}/m_ec`$, where $`B_{}=B\mathrm{cos}\chi `$ and $`\chi `$ is the particle’s pitch angle.
If the magnetic field is randomly tangled and the correlation length is less then a Larmor radius of an emitting electron, then the electron experiences random deflections as it moves through the field. Its trajectory is, in general, stochastic. This is similar to a collisional motion of an electron in a medium. Bremsstrahlung quanta are emitted in every collision. Unlike the bremsstrahlung case, here “collisions” are due to small-scale inhomogeneities of the magnetic field rather than due to electrostatic fields of other charged particles. Since the Lorentz force depends on particle’s velocity, the emergent spectrum will be somewhat different from pure bremsstrahlung. There is also an alternative physical interpretation of the process. For an ultrarelativistic electron, the method of virtual quanta applies (Rybicki & Lightman 1979). In the rest frame of the electron, the magnetic field inhomogeneity with wavenumber $`k`$ is transformed into a transverse pulse of electromagnetic radiation with frequency $`kc`$. This radiation is then Compton scattered by the electron to produce observed radiation with frequency $`\gamma ^2kc`$ in the lab frame.
Keeping this general physical picture in mind, we now analyze the problem in more details. Let’s consider a nonuniform random magnetic field with a typical correlation scale $`\lambda _B`$, the Larmor radius of the electron, $`\rho _e=\gamma m_ec^2/eB_{}`$ is less or comparable comparable to $`\lambda _B`$. The emerging spectrum depends on the relation between the particle’s deflection angle, $`\alpha `$, and the beaming angle, $`\mathrm{\Delta }\theta `$, (Landau & Lifshits 1975). For ultrarelativistic particles and small deflection angles, the latter is estimated as follows. The particle’s momentum is $`p\gamma m_ec`$. The change in the perpendicular momentum due to the Lorentz force acting on the particle during the transit time $`t\lambda _B/c`$ is $`p_{}F_LteB_{}\lambda _B/c`$. The angle $`\alpha `$ is then $`\alpha p_{}/peB_{}\lambda _B/\gamma m_ec^2`$. Thus, the ratio of the two angles is
$$\frac{\alpha }{\mathrm{\Delta }\theta }\frac{eB_{}\lambda _B}{m_ec^2}\gamma \frac{\lambda _B}{\rho _e}.$$
(1)
It is interesting to note that this ratio is independent of particle’s energy (i.e., of $`\gamma `$) and is determined by the properties of the magnetic field only, i.e., by $`B`$ and $`\lambda _B`$. It is more convenient, however, to use the wave-vector, $`k_B`$, as a measure of the magnetic field scale, instead of $`\lambda _Bk_B^1`$. We now define the deflection-to-beaming ratio as follows,
$$\delta \frac{\gamma }{k_B\rho _e}\frac{\alpha }{\mathrm{\Delta }\theta }.$$
(2)
There are two limiting cases.
First, $`\delta \alpha /\mathrm{\Delta }\theta 1`$, i.e., the deflection angle is much larger than the beaming angle, see Figure 2a. Then, an observer sees radiation coming from short segments (“patches”) of the electron’s trajectory which are nearly parallel to the line of sight (very much like for pure synchrotron radiation). The magnetic field in every patch is almost uniform but it varies from patch to patch. The radiation is pulsed with a typical duration $`\tau _p1/\omega _c`$. The characteristic frequency of the observed radiation is thus $`\omega \omega _c`$. Note that in the time- or ensemble-averaged spectrum, the instantaneous field, $`B_{}`$, which enters $`\omega _c`$, should be appropriately averaged, $`\overline{B}\sqrt{B_{}^2}`$. In this case the emergent radiation is completely identical to synchrotron radiation from large-scale weakly inhomogeneous magnetic fields.
Second, $`\delta \alpha /\mathrm{\Delta }\theta 1`$, i.e., the deflection angle is smaller than the beaming angle, so that the entire electron’s trajectory is seen by an observer, as shown in Figure 2b. The particle moves along the line of sight almost straight and experiences high-frequency jittering in the perpendicular direction due to the random Lorentz force. We therefore refer the emerging radiation to as “jitter” radiation. Its spectrum is determined by random accelerations of the particle. Let’s imagine an electron moving ultrarelativistically along the line of sight with a constant velocity, the transverse accelerations of the electron are small. In the laboratory frame, the electron passes through the magnetic field inhomogeneities having a typical scale $`\lambda _Bk_B^1`$ with the velocity $`c`$. In the particle’s frame (i.e., where its parallel velocity vanishes), the field correlation scale is $`\lambda _B^{}\lambda _B/\gamma (k_B\gamma )^1`$ due to the Lorentz transformation. The electron’s perpendicular acceleration changes significantly during $`\tau \lambda _B^{}/c(\gamma k_Bc)^1`$, so that the characteristic frequency of the emitted radiation is $`\omega _{j0}\tau ^1\gamma k_Bc`$. In the laboratory frame, this frequency is boosted to $`\omega _j=\gamma \omega _{j0}`$. Thus, the spectrum of the emergent radiation is peaked at the frequency $`\omega _j\gamma ^2k_Bc`$. This frequency is higher than the critical synchrotron frequency in the uniform magnetic field of the same strength, $`\omega _c`$, namely
$$\frac{\omega _c}{\omega _j}\frac{3}{2}\delta 1,$$
(3)
as follows from equation (2).
We should warn here that despite apparent similarity of the jitter and free electron laser emission mechanisms, they are quite different. The wiggler field in free electron lasers is appropriately adjusted for the electron motion to be in phase with the produced radiation field to emit coherent radiation. Jitter radiation is, in general, incoherent.
## 3. The structure of the magnetic field in GRB shocks
To proceed further, a model for a magnetic field in GRB shocks is required. We use the only presently available quantitative theory of the magnetic field generation in shocks proposed by Medvedev & Loeb (1999). To be specific, we focus on internal shocks which produce $`\gamma `$-ray emission. External shocks which are responsible for the delayed afterglows may be treated similarly and will be considered in a future publication.
Shock fronts are shown to be natural sites of the magnetic field generation. Right before a shock, the inflowing (in the shock frame) bulk plasma particles meet the outflowing particles which were reflected (scattered) from the shock. Such a two-stream motion is kinetically unstable. The emergent magnetic field is random with zero mean and lies in the plane of the shock front, — perpendicular to the shock velocity. In principle, all plasma species participate in the instability. We assume the protons and electrons to be the only species and discuss their contributions separately.
It is important to emphasize that the generated magnetic field fills the entire volume of a shock shell and is not located within a thin layer of order several skin depths near the front. There is a gas flow though a shock. Because of flux freezing the generated magnetic field is transported with the shocked material downstream. This material is replenished with a fresh one where a new magnetic filed is thus continuously produced. Since the magnetic field is long-lived and does not decay in a dynamical time, as indicated by numerical simulations see references in Medvedev & Loeb 1999, this field will be present in the entire ejecta.
### 3.1. Fields produced by the electrons and protons
In this subsection we briefly remind main results of the theory of Medvedev & Loeb (1999) for future reference. Since electrons are light, the instability induced by them is rapid: the typical $`e`$-folding length (i.e., the $`e`$-folding time times the shock speed) is much smaller than the characteristic shock thickness determined by the Larmor radius of heavier protons. Therefore, the magnetic field energy grows rapidly and reaches the approximate equipartition with the electron kinetic energy,
$$\frac{\overline{B}_e^2}{8\pi }=\eta _e\gamma _{\mathrm{int}}m_ec^2n=\frac{m_e}{m_p}e_t\eta _e,$$
(4)
i.e., $`ϵ_{Be}=\overline{B}_e^2/8\pi e_t=(m_e/m_p)\eta _e`$, where $`\eta _e`$ is the efficiency factor for the electrons which incorporates uncertainties due to the nonlinear phase of the instability, one infers from numerical simulations that generically $`\eta _e0.10.01`$, $`\gamma _{\mathrm{int}}`$ is the relative Lorentz factor of two colliding shells which produce an internal shock, $`\gamma _{\mathrm{int}}24`$, and $`n`$ is the number density of particles in the expanding shell, before the shock. The spatial correlation scale, $`k_{Be}`$, of the field behind the shock is
$$k_{Be}=\frac{4\gamma _{\mathrm{int}}\omega _{pe}}{2^{1/4}\overline{\gamma }_e^{1/2}c},$$
(5)
where $`\omega _{pe}^2=4\pi e^2n/m_e`$ is the electron plasma frequency squared, $`\overline{\gamma }_e`$ is the initial effective thermal Lorentz factor of the streaming electrons, and an extra factor of $`4\gamma _{\mathrm{int}}`$ is due to the relativistic shock compression.
The protons may generate magnetic fields too. Since they are much heavier than electrons, the spatial coherence length and the $`e`$-folding length are comparable to the thickness of the collisionless shock. Therefore, the field does not have enough time to grow during the flow transit through the shock. It may however grow behind the shock if the two-stream motion of the protons persists there. If this is the case, then there are two possibilities. First, if there is no energy transfer from the protons to the electrons or if it is slow compared to the rate of the field growth, then the magnetic field energy may be as large as
$$\frac{\overline{B}_p^2}{8\pi }=\eta _p\gamma _{\mathrm{int}}m_pc^2n=e_t\eta _p0.1e_t,$$
(6)
provided $`\eta _p0.1`$. Second, if the energy transfer is fast, which may be the case in fields which are in equipartition with the electrons or stronger, then the protons may efficiently damp their energy into the emitting electrons, so that the resultant field will be $`\overline{B}_p^2/8\pi \overline{B}_e^2/8\pi `$. Alternatively, if no magnetic field is generated downstream the shock, the strength of the field produced by the protons may be orders of magnitude lower. Which case realizes may be learned from numerical particle simulations or from observations. We keep $`ϵ_{Bp}=\overline{B}_p^2/8\pi e_t`$ as a parameter. The characteristic correlation scale of the generated magnetic field is
$$k_{Bp}=\frac{\omega _{pp}}{2^{1/4}\overline{\gamma }_p^{1/2}c},$$
(7)
where $`\omega _{pp}^2=4\pi e^2n/m_p`$, $`\overline{\gamma }_p2`$ is the initial effective thermal Lorentz factor of the streaming protons. The compression factor, $`4\gamma _{\mathrm{int}}`$, is absent because the field is likely produced after the compression occurs.
### 3.2. The model
From equations (2), (4), and (5), we estimate the $`\delta `$-parameter for the electrons,
$$\delta _e=\frac{1}{2^{7/4}}\left(\frac{\overline{\gamma }_e}{\gamma _{\mathrm{int}}}\right)^{1/2}\sqrt{\eta _e}\varphi \sqrt{\eta _e}1.$$
(8)
The exact value of $`\varphi =2^{7/4}(\overline{\gamma }_e/\gamma _{\mathrm{int}})^{1/2}`$ is somewhat uncertain: $`\overline{\gamma }_e`$ may evolve during the instability from its initial value $`\overline{\gamma }_e23`$ to $`\overline{\gamma }_e\gamma _{\mathrm{int}}34`$ due to nonlinear effects.<sup>1</sup><sup>1</sup>1Note, the inflowing electrons are cold; they are not Fermi accelerated yet. Note also that $`\overline{\gamma }_e`$ cannot be greater than $`\gamma _{\mathrm{int}}`$ for the instability to operate. Otherwise, no magnetic field is produced. The numerical prefactor may also be affected. Therefore we assume that possible values of $`\delta _e`$ are in the range $`1\delta _e10^2`$ (given the uncertainty in $`\eta _e`$ from 0.1 to 0.01) and generically $`\delta _e0.1`$. The $`\delta `$-parameter for the protons is
$$\delta _p=2^{1/4}\left(\overline{\gamma }_p\gamma _{\mathrm{int}}\right)^{1/2}\frac{m_p}{m_e}\sqrt{\eta _p}1,$$
(9)
unless $`\eta _p`$ is too small: $`\eta _p10^7`$ (i.e., $`ϵ_{Bp}10^3ϵ_{Be}`$).
As will be shown below in §4, a spatial spectrum of the magnetic field $`\overline{B}_e`$ with $`\delta _e<1`$ is required to calculate the spectrum of jitter radiation. This distribution of $`\overline{B}_e`$ over scales is difficult to find from the first principles because it is determined by fully nonlinear dynamics of the instability process. Some constraints may however be drawn. First, the two-stream instability produces magnetic fields in a finite range of scales, $`0kk_{\mathrm{crit},e}`$, where $`k_{\mathrm{crit},e}k_{Be}`$ to within a numerical factor of order unity. Second, the field grows until it becomes strong enough to deflect the particles in the transverse direction by $`1`$ radian on a scale of the field coherence length, i.e., $`\overline{B}_e^1\rho _ek_{Be}^1`$. If we now assume that each Fourier harmonic, $`B_k`$, is amplified independently of others, we obtain: $`B_kk`$ for $`kk_{\mathrm{crit},e}`$. This is the lower limit: the spectrum of the magnetic field can only be steeper than linear.
In reality, all Fourier harmonics are coupled. Thus, when at least one harmonic reaches the sub-equipartition strength, the streaming electrons are isotropized by random Lorentz forces. This prevents the growth of other harmonics. Therefore, the spectrum of the field will be steeper than linear and will have a maximum near $`k_{Be}`$. The simplest possible model is a power-law,
$$B_k=\{\begin{array}{cc}C_Bk^\mu ,\hfill & \text{ if }0kk_{Be},\hfill \\ 0,\hfill & \text{ otherwise},\hfill \end{array}$$
(10)
where $`C_B`$ is a normalization constant and $`\mu 1`$ is a power-law spectral index of the magnetic field being a free parameter here. The constant $`C_B`$ may be determined using Parseval’s theorem,<sup>2</sup><sup>2</sup>2We use the following definition of the Fourier transform: $`f(\omega )=f(t)e^{i\omega t}𝑑t`$. $`lim_L\mathrm{}_{L/2}^{L/2}B_e^2𝑑x=\frac{1}{\pi }_0^{\mathrm{}}B_k^2𝑑k`$, where $`L`$ is the system size. Taking into account that $`B_e^2𝑑x=\overline{B}_e^2L=\overline{B}_e^2cT`$, where $`T`$ is the total duration of the pulse, we write
$$C_B^2=\pi \left(2\mu +1\right)cT\overline{B}_e^2k_{Be}^{(2\mu +1)}.$$
(11)
The magnitude of the field, $`\overline{B}_e`$, is given by equation (4).
Strictly speaking, this model may be used for not too small $`k`$’s: $`k>k_{\mathrm{min}}`$, where $`k_{\mathrm{min}}`$ is set by the condition: $`\gamma /(\rho _ek_{\mathrm{min}})=1`$, see equation (2). The field harmonics with $`k<k_{\mathrm{min}}`$ are large-scale ones and they contribute to synchrotron radiation. Using equation (2), it is easy to obtain that $`k_{\mathrm{min}}=\delta _ek_{Be}`$. It is also useful to introduce the small-scale field component as follows,
$$\overline{B}_{SS}^2=_{k_{\mathrm{min}}}^{k_{Be}}B_k^2𝑑k,\frac{\overline{B}_{SS}^2}{\overline{B}_e^2}=1\delta ^{2\mu +1}.$$
(12)
Hereafter, we omit the subscripts by $`\delta `$; this should not cause any confusion.
To summarize, we assume the following structure of the magnetic field sketched in Figure 3.2. First, there is a magnetic field produced by the electrons, $`\overline{B}_e`$, the magnitude of which is given by equation (4). A large fraction of this field contributes to the small-scale component, in accordance with equation (12). Its distribution over spatial scales is described by a power-law with the index $`\mu `$, see equation (10). The rest, $`\overline{B}_e^2\overline{B}_{SS}^2`$, contributes to the large-scale component. Second, there is a magnetic field produced by the protons, $`\overline{B}_p`$. This field is a large-scale field. Both $`\overline{B}_e`$ and $`\overline{B}_p`$ are random with zero mean. Third, there could be an ordered magnetic field left from a magnetized progenitor, $`B_{}`$. We define the total large-scale magnetic field as follows,
$$\overline{B}_{LS}^2=\overline{B}_0^2+\left(\overline{B}_e^2\overline{B}_{SS}^2\right),\overline{B}_0^2=B_{}^2+\overline{B}_p^2,$$
(13)
where $`\overline{B}_0`$ is the fraction of the large-scale field which is not produced by the electrons.
## 4. Radiation from Small-Scale Fields: Quantitative Theory
In §2 we qualitatively demonstrated that there are two regimes of radiation. An ultrarelativistic electron propagating through small-scale fields generated only by the electrons ($`\delta 1`$) emits jitter radiation. The electron moving through larger-scale fields generated only by the protons ($`\delta 1`$) emits synchrotron radiation. The radiation spectrum from a magnetic field with a broadband distribution over scales is neither of the above and must be calculated by appropriate scale averaging of a particle trajectory. However, for a bimodal field distribution of §3.2 such that the magnetic energy in $`B_k`$ harmonics for which $`\delta 1`$ (i.e., where the transition from a jitter to synchrotron regime occurs) is small, the separation of scales is possible. The resultant radiation will approximately be a composition of the jitter and synchrotron spectra.<sup>3</sup><sup>3</sup>3Physically, a particle moves along a helical trajectory about field lines of a large-scale field. This trajectory is slightly perturbed by high frequency jittering in the (instantaneous) transverse direction due to a small-scale field component. The intensities of the spectral sub-components are determined by the magnetic field energy densities at large and small scales. For clarity and illustration purposes, we derive the radiation properties and spectra in both limits separately. A rigorous treatment of a general case will be given elsewhere. All calculations in this section are done in the reference frame of the shocked material, unless stated otherwise.
### 4.1. Jitter radiation, $`\delta 1`$
For sufficiently small $`\delta `$’s such that $`\alpha \mathrm{\Delta }\theta `$, the velocity $`𝐯`$ of a particle is almost constant whereas its acceleration $`𝐰\dot{𝐯}`$ varies with time. Calculating the Fourier component of the electric field using the Liénard-Wiechart (retarded) potentials, one arrives at the following expression for the total energy emitted per unit solid angle $`d\mathrm{\Omega }`$ per unit frequency $`d\omega `$ Landau & Lifshits 1975, §77; see also Rybicki & Lightman 1979, §3.2:
$$dW=\frac{e^2}{2\pi c^3}\left(\frac{\omega }{\omega ^{}}\right)^4\left|𝐧\times \left[\left(𝐧\frac{𝐯}{c}\right)\times 𝐰_\omega ^{}\right]\right|^2d\mathrm{\Omega }\frac{d\omega }{2\pi },$$
(14)
where $`𝐰_\omega ^{}=𝐰e^{i\omega ^{}t}𝑑t`$ is the Fourier component of the particle’s acceleration, $`\omega ^{}=\omega \left(1𝐧𝐯/c\right)`$, and $`𝐧`$ is the unit vector pointing towards the observer.
First, we can simplify the vector expression in (14). Indeed, in the ultrarelativistic case, the longitudinal component of the acceleration is small compared to the transverse component, $`w_{}/w_{}1/\gamma ^21`$. Therefore $`𝐯`$ and $`𝐰`$ are approximately perpendicular to each other. Second, the dominant contribution to the integral over $`d\mathrm{\Omega }`$ comes from small angles $`\theta 1/\gamma `$ with respect to the particle’s velocity. Therefore, we approximately write $`\omega ^{}\omega \left(1v/c+\theta ^2/2\right)\frac{1}{2}\omega \left(1v^2/c^2+\theta ^2\right)=\frac{1}{2}\omega \left(\theta ^2+\gamma ^2\right)`$. We now can replace integration over the solid angle $`d\mathrm{\Omega }\theta d\theta d\varphi `$ with integration over $`d\varphi d\omega ^{}/\omega `$ and integrate equation (14) over the azimuthal angle, $`\varphi `$, from $`0`$ to $`2\pi `$. The spectral energy finally becomes
$$\frac{dW}{d\omega }=\frac{e^2\omega }{2\pi c^3}_{\omega /2\gamma ^2}^{\mathrm{}}\frac{\left|𝐰_\omega ^{}\right|^2}{\omega ^2}\left(1\frac{\omega }{\omega ^{}\gamma ^2}+\frac{\omega ^2}{2\omega ^2\gamma ^4}\right)𝑑\omega ^{}.$$
(15)
The leading term in the above equation is due to high-frequency “jittering” of the electron as it moves through the random magnetic field. The second and third terms in the brackets are corrections due to the angular structure of the radiation field convolved with the relativistic beaming.
The acceleration $`𝐰`$ is found from the equation of motion, $`\dot{𝐩}=(e/c)𝐯\times 𝐁`$. In general, $`𝐁`$ may vary both in amplitude and in direction. In relativistic GRB shocks, radiation is beamed; only a conical section bound by the opening angle $`\theta _b1/\gamma _{\mathrm{sh}}`$ is seen by an observer, where $`\gamma _{\mathrm{sh}}`$ is the bulk Lorentz factor of the expanding shell. At a particular observing time, most of the radiation is coming from the brightened limb (Sari 1998; Panaitescu & Mészáros 1998). Because of relativistic aberration (Rybicki & Lightman 1979), the shock front is seen almost edge-on at the limb. On the other hand, the magnetic fields in the shock shell are random but always lie in the plane of the shock front (Medvedev & Loeb 1999). Choosing a coordinate system with $`\widehat{x}`$-direction pointing towards the observer (in the shock frame) and $`\widehat{x}`$-$`\widehat{y}`$-plane being the shock front plane at the limb, we have $`𝐯=c\widehat{x}`$, $`𝐁_{}=B\widehat{y}`$, and $`\dot{𝐩}\gamma m_ew\widehat{z}`$. Then, the Fourier component of $`w`$ is
$$w_\omega ^{}=\frac{eB_\omega ^{}}{\gamma m_e}=\frac{eB_k^{}}{\gamma m_ec}.$$
(16)
For the magnetic field distribution, we use the model (10).
At last, the spectral power $`P(\omega )dW/dtd\omega `$ (i.e., the energy emitted per unit time per unit frequency interval, measured in ergs s<sup>-1</sup> Hz<sup>-1</sup>) is obtained dividing the energy spectrum $`dW/d\omega `$ by the pulse duration, $`T`$ (Rybicki & Lightman 1979). Combining equations (16), (10), (11) together and substituting into (15), we obtain
$$P(\omega )=r_e^2c\gamma ^2\frac{\overline{B}_{SS}^2}{2\omega _j}J\left(\frac{\omega }{\omega _j}\right),$$
(17)
where $`\omega /\omega _j2`$, and $`r_e=e^2/m_ec^2`$ is the classical electron radius. As is expected, the characteristic frequency of the emergent radiation is
$$\omega _j=\gamma ^2k_{Be}c=2^{7/4}\gamma ^2\gamma _{\mathrm{int}}\overline{\gamma }_e^{1/2}\omega _{pe},$$
(18)
cf., equations (5) and (8). The function $`J`$ is defined as
$$J(\xi )=(2\mu +1)\xi ^{2\mu }\left[I\left(\text{min}[2;\frac{\xi }{\delta }]\right)I(\xi )\right],$$
(19)
where $`I`$ is the integral,<sup>4</sup><sup>4</sup>4 We keep this integral in a general form because it contains logarithmic terms for $`\mu =0.5,1,1.5`$. $`I(\xi )=\xi ^{2\mu }\left(1\xi +\frac{1}{2}\xi ^2\right)𝑑\xi `$, and $`\text{min}[a;b]`$ denotes the smallest of $`a`$ and $`b`$. From equations (2), (4), (18), we have $`\omega _j=\gamma ^2\overline{B}_e/\delta m_ec`$. Equation (17) may now be cast into the form,
$$P(\omega )=\frac{e^2}{2c}\delta ^2\frac{\omega _j}{\gamma ^2}\frac{\overline{B}_{SS}^2}{\overline{B}_e^2}J\left(\frac{\omega }{\omega _j}\right).$$
(20)
This spectrum is shown in Figure 4.1 for several values of $`\mu `$. In general, the steeper the field distribution, $`B_kk^\mu `$, the more peaked is the radiation spectrum and the closer the peak frequency to $`2\omega _j`$. Note, the discontinuity of the slope is artificial. This is due to the $`k_{\mathrm{min}}`$-cutoff because the model we use does not continuously interpolate between the jitter and synchrotron radiation limits (i.e., between $`\delta 1`$ and $`\delta 1`$ cases). This discontinuity is less prominent for large $`\mu `$’s and small $`\delta `$’s.
The total emitted power may be obtained by integrating (17) over frequencies. A simpler way is to express it in terms of electron’s acceleration, i.e., $`dW/dt=(2e^2\gamma ^4/3c^3)(w_{}^2+\gamma ^2w_{}^2)`$ (Rybicki & Lightman 1979) and substitute $`w_{}=0`$ and $`w_{}`$ given by (16). The result is
$$dW/dt=(2/3)r_e^2c\gamma ^2\overline{B}_{SS}^2.$$
(21)
Note, this expression is identical to that for a synchrotron radiation in which a uniform field, $`B_{}^2`$, is replaced with its average, $`\overline{B}_{SS}^2=B_{}^2`$.
We are now able to calculate the jitter radiation spectrum from an ensemble of electrons. We assume that electrons are accelerated in the shock<sup>5</sup><sup>5</sup>5Small-scale magnetic fields present in the shock provide effective collisions in the otherwise collisionless plasma and make the Fermi acceleration to operate (Medvedev & Loeb 1999). to a power-law distribution $`N(\gamma )=C_N\gamma ^p`$ (where $`C_N`$ is a normalization constant) with a minimum Lorentz factor, i.e., $`\gamma _{\mathrm{min}}\gamma <\mathrm{}`$. The index $`p`$ must be $`p2`$ so that energy does not diverge at large $`\gamma `$’s. We assume the standard value,<sup>6</sup><sup>6</sup>6Resent studies indicate that the high-frequency spectral power index in the prompt GRBs is $`1`$ which translates to $`p3`$. We keep $`p=2.5`$ for illustrative purposes. $`p=2.5`$ (Sari, Narayan, & Piran 1996), unless stated otherwise. The spectrum is found from
$$P_{\mathrm{ens}}(\omega )=_{\gamma _{\mathrm{min}}}^{\mathrm{}}N(\gamma )P(\omega )𝑑\gamma ,$$
(22)
noting that $`\omega _j\omega _j(\gamma )\gamma ^2`$, i.e., in exact analogy with synchrotron radiation, $`\omega _c\gamma ^2`$. We obtain
$$P_{\mathrm{ens}}(\omega )=\frac{C_N}{2\gamma _{\mathrm{min}}^{p1}}\left(\frac{\omega }{\omega _{jm}}\right)^{\frac{p1}{2}}_0^{\omega /\omega _{jm}}\xi ^{(p1)/2}P(\xi )𝑑\xi ,$$
(23)
where we introduced the characteristic frequency of radiation $`\omega _{jm}\omega _j(\gamma _{\mathrm{min}})`$, cf. equation (18),
$$\omega _{jm}=2^{7/4}\gamma _{\mathrm{min}}^2\gamma _{\mathrm{int}}\overline{\gamma }_e^{1/2}\omega _{pe},$$
(24)
At low and high frequencies, the spectra may be obtained analytically as follows,
$$P_{\mathrm{ens}}(\omega )\{\begin{array}{cc}\left(\frac{\omega }{\omega _{jm}}\right)^1,\hfill & \omega \omega _{jm},\hfill \\ \left(\frac{\omega }{\omega _{jm}}\right)^{(p1)/2},\hfill & \omega \omega _{jm},\hfill \end{array}$$
(25)
At high frequencies, the spectrum is analogous to the synchrotron case. The electron power-law index determines the slope. At low frequencies, the spectrum is linear in frequency, in contrast to the synchrotron spectrum which is softer and scales as $`P_{\mathrm{syn}}(\omega )\omega ^{1/3}`$. The spectrum of radiation from single-speed electrons is peaked at $`\omega 2\omega _j`$ for large $`\mu `$’s, i.e., for a peaked magnetic field distribution, as discussed above (see also Figure 4.1). Therefore, the spectral break due to the $`\gamma _{\mathrm{min}}`$-cutoff occurs at
$`\nu _{j,\mathrm{break}}`$ $``$ $`2\omega _{jm}/2\pi =(2^{7/4}/\pi )\gamma _{\mathrm{min}}^2\gamma _{\mathrm{int}}\overline{\gamma }_e^{1/2}\omega _{pe}`$ (26)
$``$ $`6.0\times 10^9\gamma _{\mathrm{min}}^2\gamma _{\mathrm{int}}\overline{\gamma }_e^{1/2}n_{10}^{1/2}\text{ Hz},`$
where $`n_{10}n/10^{10}\text{ cm}^3`$. The above frequency is calculated in the frame of the relativistic expanding shell. This frequency is boosted in observer’s frame by a factor of $`\gamma _{\mathrm{sh}}`$.
The model of jitter radiation contains two extra parameters, compared to the model of synchrotron radiation. These parameters are the deflection-to-beaming ratio, $`\delta `$, and the magnetic field index, $`\mu `$, which determines the peakedness of the magnetic field over spatial scales. We now investigate how the spectrum of the emergent radiation depends on these parameters.
Figure 4.1 represents the spectral power, $`P(\omega )`$, given by equations (20), (23), as a function of $`\omega /\omega _{cm}`$, where $`\omega _{cm}\omega _c(\gamma _{\mathrm{min}})`$, for various values of $`\delta `$. The dashed curve represents the synchrotron spectrum which corresponds to the same magnetic field strength. Note, the frequency is normalized by the synchrotron break frequency, $`\omega _{cm}`$ (not the jitter break frequency, $`\omega _{jm}`$), to emphasize that, for the fixed field magnitude, $`\overline{B}_{SS}^2`$, the jitter frequency increases with decreasing $`\delta `$, in accordance with equation (3). One can see that jitter radiation is well described by a broken power-law given by equation (25). Notice also little change in the spectrum shape near the break frequency and the overall decrease of $`P(\omega )`$ as $`\delta `$ decreases.
Figure 4.1 shows $`P(\omega )`$ vs. $`\omega /\omega _{jm}`$ for a few values of the magnetic field index, $`\mu `$. In general, the ratio $`\overline{B}_{SS}^2/\overline{B}_e^2`$ depends on $`\mu `$. To highlight the effect of the magnetic spectrum shape alone, we keep $`\overline{B}_{SS}^2/\overline{B}_e^2`$ fixed for all $`\mu `$’s. One can see from Figure 4.1 that $`\mu `$ has little effect on the overall shape of the radiation spectrum. However, as one goes from a flat ($`\mu =0`$) to a peaked ($`\mu =10`$) magnetic field distribution, the change in slope at the break frequency becomes more abrupt.
### 4.2. Synchrotron radiation, $`\delta 1`$
For completeness, we also consider the case of $`\delta 1`$. In this case, $`\alpha /\mathrm{\Delta }\theta 1`$ while still $`\lambda _B/\rho _L1`$ for energetic electrons. The spectral power emitted by a single electron is given as (Landau & Lifshits 1975, §77 and §74)
$$\frac{dW}{dtd\omega }=\frac{e^2}{\sqrt{3}\pi c\gamma ^2}\stackrel{~}{\omega }_cF\left(\frac{\omega }{\stackrel{~}{\omega }_c}\right),$$
(27)
where $`\stackrel{~}{\omega }_c=\frac{3}{2}\gamma ^2e\stackrel{~}{B}_{}/m_ec`$$`\mathrm{}`$ denotes the average over the electron’s trajectory, and “tilde” denotes the local (instantaneous) quantity. Here $`F(\xi )=\xi _\xi ^{\mathrm{}}K_{5/3}(\xi ^{})𝑑\xi ^{}`$ and $`K_{5/3}(x)`$ is the modified Bessel function of $`5/3`$ order. As for the spectrum from a plasma, the observed radiation comes from regions with different field strength. Therefore, one has to replace $`\stackrel{~}{B}_{}`$ in (27) with the ensemble average, $`\overline{B}=\sqrt{\stackrel{~}{B}_{}^2}`$, which rigorously yields the standard synchrotron spectrum.
## 5. Jitter+Synchrotron model of GRB emission
We now use the model of magnetic fields discussed in $ 3.2 to construct a GRB spectral model All calculations are performed in the frame of the expanding shell. The transition to observer’s frame is obvious. The magnetic field in the shock shell was shown to be sub-divided into small-scale and large-scale components. The small-scale field, for which $`\delta 1`$, yields jitter radiation. Its spectral power is given by equations (20), (23). The large-scale field, for which $`\delta 1`$, yields synchrotron radiation. The relevant theory may be found in Rybicki & Lightman (1979). The composite spectrum, thus, contains both jitter and synchrotron components and may be schematically represented as follows,
$$P_{J+S}(\omega )=P_J(\omega ;\overline{B}_{SS},\delta ,\mu )+P_S(\omega ;\overline{B}_{LS}).$$
(28)
This approximation is valid for the field distribution from §3.2 unless $`\delta 1`$, as discussed in the beginning of §4. It is convenient to normalize frequencies onto the jitter frequency, $`\omega _{jm}`$, which is independent of the magnetic field strength, cf., equation (26). The synchrotron-to-jitter frequency ratio is then
$$\frac{\omega _{cm}}{\omega _{jm}}=\frac{\omega _c}{\omega _j}\frac{3}{2}\frac{\overline{B}_{LS}}{\overline{B}_{SS}}\delta ,$$
(29)
as follows from equations (3), (13).
Figure 5 represents the spectrum emitted by single-speed electrons for the same values of $`\mu `$ as in Figure 4.1. One can clearly see a sharp feature on top of the broad synchrotron spectrum. Integrating this spectrum over the power-law distribution of $`\gamma `$’s, we obtain the composite “jitter+synchrotron” (JS) model of GRBs. A typical example is shown in Figure 5. In general, there are two bumps: the sharp one is near the jitter frequency and the other, broad bump, is associated with synchrotron emission. Depending on the ratio $`\overline{B}_{LS}^2/\overline{B}_{SS}^2`$, these bumps may overlap to produce either featureless broad or sharply peaked spectra, as well. The high-frequency tail always scales as $`P_{J+S}(\omega )\omega ^{(p1)/2}`$. The dependence of the spectrum from $`\delta `$ is displayed in Figure 5. As an example, we take $`\overline{B}_0^2/8\pi =\overline{B}_{SS}^2/8\pi `$ and $`\mu =10`$ (for such $`\mu `$’s $`\overline{B}_{SS}^2/8\pi \overline{B}_e^2/8\pi `$, and, therefore $`\overline{B}_{LS}^2/8\pi \overline{B}_0^2/8\pi `$). A sharp jitter feature is easily seen in the spectrum. As $`\delta `$ decreases, the synchrotron bump moves towards lower frequencies and decreases in amplitude. The jitter peak decreases even faster and the spectral feature becomes less prominent. The other parameter, $`\mu `$, determines the ratio $`\overline{B}_{LS}^2/\overline{B}_{SS}^2`$ but, besides this, its effect onto the spectrum is weak (cf., Figure 4.1) and is, therefore, not shown.
The spectrum of radiation depends on the relative magnitudes of the large- and small-scale field components. Figure 5 represents spectra for $`\mu =10`$ and various values of the ratio $`\overline{B}_0^2/\overline{B}_{SS}^2\overline{B}_{LS}^2/\overline{B}_{SS}^2`$. Figure 5a shows the spectral power vs. frequency. For strong large-scale fields, the synchrotron spectral component dominates. As $`\overline{B}_{LS}^2/\overline{B}_{SS}^2`$ decreases, the synchrotron peak moves towards lower frequencies. The amplitude of the synchrotron peak decreases too. The position and the amplitude of the jitter peak remain almost unchanged. Thus, in the limit $`\overline{B}_{LS}^2/\overline{B}_{SS}^20`$, the spectrum becomes purely jitter. It is illustrative to depict the spectral flux, $`F(\omega )P(\omega )/\omega `$ \[observationally, this quantity is proportional to the photon spectral flux, $`N(E)`$, measured in units: photons s<sup>-1</sup> cm<sup>-2</sup> keV<sup>-1</sup>\], as presented in Figure 5b. It is seen that the slope of the flux below the jitter break continuously decreases as $`\overline{B}_{LS}^2/\overline{B}_{SS}^2`$ decreases. The synchrotron asymptotic slope $`\omega ^{2/3}`$ is shown for comparison. It is interesting that for $`\overline{B}_{LS}^2/\overline{B}_{SS}^210^1`$ (the actual value depends on $`\delta `$; the larger $`\delta `$, the larger the ratio), there is a portion in the spectrum which has the power-law index being less than $`2/3`$ and approaching zero in the limit of the vanishing large-scale field. We discuss this property in §6 in the context of the “line of death” for synchrotron radiation in GRBs.
Finally, it is worthwhile to compare the peak spectral fluxes of jitter, $`F_{J,\mathrm{max}}F_J(\omega _{jm})`$, and synchrotron, $`F_{S,\mathrm{max}}F_S(\omega _{cm})`$, radiation,
$$\frac{F_{J,\mathrm{max}}}{F_{S,\mathrm{max}}}=f(p,\mu )\delta ^2,$$
(30)
where $`f(p,\mu )`$ is a function of two power-law indices, $`p`$ and $`\mu `$. The above equation may be readily obtained from equations (23), (20), and, for instance, (27). For $`\mu 1`$, the function $`f`$ depends on $`\mu `$ only weakly while, given the spectrum, $`p`$ is found by fitting the large frequency slope. Therefore, for a given spectrum with two sub-components, both significant parameters, namely $`\overline{B}_{LS}/\overline{B}_{SS}`$ and $`\delta `$, are uniquely found from equations (29) and (30).
## 6. Comparison with observations
In §5 we constructed a composite, jitter+synchrotron model of radiation emitted from internal shocks of GRBs, assuming that the magnetic fields are produced in these shocks via the relativistic two-stream instability (Medvedev & Loeb 1999). Several observational predictions can now be drawn. We compare them with presently available observational data. For future reference, we define the photon spectral index, $`s`$, as follows,
$$F(\omega )\omega ^s.$$
(31)
Here we remind that $`F(\omega )P(\omega )/\omega N(E)`$, where $`N(E)`$ is the number of observed photons per unit time, per unit energy range, per unit area.
We have shown that prompt GRB spectra consist of two spectral sub-components, namely synchrotron and jitter. The synchrotron component is, as usual, well approximated by the smoothly broken power-law, also referred to as the Band function (Band et al. 1993) or the GRB function. The jitter component is better approximated by a sharply broken power-law with the hard and soft photon indexes being equal to $`s=(p+1)/2`$ and $`s=0`$ respectively. The position of the jitter peak is independent of the magnetic field strength, in contrast to the synchrotron peak, but depends on the particle (electron) density in the relativistic expanding shell, see equation (26).
### 6.1. Low-frequency spectra and the “Line of Death”
The optically thin synchrotron model of GRBs makes a solid prediction. Namely, the soft photon spectral index must be in the range $`3/2s2/3`$, depending on the strength of the magnetic field if strong synchrotron cooling of the emitting electrons in the fireball is taken into account (Katz 1994; Sari, Narayan, & Piran 1996). Thus, the photon index in this model can newer be greater than $`2/3`$, creating a testable “line of death” for the synchrotron shock model (Katz 1994) . There is, however, a growing observational evidence that many bursts violate this prediction (Crider et al. 1997; Strohmayer et al. 1998; Preece et al. 1998; Frontera et al. 1999). For instance, Preece et al. (1998) have studied time-resolved spectra of the bursts collected by the LAD from the BATSE instrument. They found 23 out of 137 bursts which violate the optically thin synchrotron model. Frontera et al. (1999) demonstrated that about $`50\%`$ of time-resolved spectra the bursts observed by the Wide Field Camera on board of the BeppoSAX before May 1998 also violate this model. The proposed explanations, which include Compton up-scattering of low-energy photons (Liang et al. 1997), synchrotron self-absorption (Papathanassiou 1999), and the influence of the pair annihilation in the fireball photosphere (Eichler & Levinson 1999), though possible, seem rather ad hoc and suffer from drawbacks. The Compton up-scattering model strictly requires a single up-scattering event per photon which requires the column density to self-adjust to a few g/cm<sup>2</sup>. The self-absorption model results in large optical depths and, thus, weak emitted flux and very low radiation efficiency. The photospheric model requires an extremely low baryonic load. For more discussion see Mészáros (1999).
The JS model provides a natural resolution of this puzzle. Figure 6.1 represents the photon index as a function of a logarithm of frequency for the synchrotron and jitter cases<sup>7</sup><sup>7</sup>7Little spikes on the lower curve are numerical artifacts. When jitter radiation dominates, the low-energy photon index tends towards its asymptotic value of $`s=0`$, thus shifting the “line of death” to harder spectra. In principle, the indexes as large as $`s=1/2`$ are allowed by our model. Positive $`s`$, however, are quite unlikely. These predictions are in good agreement with observational data. Namely, only two bursts out of those studied by Preece et al. (1998) are above the $`s=0`$ line by more than $`1\sigma `$ and all bursts studied by Frontera et al. (1999) except the peculiar one, GRB 970111, fall below this line. At last, these three bursts are still consistent with $`s=1/2`$ within statistical uncertainties.
### 6.2. Sharply broken power-law spectra of GRBs
Most of the observed bursts are well-fitted by the GRB function (Band et al. 1993). However, there are bursts with a sharp curvature of the spectrum at the break energy. For such bursts a broken power-law model generates better $`\chi ^2`$ fits than the GRB function (see, e.g., Preece et al. 1998). In terms of the JS model, there are two cases in which such spectra occur. They can be observationally distinguished by the value of the soft (i.e., low-energy) photon index. First, it is the case of purely jitter radiation: $`\overline{B}_{LS}^2/\overline{B}_{SS}^20`$ and there is no synchrotron peak (or this peak is too weak and is at low energies, outside the detector range), see Figure 5. These spectra are flat, $`s0`$, at low energies. Second, for large $`\delta 1`$ and the equipartition between the large- and small-scale fields, $`\overline{B}_{LS}^2/\overline{B}_{SS}^21`$, similar spectral shapes are also obtained, see Figure 5. In this case, the soft photon index is $`s2/3`$, or it is in the range $`3/2s2/3`$ if synchrotron cooling of the electrons is taken into account (Sari, Narayan, & Piran 1996). Thus, we expect fewer broken power-law bursts which have the soft photon indexes in the range $`2/3s0`$. Figure 2 of Preece et al. (1998) indeed reveals a gap between the power-law bursts with $`s0`$ and $`s2/3`$. However, a larger number of such bursts is required to draw a statistically significant conclusion. We also should point out that if no ordered field is present in the ejecta, then $`\overline{B}_{LS}\overline{B}_{SS}`$ is equivalent to $`\overline{B}_p^2/8\pi \overline{B}_e^2/8\pi `$. This, in turn, naturally implies strong energy coupling and rough equipartition between the protons and the electrons, as discussed in §3.1.
### 6.3. Two-component spectral structure of GRB emission
The analysis done by Pendleton et al. (1994) of the bursts from the first BATSE catalog collected by the Large Area Detector (LAD) demonstrates that there is a large number of bursts which have the high-energy photon indexes ($`50300`$ keV range) being much larger than the low-energy ones ($`20100`$ keV range). No such behavior is observed at other energy ranges. Thus, this result may indicate the presence of more than one spectral component in the energy range $`40100`$ keV. (The absence of similar behavior at higher energies seems to rule out the inverse-Compton origin of a second spectral component.) On the other hand, such spectra are completely consistent with the JS model. They generally correspond to the well-separated synchrotron and jitter peaks and are characterized by small $`\delta `$’s, $`\delta 0.1`$, and weak large-scale fields, $`\overline{B}_{LS}^2/\overline{B}_{SS}^20.01`$, as is clearly seen from Figure 5.
### 6.4. “Lines” in GRB spectra
Whether emission and/or absorption features, often referred to as lines, in GRB spectra are real remains an open question for a long time see, e.g., Briggs 1999 and Ryde 1999 for discussion. These features have been observed, for instance, by the KONUS experiment, Venera mission (Mazets et al. 1981) and by Ginga (Murakami et al. 1988). At that time, it was widely believed that these features are cyclotron lines due to a highly magnetized progenitor. The BATSE Spectroscopy Detectors (SD) refined these results. Palmer et al. (1994) studied almost two hundred bursts detected by SDs. No convincing line features were found. In a more recent analysis by Briggs (1999) of more than a hundred bursts observed by SDs before May 31, 1996, about 10 highly significant line features have been found. These are low-energy ($`50`$ keV) emission features. Clearly, this controversy requires further studies. However, it seems likely that if lines exist, they are rare.
Not too surprisingly, the two-component JS model is able to produce spectra with emission line-like features. These are not lines in a strict sense because at energies higher than the “line” peak, the spectrum has no curvature and falls down as a power-law with $`s=(p+1)/2`$. Few such spectra are seen in Figure 5. These spectra are obtained (i) for rather low values of $`\delta `$ ($`\delta 0.3`$), otherwise the synchrotron and jitter peaks overlap, and (ii) for a narrow range of $`\overline{B}_{LS}^2/\overline{B}_{SS}^2`$ ($`0.1\overline{B}_{LS}^2/\overline{B}_{SS}^21`$), otherwise either the spectral components are too well separated if the value of this ratio is small, or the jitter peak is too week and unobservable if the ratio is large. We thus conclude that the spectra with emission features are not very common. They require some parameter tuning and, therefore, must be rare. Moreover, inhomogeneities and highly variable conditions in the GRB shocks may smear out weak spectral features completely.
At last, we present an illustrative example. Figure 1 shows the spectrum of GRB 910503 obtained using all capability of CGRO’s four experiments (Schaefer et al. 1998). Here the energy flux $`\nu F_\nu E^2N(E)\omega P(\omega )`$ in units of (photons cm<sup>-2</sup> s<sup>-1</sup> keV<sup>-1</sup>)$`\times (E/100`$ keV)<sup>2</sup> is plotted vs. frequency in units of keV. A sharp second peak at $`2`$ MeV is clearly seen. The solid curve in the figure represents a visual (i.e., not $`\chi ^2`$) spectral shape fit using the JS model. Here we fit the separation of the two peaks and their relative amplitudes; we didn’t fit the absolute flux and position of the synchrotron spectral break which depend on the energetics of the fireball. Note that the intensity of a second, high-frequency (jitter) component is low. The values inferred from this fit are the following: $`\overline{B}_{LS}^2/\overline{B}_{SS}^2=7`$$`\delta =0.07`$$`p=3.9`$, and we took $`\mu =10`$. We calculate the efficiency of the magnetic field generation in the shock from equation (2) assuming $`\varphi =1/4`$ as a typical value. We obtain $`\eta _e0.08`$. The total magnetic field energy is calculated as $`ϵ_B=ϵ_{Be}(1+\overline{B}_{LS}^2/\overline{B}_{SS}^2)`$, where $`ϵ_{Be}`$ is found from equation (4) and we used $`\overline{B}_{SS}\overline{B}_e`$ for large $`\mu `$. We obtain $`ϵ_B4\times 10^4`$, which is in agreement with the conclusion drawn by Chiang & Dermer (1999) from a completely different analysis that the magnetic field in this burst is well below the equipartition, $`ϵ_B10^2`$. Their argument goes from the fact that the low-energy spectrum of this burst is consistent with $`F(\omega )\omega ^{2/3}`$. Evidently, synchrotron losses are small and the emitting electrons do not form a cooled distribution, which otherwise would result in $`F(\omega )\omega ^{3/2}`$. This is possible only for very small values of $`ϵ_B`$.
## 7. Conclusion
In this paper we have shown that radiation produced by relativistic electrons in magnetic fields may differ quite substantially from synchrotron. This radiation, referred here to as jitter radiation, is produced in the magnetic fields which are highly inhomogeneous on very small spatial scales. Such fields are likely to be present in GRB shocks. We developed a quantitative theory of jitter radiation. Jitter radiation has a different spectrum and its peak frequency is independent of the field strength. However, the total (i.e., frequency integrated) emitted power of jitter radiation depends on the field strength and is exactly identical to that of synchrotron radiation. We also constructed a composite, two-component jitter+synchrotron spectral model of the prompt GRB emission. Predictions of this model seem to be in excellent agreement with presently available data and, likely, resolve some puzzling spectral properties of the prompt $`\gamma `$-ray emission. All this, we think, strongly supports that (i) the proposed jitter radiation mechanism operates in astrophysical objects and (ii) the magnetic field is generated in shocks by the two-stream instability. We emphasize that a reliable identification/detection of the jitter spectral features will provide a direct evidence that the magnetic field in GRBs is due to the two-stream instability, since we presently unaware of any other mechanism which is capable of producing the required small-scale, large-amplitude fields. In general, the detection of both spectral components in GRB spectra would be a powerful and presize tool to investigate the properties of cosmological fireballs.
It is important to emphasize that the phenomenon of jitter radiation is intrinsic not only to internal shocks. Similar conditions (i.e., strong, small-scale fields) are expected to occur in external shocks which produce delayed afterglows, as well as in more conventional supernova shocks and relativistic jets. More speculatively, the above mechanism could allow to study the magnetic and electric fields in reconnection regions (remember, reconnection occurs on the electron skin depth scales) and small-scale structure of magnetic turbulence and cascade (e.g., in the interstellar medium).
The theory presented in this paper is only a “first step”. It is incomplete in a sence that it does not include the effects of both ordered and large-scale random magnetic fields in a self-consistent way and, hence, does not smoothly interpolate between the two extremes of synchrotron and jitter mechanisms. The synchrotron self-absorption has been omitted. The model also does not consider some related effects. For instance, it is likely that random electric fields (i.e., strong Langmuir turbulence) are present in the cosmological collisionless shocks. These fields will definitely affects the radiation spectrum via a similar mechanism. It is also unclear now whether particle’s motion is ergodic in random fields (i.e., whether a particle “samples” strong and weak fields statistically homogeneously), and what could be the effect of nonergodicity on the observed spectrum. All these issues will be addressed in future publications.
The author is grateful to Ramesh Narayan and Norm Murray for their interest in this work, various insightful comments, and useful discussions and to George Rybicki and Avi Loeb for discussions. This work has been supported by NSF grant PHY 9507695. |
warning/0001/hep-ph0001258.html | ar5iv | text | # 𝐾_{𝑙4} decays
## 1 Introduction
There is rich physics in kaon decays. Study on rare kaon decays are still active. The theoretical study of $`K_{l4}`$ decays has a long history.
In Ref. we have proposed an effective theory of large $`N_C`$ QCD of mesons. In this theory the diagrams at the tree level are at the leading order in large $`N_C`$ expansion and the loop diagrams of mesons are at higher orders. This theory is phenomenologically successful. We have used this theory to study $`K_{l3}`$, $`Ke\nu \gamma `$, kaon form factors, and $`\pi `$K scattering in the chiral limit. Theoretical results agree well with data. In these studies VMD and PCAC are satisfied. There are five parameters in this theory: three current quark masses, a parameter related to the quark condensate, and a universal coupling constant g which is determined to be 0.39 by fitting $`\rho ee^+`$. All parameters have been fixed by previous studies.
In this paper we use this theory of pseudoscalar, vector, and axial vector mesons to study $`K^{}\pi ^+\pi ^{}l\nu ,\pi ^0\pi ^0l\nu `$, and $`K_L\pi ^\pm \pi ^0l^{}\nu `$. In the study of $`K_{l4}`$ there is no adjustable parameter.
The Lagrangian of the theory of Ref. 3Y is
$`=\overline{\psi }(x)(i\gamma +\gamma v+\gamma a\gamma _5mu(x))\psi (x)+{\displaystyle \frac{1}{2}}m_1^2(\rho _i^\mu \rho _{\mu i}+\omega ^\mu \omega _\mu +a_i^\mu a_{\mu i}+f^\mu f_\mu )`$
$`+{\displaystyle \frac{1}{2}}m_2^2(K_\mu ^a\overline{K}^{a\mu }+K_1^\mu K_{1\mu })+{\displaystyle \frac{1}{2}}m_3^2(\varphi _\mu \varphi ^\mu +f_s^\mu f_{s\mu })`$
$`+\overline{\psi }(x)_L\gamma W\psi (x)_L+_W+_{lepton}\overline{\psi }M\psi ,`$ (1)
where $`a_\mu =\tau _ia_\mu ^i+\lambda _aK_{1\mu }^a+(\frac{2}{3}+\frac{1}{\sqrt{3}}\lambda _8)f_\mu +(\frac{1}{3}\frac{1}{\sqrt{3}}\lambda _8)f_{s\mu }`$($`i=1,2,3`$ and $`a=4,5,6,7`$), $`v_\mu =\tau _i\rho _\mu ^i+\lambda _aK_\mu ^{}+(\frac{2}{3}+\frac{1}{\sqrt{3}}\lambda _8)\omega _\mu +(\frac{1}{3}\frac{1}{\sqrt{3}}\lambda _8)\varphi _\mu `$, $`W_\mu ^i`$ is the W boson, and $`u=exp\{\gamma _5i(\tau _i\pi _i+\lambda _aK^a+\eta +\eta ^{})\}`$, $`m`$ is a parameter, and M is the mass matrix of u, d, s quarks, The masses $`m_1^2`$, $`m_2^2`$, and $`m_3^2`$ have been determined theoretically.
Using the notations of Ref., we have
$`<\pi ^i\pi ^j|A_\mu |K>={\displaystyle \frac{i}{m_K}}\{(p_1+p_2)_\mu F^{ij}+(p_1p_2)_\mu G^{ij}+q_\mu R^{ij}\},`$ (2)
$`<\pi ^i\pi ^j|V_\mu |K>={\displaystyle \frac{H^{ij}}{m_K^3}}\epsilon ^{\mu \nu \lambda \rho }p_\nu (p_1+p_2)_\lambda (p_1p_2)_\rho ,`$
where $`p_1,p_2,p`$ are momenta of two pions and kaon respectively, $`q=pp_1p_2`$, and $`i,j=+,,0`$. We define
$$q_1^2=(pp_1)^2,q_2^2=(pp_2)^2,q_3^2=(p_1+p_2)^2.$$
The form factors, $`F^{ij},G^{ij},R^{ij}`$ and $`H^{ij}`$ are functions of $`q^2,q_1^2,q_2^2`$, and $`q_3^2`$. These four variables satisfy
$$q_1^2+q_2^2+q_3^2=m_K^2+2m_\pi ^2+q^2.$$
The paper is organized as: 1)introduction; 2)isospin relation; 3)form factors of vector current; 4)$`K^{}K\pi \pi `$ decay; 5)form factors of axial-vector current; 6)decay rates; 7)conclusions.
## 2 Isospin relation
For the decay modes $`K^{}\pi ^+\pi ^{}l\nu ,\pi ^0\pi ^0l\nu `$ and $`\overline{K^0}\pi ^+\pi ^0l\nu `$ there are isospin relations between the form factors denoted as $`A^{ij}`$. We take -$`\pi ^+`$, $`\pi ^0`$, and $`\pi ^{}`$ as isospin triplet and -$`\overline{K^0}`$ and $`K^{}`$ as isospin doublet. The isospin relation is obtained as
$$A^+=A^{00}\frac{1}{\sqrt{2}}A^{+0},$$
(3)
where $`A^{ij}=F^{ij},G^{ij},R^{ij},H^{ij}`$ respectively.
## 3 Form factors of vector current
The Vector Meson Dominance(VMD) is a natural result of this theory. The coupling between the W bosons and the bosonized vector current($`\mathrm{\Delta }s=1`$) has been derived as
$`^V`$ $`=`$ $`{\displaystyle \frac{g_W}{4}}sin\theta _Cg\{{\displaystyle \frac{1}{2}}(_\mu W_\nu ^+_\nu W_\mu ^+)(_\mu K_\nu ^{}_\nu K_\mu ^{}){\displaystyle \frac{1}{2}}(_\mu W_\nu ^{}_\nu W_\mu ^{})(_\mu K_\nu ^+`$ (5)
$`_\nu K_\mu ^+)+W_\mu ^+j_\mu ^{}+W_\mu ^{}j_\mu ^+\},`$
where $`j_\mu ^\pm `$ is obtained by substituting
$$K_\mu ^\pm \frac{g_W}{4}sin\theta _CgW_\mu ^\pm $$
into the vertex in which $`K_\mu `$ field is involved.
The matrix elements of the vector current of $`K_{l4}`$ are resulted in anomalous vertices of mesons. The two subprocesses are shown in Fig.1(a,b). There is contact term. Three kinds of vertices are involved: the contact term $`_{K^{}K\pi \pi }`$, $`_{K^{}K^{}\pi }`$ and $`_{K^{}K\pi }`$ , and $`_{K^{}K\rho }`$ and $`_{\rho \pi \pi }`$. In the chiral limit, $`m_q0`$, all these vertices have been derived from the Lagrangian (1) and are listed below
$`_{K^{}K^{}\pi }={\displaystyle \frac{N_C}{\pi ^2g^2f_\pi }}\epsilon ^{\mu \nu \alpha \beta }d_{aci}K_\mu ^a_\nu K_\alpha ^c_\beta \pi ^i`$ , (7)
$`_{K^{}K\pi }={\displaystyle \frac{2}{g}}f(q^2)f_{abi}K_\mu ^a(_\mu \pi ^iK^b\pi ^i_\mu K^b),`$
$`f(q^2)=1+{\displaystyle \frac{q^2}{2\pi ^2f_\pi ^2}}[(1{\displaystyle \frac{2c}{g}})^24\pi ^2c^2)],`$
$`c={\displaystyle \frac{f_\pi ^2}{2gm_\rho ^2}},`$
$`_{K^{}\rho K}={\displaystyle \frac{N_C}{\pi ^2g^2f_\pi ^2}}\epsilon ^{\mu \nu \alpha \beta }d_{abi}K_\mu ^a_\nu \rho _\alpha ^i_\beta K^b,`$
$`_{\rho \pi \pi }={\displaystyle \frac{2}{g}}f(q^2)ϵ_{ijk}\rho _\mu ^i\pi ^j_\mu \pi ^k,`$
$`_{K^{}K\pi \pi }={\displaystyle \frac{2}{g\pi ^2f_\pi ^3}}(1{\displaystyle \frac{6c}{g}}+{\displaystyle \frac{6c^2}{g^2}})d_{abe}f_{cde}\epsilon ^{\mu \nu \alpha \beta }K_\mu ^a_\nu P^b_\alpha P^c_\beta P^d,`$
From theses vertices the form factors $`h^{ij}`$ are found
$`H^+={\displaystyle \frac{m_K^3m_K^{}^2}{q^2m_K^{}^2}}\{{\displaystyle \frac{1}{\pi ^2f_\pi ^3}}(1{\displaystyle \frac{6c}{g}}+{\displaystyle \frac{6c^2}{g^2}}){\displaystyle \frac{N_C}{g^2\pi ^2f_\pi }}{\displaystyle \frac{f(q_2^2)}{q_2^2m_K^{}^2}}`$ (10)
$`{\displaystyle \frac{N_C}{2g^2\pi ^2f_\pi }}{\displaystyle \frac{f(q_3^2)}{q_3^2m_\rho ^2+i\sqrt{q_3^2}\mathrm{\Gamma }(q_3^2)}}\},`$
$`H^{00}={\displaystyle \frac{m_K^3m_K^{}^2}{q^2m_K^{}^2}}{\displaystyle \frac{N_c}{2g^2\pi ^2f_\pi }}\{{\displaystyle \frac{f(q_2^2)}{q_2^2m_K^{}^2}}{\displaystyle \frac{f^2(q_1^2)}{q_1^2m_K^{}^2}}\},`$
$`H^{+0}={\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \frac{m_K^3m_K^{}^2}{q^2m_K^{}^2}}\{{\displaystyle \frac{2}{\pi ^2f_\pi ^3}}(1{\displaystyle \frac{6c}{g}}+{\displaystyle \frac{6c^2}{g^2}})+{\displaystyle \frac{N_C}{\pi ^2g^2f_\pi }}[{\displaystyle \frac{f(q_1)}{q_1^2m_K^{}^2}}+{\displaystyle \frac{f(q_2)}{q_2^2m_K^{}^2}}`$
$`+{\displaystyle \frac{f(q_3)}{q_3^2m_\rho ^2+i\sqrt{q_3^2}\mathrm{\Gamma }_\rho (q_3^2)}}]\},`$
where $`\mathrm{\Gamma }_\rho `$ is the decay width of $`\rho `$ meson
$$\mathrm{\Gamma }_\rho (q_3^2)=\frac{\sqrt{q_3^2}f^2(q_3^2)}{12g^2\pi }(1\frac{4m_\pi ^2}{q_3^2})^{\frac{3}{2}}.$$
(11)
The equations(8-10) show that the isospin relation(3) is satisfied.
## 4 $`K^{}K\pi \pi `$ decay
The vertices(6,7) are responsible for the decay of $`K^{}K\pi \pi `$. As a test the decay widths of $`K^{}K\pi \pi `$ are calculated
$$\mathrm{\Gamma }(K^{}K^{}\pi ^+\pi ^{})=\frac{1}{96(2\pi )^3m_K^{}}𝑑k_1^2𝑑k_2^2\{p_1^2p_2^2(\stackrel{}{p}_1\stackrel{}{p}_2)^2\}|A|^2=0.29\times 10^5GeV$$
(12)
which is less than the experimental upper limit, where A is the amplitude
$`A`$ $`=`$ $`{\displaystyle \frac{4}{g\pi ^2f_\pi ^3}}(1{\displaystyle \frac{6c}{g}}+{\displaystyle \frac{6c^2}{g^2}}){\displaystyle \frac{4N_c}{g^3\pi ^2f_\pi }}{\displaystyle \frac{f(k_2^2)}{k_2^2m_K^{}^2+i\sqrt{k_2^2}\mathrm{\Gamma }_K^{}(k_2^2)}}`$ (13)
$`{\displaystyle \frac{2N_c}{g^3\pi ^2f_\pi }}{\displaystyle \frac{f(k_3^2)}{k_3^2m_\rho ^2+i\sqrt{k_3^2}\mathrm{\Gamma }_\rho (k_3^2)}}`$
where $`k_1^2=(p+p_1)^2`$, $`k_2^2=(p+p_2)^2`$, $`k_3^2=(p_1+p_2)^2`$, and $`p_1,p_2,p`$ are momenta of $`\pi ^+,\pi ^{}`$ and $`K^{}`$ respectively, $`\mathrm{\Gamma }_K^{}`$ is the decay width of $`K^{}`$
$$\mathrm{\Gamma }_K^{}(k_2^2)=\frac{f^2(k_2^2)}{2\pi g^2k_2^2}\{\frac{1}{4k_2^2}(k_2^2+m_K^2m_\pi ^2)^2m_K^2\}^{\frac{3}{2}}.$$
(14)
$`\mathrm{\Gamma }(K^{}K^{}\pi ^0\pi ^0)={\displaystyle \frac{1}{192(2\pi )^3m_K^{}}}{\displaystyle 𝑑k_1^2𝑑k_2^2\{p_1^2p_2^2(\stackrel{}{p}_1\stackrel{}{p}_2)^2\}}`$ (15)
$`{\displaystyle \frac{36}{\pi ^4g^6f_\pi ^2}}\{{\displaystyle \frac{f(k_1)}{k_1^2m_K^{}^2+i\sqrt{k_1^2}\mathrm{\Gamma }_K^{}(k_1^2)}}{\displaystyle \frac{f(k_2)}{k_2^2m_K^{}^2+i\sqrt{k_2^2}\mathrm{\Gamma }_K^{}(k_2^2)}}\}^2`$
$`=0.61\times 10^6GeV.`$
$$\mathrm{\Gamma }(K^{}\overline{K}^0\pi ^{}\pi ^0)=\frac{1}{96(2\pi )^3m_K^{}}𝑑k_1^2𝑑k_2^2\{p_1^2p_2^2(\stackrel{}{p}_1\stackrel{}{p}_2)^2\}|B|^2=0.38\times 10^4GeV,$$
(16)
where
$`B={\displaystyle \frac{8}{\sqrt{2}gf_\pi ^3}}(1{\displaystyle \frac{6c}{g}}+{\displaystyle \frac{6c^2}{g^2}})+{\displaystyle \frac{12}{\sqrt{2}\pi ^2g^3f_\pi }}\{{\displaystyle \frac{f(k_1)}{k_1^2m_K^{}^2+i\sqrt{k_1^2}\mathrm{\Gamma }_K^{}(k_1^2)}}`$ (17)
$`+{\displaystyle \frac{f(k_2)}{k_2^2m_K^{}^2+i\sqrt{k_2^2}\mathrm{\Gamma }_K^{}(k_2^2)}}+{\displaystyle \frac{f(k_3)}{k_3^2m_\rho ^2+i\sqrt{k_3^2}\mathrm{\Gamma }_\rho (k_3^2)}}\}.`$
Eq.(16) is compatible with the data 9Y.
## 5 Form factors of axial-vector current
In the chiral limit, the axial-vector part of the interaction between W-boson and mesons is expressed as
$`^{As}={\displaystyle \frac{g_W}{4}}{\displaystyle \frac{1}{f_a}}sin\theta _C\{{\displaystyle \frac{1}{2}}(_\mu W_\nu ^\pm _\nu W_\mu ^\pm )(^\mu K_a^\nu ^\nu K_a^\mu )+W^{\pm \mu }j_\mu ^{}\}`$ (18)
$`+{\displaystyle \frac{g_W}{4}}sin\theta _C\mathrm{\Delta }m^2f_aW_\mu ^\pm K^\mu +{\displaystyle \frac{g_w}{4}}\mathrm{sin}\theta _Cf_KW_\mu ^\pm ^\mu K^{},`$
where $`j_\mu ^\pm `$ are obtained by substituting $`K_{a\mu }^\pm \frac{g_W}{4f_a}sin\theta _CW_\mu ^\pm `$ into the vertex in which $`K_a`$ fields are involved,
$`f_a=g^1(1{\displaystyle \frac{1}{2\pi ^2g^2}})^{\frac{1}{2}},`$ (21)
$`\mathrm{\Delta }m^2=6m^2g^2=f_\pi ^2(1{\displaystyle \frac{f_\pi ^2}{g^2m_\rho ^2}})^1,`$
$`c={\displaystyle \frac{f_\pi ^2}{2gm_\rho ^2}}.`$
The mass of $`K_1`$ meson is determined by
$$(1\frac{1}{2\pi ^2g^2})m_{K_1}^2=6m^2+m_K^{}^2.$$
(22)
The numerical value is $`m_{k_1}=1.322GeV`$ which is compatible with the data .
Two subprocesses contribute to the matrix element of the axial-vector current. They are shown in Fig.2(a,b). The vertices of mesons involved in these processes are $`_{K_1K^{}\pi }`$, $`_{K^{}K\pi }`$ and $`_{K_1\rho K}`$, $`_{\rho \pi \pi }`$. There is a contact term $`_{K_1K\pi \pi }`$ too. However, the calculation shows that the contribution of the contact term is very small and negligible. In the chiral limit, these vertices have been derived from the Lagrangian(1)
$`_{K_1K^{}\pi }=f_{abi}\{A(p^2)K_{1\mu }^aK_\mu ^b\pi ^iBK_{1\mu }^aK_\nu ^b_{\mu \nu }\pi ^i+DK_{1\mu }^a^\mu (K_\nu ^b^\nu \pi ^i)\}`$ , (24)
$`_{K_1\rho K}=f_{abi}\{A(p^2)K_{1\mu }^a\rho _\mu ^iK^bBK_{1\mu }^a\rho _\nu ^i_{\mu \nu }K^b+DK_{1\mu }^a^\mu (\rho _\nu ^i^\nu K^b)\},`$
where
$`A(p^2)={\displaystyle \frac{2}{f_\pi }}gf_a\{{\displaystyle \frac{F^2}{g^2}}+p^2[{\displaystyle \frac{2c}{g}}+{\displaystyle \frac{3}{4\pi ^2g^2}}(1{\displaystyle \frac{2c}{g}})]`$ (28)
$`+q^2[{\displaystyle \frac{1}{2\pi ^2g^2}}{\displaystyle \frac{2c}{g}}{\displaystyle \frac{3}{4\pi ^2g^2}}(1{\displaystyle \frac{2c}{g}})]\},`$
$`F^2=f_\pi ^2(1{\displaystyle \frac{2c}{g}})^1,`$
$`B={\displaystyle \frac{2}{f_\pi }}gf_a{\displaystyle \frac{1}{2\pi ^2g^2}}(1{\displaystyle \frac{2c}{g}}),`$
$`D={\displaystyle \frac{2}{f_\pi }}f_a\{2c+{\displaystyle \frac{3}{2\pi ^2g}}(1{\displaystyle \frac{2c}{g}})\},`$
where q and p are the momentum of $`K_1`$ and the vector meson respectively.
$`_{K^{}K\pi }={\displaystyle \frac{2}{g}}f_{abi}f(p^2)K_\mu ^a(K^b_\mu \pi ^i\pi ^i_\mu K^b),`$ (31)
$`_{\rho \pi \pi }={\displaystyle \frac{2}{g}}ϵ_{ijk}f(p^2)\rho _\mu ^i\pi ^j_\mu \pi ^k,`$
$`f(p^2)=1+{\displaystyle \frac{p^2}{2\pi f_\pi ^2}}[(1{\displaystyle \frac{2c}{g}})^24\pi ^2c^2],`$
where p is the momentum of the vector meson.
By using Eqs.(18,23,24), we obtain
$`<\pi ^+\pi ^{}|A_\mu |K^{}>`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}({\displaystyle \frac{q_\mu q_\nu }{q^2}}g_{\mu \nu }){\displaystyle \frac{g^2f_am_K^{}^2}{q^2m_{K_1}^2}}<\pi ^+\pi ^{}|\{A(p_K^{})\overline{K^0}_\nu \pi ^{}B\overline{K^0}_\lambda _{\lambda \nu }\pi ^{}\}`$ (32)
$`{\displaystyle \frac{1}{\sqrt{2}}}\{A(p_\rho )\rho _\nu ^0K^{}B\rho _\lambda ^0_{\lambda \nu }K^{}\}|K^{}>.`$
In the chiral limit PCAC is satisfied. The reason is that the Lagrangian(1) is chiral symmetric in the limit $`m_q0`$. On the other hand, the satisfaction of PCAC is resulted in the cancellations between the four terms of Eq.(18). The Eq.(18) shows that the axial-vector current has more complicated structure than the vector current does(5). Because of the PCAC the form factor R(2) is not an independent quantity and determined as
$$R=\frac{1}{q^2}\{q(p_1+p_2)F+q(p_1p_2)G\}.$$
(33)
Substituting the vertices(29,30) into Eq.(32), the three form factors are obtained
$`F^+={\displaystyle \frac{gf_am_K^{}^2m_K}{q^2m_{K_1}^2}}\{{\displaystyle \frac{f(q_2^2)}{q_2^2m_K^{}^2}}[{\displaystyle \frac{3}{2}}A(q_2^2)+{\displaystyle \frac{1}{2}}Bp_1(p+p_2)]`$ (35)
$`+{\displaystyle \frac{f(q_3^2)}{q_3^2m_\rho ^2+i\sqrt{q_3^2}\mathrm{\Gamma }_\rho (q_3^2)}}Bp(p_2p_1)\},`$
$`G^+={\displaystyle \frac{gf_am_K^{}^2m_K}{q^2m_{K_1}^2}}\{{\displaystyle \frac{f(q_2^2)}{q_2^2m_K^{}^2}}[{\displaystyle \frac{1}{2}}A(q_2^2)+{\displaystyle \frac{1}{2}}Bp_1(p+p_2)]`$
$`{\displaystyle \frac{f(q_3^2)}{q_3^2m_\rho ^2+i\sqrt{q_3^2}\mathrm{\Gamma }_\rho (q_3^2)}}A(q_3^2)\}.`$
In the same way the form factors of other two decay modes are obtained
$`F^{00}={\displaystyle \frac{1}{2}}{\displaystyle \frac{gf_am_K^{}^2m_K}{q^2m_{K_1}^2}}\{{\displaystyle \frac{f(q_1^2)}{q_1^2m_K^{}^2}}[{\displaystyle \frac{3}{2}}A(q_1^2)+{\displaystyle \frac{1}{2}}B(p_2p+p_2p_1)]`$ (39)
$`+{\displaystyle \frac{f(q_2^2)}{q_2^2m_K^{}^2}}[{\displaystyle \frac{3}{2}}A(q_2^2)+{\displaystyle \frac{1}{2}}B(p_1p+p_1p_2)]\},`$
$`G^{00}={\displaystyle \frac{1}{2}}{\displaystyle \frac{gf_am_K^{}^2m_K}{q^2m_{K_1}^2}}\{{\displaystyle \frac{f(q_1^2)}{q_1^2m_K^{}^2}}[{\displaystyle \frac{1}{2}}A(q_1^2){\displaystyle \frac{1}{2}}B(p_2p+p_2p_1)]`$
$`+{\displaystyle \frac{f(q_2^2)}{q_2^2m_K^{}^2}}[{\displaystyle \frac{1}{2}}A(q_2^2)+{\displaystyle \frac{1}{2}}B(p_1p+p_1p_2)]\},`$
$`F^{+0}={\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \frac{gf_am_K^{}^2m_K}{q^2m_{K_1}^2}}\{{\displaystyle \frac{f(q_1^2)}{q_1^2m_K^{}^2}}[{\displaystyle \frac{3}{2}}A(q_1^2)+{\displaystyle \frac{1}{2}}B(p_2p+p_2p_1)]`$
$`{\displaystyle \frac{f(q_2^2)}{q_2^2m_K^{}^2}}[{\displaystyle \frac{3}{2}}A(q_2^2)+{\displaystyle \frac{1}{2}}B(p_1p+p_1p_2)]`$
$`+{\displaystyle \frac{2f(q_3^2)}{q_3^2m_\rho ^2+i\sqrt{q_3^2}\mathrm{\Gamma }_\rho (q_3^2)}}Bp(p_1p_2)\},`$
$`G^{+0}={\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \frac{gf_aM_K^{}^2m_K}{q^2m_{K_1}^2}}\{{\displaystyle \frac{f(q_1^2)}{q_1^2m_K^{}^2}}[{\displaystyle \frac{1}{2}}A(q_1^2){\displaystyle \frac{1}{2}}B(p_2p+p_2p_1)]`$
$`{\displaystyle \frac{f(q_2^2)}{q_2^2m_K^{}^2}}[{\displaystyle \frac{1}{2}}A(q_2^2)+{\displaystyle \frac{1}{2}}B(p_1p+p_1p_2)]+{\displaystyle \frac{2f(q_3^2)}{q_3^2m_\rho ^2+i\sqrt{q_3^2}\mathrm{\Gamma }_\rho (q_3^2)}}A(q_3^2)\}.`$
The isospin relations(3) between these form factors are satisfied.
The partial wave analysis of these form factors can be done. The decay channel $`\rho \pi \pi `$ contributes to the decay modes of $`\pi ^+\pi ^{}`$ and $`\pi ^+\pi ^0`$. The range of the variable $`q_3^2`$ is $`4m_\pi ^2<q_3^2<(m_Km_l)^2`$ in which the decay width $`\mathrm{\Gamma }_\rho (q_3^2)`$ is not zero. The form factors, $`A^+`$ and $`A^{+0}`$ are complex functions of $`q_3^2`$. The $`\rho \pi \pi `$ doesn’t contribute to $`\pi ^0\pi ^0`$ mode. Therefore, $`F^{00}`$ and $`G^{00}`$ are real. $`K_{l4}`$ are decays at low energies. s- and p- waves are major partial waves. The $`q_1^2`$ and $`q_2^2`$ variables are expressed as
$`q_1^2={\displaystyle \frac{1}{2}}(m_K^2+2m_\pi ^2+q^2q_3^2)+(1{\displaystyle \frac{4m_\pi ^2}{q_3^2}})^{\frac{1}{2}}Xcos\theta _\pi ,`$ (41)
$`q_2^2={\displaystyle \frac{1}{2}}(m_K^2+2m_\pi ^2+q^2q_3^2)(1{\displaystyle \frac{4m_\pi ^2}{q_3^2}})^{\frac{1}{2}}Xcos\theta _\pi ,`$
where $`X=\{\frac{1}{4}(m_K^2q^2q_3^2)^2q^2q_3^2\}^{\frac{1}{2}}`$ and $`\theta _\pi `$ is the angle between $`\stackrel{}{p}_1`$ and $`\stackrel{}{p}`$ in the rest frame of the two pions.
The s- and p- wave amplitudes are obtained from Eqs.(34-39)
1. $`F_s^+`$ is real. Only Fig.2(a) contributes to it. $`F_p^+`$ is a complex function of $`q_3^2`$ resulted by $`\rho \pi \pi `$. $`F_p^+`$ has a phase shift.
$$F^+=F_s^++|F_p^+|e^{i\delta _p^+}(1\frac{4m_\pi ^2}{q_3^2})^{\frac{1}{2}}\frac{X}{m_K^2}cos\theta _\pi .$$
(42)
2. $`G_s^+`$ is complex and has a phase shift. $`G_p^+`$ is real.
$$G^+=|G_s^+|e^{i\delta _s^+}+G_p^+(1\frac{4m_\pi ^2}{q_3^2})^{\frac{1}{2}}\frac{X}{m_K^2}cos\theta _\pi .$$
(43)
3. Both $`G_s^{00}`$ and $`G_p^{00}`$ are real.
$`F^{00}`$ $`=`$ $`F_s^{00},`$
$`G^{00}`$ $`=`$ $`|G_p^{00}|(1{\displaystyle \frac{4m_\pi ^2}{q_3^2}})^{\frac{1}{2}}{\displaystyle \frac{X}{m_K^2}}cos\theta _\pi .`$ (44)
4. The isospin of the two pions of the $`\pi ^+\pi ^0`$ mode is one. Because of Bose statistics $`F^{+0}`$ only has p-wave which is complex and has phase shift. $`G^{+0}`$ has s wave only. $`G_s^{+0}`$ is complex and it has phase shift.
$`F^{+0}`$ $`=`$ $`|F_p^{+0}|e^{i\delta _p^{+0}}(1{\displaystyle \frac{4m_\pi ^2}{q_3^2}})^{\frac{1}{2}}{\displaystyle \frac{X}{m_K^2}}cos\theta _\pi ,`$
$`G^{+0}`$ $`=`$ $`|G_s^{+0}|e^{\delta _s^{+0}}.`$ (45)
All the phase shifts are caused by the decay $`\rho \pi \pi `$ and functions of $`q^2`$ and $`q_3^2`$.
## 6 Decay rates
The decay rates of the three modes of $`K_{e4}`$ and $`K_{\mu 4}`$ are calculated. As mentioned above, all the form factors are derived in the chiral limit. Therefore, only the leading terms of the masses of kaon and pions are kept in the calculation of the decay rates.
Ignoring $`m_e`$, only the form factors F, G, and H contribute to the decay rates of $`K_{e4}`$. By using the formula of Ref. we obtain
$$\mathrm{\Gamma }(K^{}\pi ^+\pi ^{}e\nu )=2.06\times 10^{21}GeV,B=3.87\times 10^5.$$
$$\mathrm{\Gamma }(K^{}\pi ^0\pi ^0e\nu )=0.221\times 10^{21}GeV,B=0.42\times 10^5.$$
$$\mathrm{\Gamma }(K^{}\pi ^+\pi ^0e\nu )=3.24\times 10^{21}GeV,B=2.55\times 10^4.$$
The experimental data are
$$B(\pi ^+\pi ^{})=(3.91\pm 0.17)\times 10^5[10],$$
$$B(\pi ^0\pi ^0)=(2.54\pm 0.89)\times 10^5(10events)[11],$$
$$B(\pi ^{}\pi ^0)=(5.16\pm 0.20\pm 0.22)\times 10^5[12],$$
$$B(\pi ^{}\pi ^0)=(6.2\pm 2.0)\times 10^5[13],$$
$$B(\pi ^{}\pi ^0)<200\times 10^5[14].$$
Theoretical result of $`\pi ^+\pi ^{}`$ mode agrees well with the data.
The form factors of the vector current are determined by anomalous vertices. The numerical calculation shows that the contribution of the form factor H is only $`0.5\%`$ of the total decay rate of $`K^{}\pi ^+\pi ^{}e\nu `$. Therefore, the axial-vector current dominates the $`K_{l4}`$ decays.
As shown in Fig.2(a,b) there are two channels in $`K_{l4}`$ decays. Numerical calculation of $`K^{}\pi ^+\pi ^{}e\nu `$ shows that the contribution of $`\rho \pi \pi `$(Fig.2(b)) is twice of the process, $`K^{}K\pi `$, (Fig.2(a)). Only the process(Fig.2(a)) contribute to $`K^{}\pi ^0\pi ^0e\nu `$. Because of Bose statistics there is an additional factor of $`\frac{1}{2}`$ in the formula of the decay rate of this mode. Therefore, this theory predicts smaller decay rate for this decay mode. On the other hand, the numerical calculation shows that the process(Fig.2(b)) is the major contributor of the decay $`\overline{K^0}\pi ^+\pi ^0e\nu `$. The theory predicts a larger branching ratio for $`\overline{K}^0\pi ^+\pi ^0e\nu `$.
All the form factors contribute to $`K_{\mu 4}`$ decays. Eq.(33) shows that in the chiral limit PCAC predicts that the form factor R is determined by other two form factors, F and G. The branching ratio of $`K_{\mu 4}`$ provides a test on this prediction. The numerical results are
$$\mathrm{\Gamma }(K^{}\pi ^+\pi ^{}\mu \nu )=0.634\times 10^{21}GeV,B=1.19\times 10^5.$$
$$\mathrm{\Gamma }(K^{}\pi ^0\pi ^0\mu \nu )=0.673\times 10^{22}GeV,B=0.126\times 10^5.$$
$$\mathrm{\Gamma }(\overline{K}^0\pi ^+\pi ^0\mu \nu )=1.01\times 10^{21}GeV,B=0.793\times 10^4.$$
The experimental data is
$$B(K^{}\pi ^+\pi ^{}\mu \nu )=(1.4\pm 0.9)\times 10^5.$$
Theory agrees with the data well.
## 7 Conclusion
All the four form factors of $`K_{l4}`$ have been derived from an effective theory of large $`N_C`$ QCD in the chiral limit. It has been found that the contribution of the vector current is negligible and the axial-vector current is dominant in $`K_{l4}`$ decays. PCAC is revealed from the theory. In the chiral limit it has been predicted that the form factor R is determined by the form factors F and G. The prediction has been tested by $`K^{}\pi ^+\pi ^{}\mu \nu `$. Theory agrees with the data. The partial wave analysis has been done. Non-zero phase shifts originate in the decay $`\rho \pi \pi `$. The process $`K_1\rho K`$ and $`\rho \pi \pi `$(Fig.2(b)) plays important role in $`K_{l4}`$ decays. Because of this channel the theory predicts larger branching ratio for $`K^{}\pi ^+\pi ^{}e\nu `$ and $`\overline{K}^0\pi ^+\pi ^0e\nu `$. The former agrees well with the data. $`\rho `$ resonance doesn’t contribute to $`K^{}\pi ^0\pi ^0e\nu `$. Therefore, the branching ratio of this decay mode is predicted to be smaller.
This research was partially supported by DOE Grant No. DE-91ER75661.
Figure Captions
Fig.1 Feynman Diagrams of vector current
Fig.2 Feynman diagrams of axial-vector current
Fig.3 Phase shifts
Fig.4 Phase shifts Fig.5 Phase shifts
Fig.5 Phase shifts
Fig.6 Phase shifts
Fig.7 Form factors od $`\pi ^+\pi ^{}`$ mode
Fig.8 Form factors od $`\pi ^+\pi ^{}`$ mode
Fig.9 Form factors od $`\pi ^+\pi ^0`$ mode
Fig.10 Form factors od $`\pi ^+\pi ^0`$ mode
Fig.11 Form factors od $`\pi ^0\pi ^0`$ mode
Fig.12 Form factors od $`\pi ^0\pi ^0`$ mode |
warning/0001/astro-ph0001280.html | ar5iv | text | # Untitled Document
ANALYTICAL STUDY OF THERMONUCLEAR REACTION PROBABILITY INTEGRALS
M.A. Chaudhry
Department of Mathematical Sciences, King Fahd University of Petroleum and Minerals, Dhahran 31261, Saudi Arabia
H.J. Haubold
Outer Space Office, United Nations,Vienna International Centre, P.O. Box 500, 1400 Vienna, Austria
and
A.M. Mathai
Department of Mathematics and Statistics, McGill University, 805 Sherbrooke Street West,
Montreal, Quebec, Canada H3A 2K6
Abstract. An analytic study of the reaction probability integrals corresponding to the various forms of the slowly varying cross-section factor $`S(E)`$ is attempted. Exact expressions for reaction probability integrals are expressed in terms of the extended gamma functions.
1. Introduction
Nuclear reactions govern major aspects of the chemical evolution of galaxies and stars (Fowler, 1984, Morel et al., 1999). The analytic study of the reaction rates and reaction probability integrals was undertaken by Critchfield (1972), Anderson et al. (1994), Haubold and Mathai (1998), and Chaudhry (1999). A proper understanding of the nuclear reactions that are going on in hot cosmic plasma, and those in the laboratories as well, requires a sound theory of nuclear-reaction dynamics (Clayton, 1983, Bergstroem et al., 1999). The rate $`r_{ij}`$ of unlike reacting nuclei $`i`$ and $`j`$ in the case of nonrelativistic nuclear reactions taking place in nondegenerate environment is expressed as (Clayton, 1983, Lang, 1999)
$$r_{ij}=n_in_j\left(\frac{8}{\pi \mu }\right)^{1/2}\left(\frac{1}{kT}\right)^{3/2}_0^{\mathrm{}}E\sigma (E)e^{E/KT}𝑑E,$$
(1.1)
where $`n_i`$ and $`n_j`$ denote the particle number densities of the reacting nuclei $`i`$ and $`j`$, $`\mu ={\displaystyle \frac{m_im_j}{m_i+m_j}}`$ is the reduced mass of the reacting nuclei, $`T`$ is the temperature, $`k`$ is the Boltzmann constant, $`\sigma (E)`$ is the cross-section for the reaction under consideration, and $`v=\left(\frac{2E}{\mu }\right)^{1/2}`$ is the relative velocity (for reactions between like nuclei $`i`$ = $`j`$, $`n_in_j`$ has to be replaced by $`n^2`$). Thus (Clayton, 1983, Lang, 1999)
$$r_{ij}=n_in_j\lambda :=n_in_j\sigma v:=_0^{\mathrm{}}\sigma (E)v(E)\psi (E)𝑑E$$
(1.2)
is the definition of the reaction probability integral $`\lambda `$, that is, the probability per unit time that two nuclei, confined to a unit volume, will react with each other. The reaction probability is written in the significant form $`\sigma v`$ to indicate that it is an appropriate average of the product of the reaction cross section and relative velocity of the interacting nuclei. If the reacting mixture is in thermal equilibrium this quantity depends only on the temperature.
For nonresonant nuclear reactions between nuclei of charges $`z_i`$ and $`z_j`$ at low energies (below the Coulomb barrier), the reaction cross-section has the form (Clayton, 1983, Rowley and Merchant, 1991, Lang, 1999)
$$\sigma (E)=\frac{S(E)}{E}e^{2\pi \eta (E)}$$
(1.3)
where
$$\eta (E)=\left(\frac{\mu }{2}\right)^{1/2}\frac{z_iz_je^2}{\mathrm{}E^{1/2}}$$
(1.4)
is the Sommerfeld parameter, $`\mathrm{}`$ is Planck’s quantum of action, and $`e`$ is the quantum of electric charge. Eq. (1.3) allows the extrapolation of measured reaction cross sections down to astrophysical energies by introducing the S-factor.
It is to be noted that $`S(E)`$, a residual function of energy, represents intrinsically nuclear parts of the probability for the occurrence of a nuclear reaction (Clayton, 1983). It is often found to be constant or a slowly varying function over a limited energy range when the interaction energy of the pair of nuclei is not nearly equal to an energy at which the two nuclei resonate in a quasi-stationary state.
The normalized energy distribution, representing the isotropic velocity distribution of the reacting nuclei, is given by (Clayton, 1983, Lang, 1999)
$$\psi (E)dE=\frac{2}{\sqrt{\pi }}\frac{E}{kT}\mathrm{exp}\left(\frac{E}{kT}\right)\frac{dE}{(kTE)^{1/2}}$$
(1.5)
(for reaction rates with anisotropic distributions see Imshennik, 1990). The substitution of $`\psi (E)dE`$ in (1.2) yields
$$\lambda =\left(\frac{8}{\mu \pi }\right)^{1/2}\frac{1}{(kT)^{3/2}}_0^{\mathrm{}}S(E)\mathrm{exp}\left(\frac{E}{kT}\frac{b}{E^{1/2}}\right)𝑑E.$$
(1.6)
where $`b=\left(\frac{\mu }{2}\right)^{1/2}\frac{z_iz_je^2}{\mathrm{}}.`$ If the residual function $`S(E)`$ is considered to be constant, $`S_0`$, the corresponding reaction probability integral is given by
$$\lambda =\left(\frac{8}{\mu \pi }\right)^{1/2}\frac{S_0}{(kT)^{3/2}}_0^{\mathrm{}}\mathrm{exp}\left(\frac{E}{kT}\frac{b}{E^{1/2}}\right)𝑑E.$$
(1.7)
In effect, the reaction probability is governed by the average of the Gamow penetration factor over the Maxwell-Boltzmann distribution. Anderson et al. (1994) considered for $`S(E)`$ the Maclaurin series expansion
$$S(E)=S(0)+\frac{dS(0)}{dE}E+\frac{1}{2}\frac{d^2S(0)}{dE^2}E^2$$
(1.8)
and solved (1.6) in terms of $`G`$\- and $`H`$-functions (Mathai and Saxena, 1973, 1978).
It is to be noted that the residual function $`S(E)`$ may not admit the representation (1.8) if the thermonuclear fusion plasma is not in the thermodynamic equilibrium, that is, the case in which there is a depletion or cut-off of the high energy tail of the Maxwell-Boltzman distribution in (1.5) (Lapenta and Quarati, 1993). Some of these cases have been considered in Anderson et al. (1994) and Haubold and Mathai (1998). In this paper we consider representations for the residual function given by
$`S_1(E)`$ $`=`$ $`S_0\delta (EE_0),`$ (1.9)
$`S_2(E)`$ $`=`$ $`S_0\left({\displaystyle \frac{E}{E_0}}\right)^{\alpha 1},`$ (1.10)
$`S_3(E)`$ $`=`$ $`S_0\left({\displaystyle \frac{E}{E_0}}\right)^{\alpha 1}H(EE_1),`$ (1.11)
$`S_4(E)`$ $`=`$ $`S_0\left({\displaystyle \frac{E}{E_0}}\right)^{\alpha 1}H(E_1E),`$ (1.12)
$`S_5(E)`$ $`=`$ $`S_0\left({\displaystyle \frac{E}{E_0}}\right)^{\alpha 1}\mathrm{exp}(C(EE_1)),`$ (1.13)
where
$$H(EE_1):=\{\begin{array}{cc}1,\hfill & \text{ if }E>E_1,\hfill \\ 0,\hfill & \text{ if }E<E_1\hfill \end{array}$$
(1.14)
is the unit step function and
$$\delta (EE_1):=\frac{d}{dE}\left(H(EE_1)\right)$$
(1.15)
is the Dirac delta function (Jennings and Karataglidid, 1998).
We discuss the analytic representations of the corresponding reaction probability integrals. The integral corresponding to (1.9) may be called instantaneous reaction probability integral.
2. The Astrophysical Thermonuclear Functions
The theoretical and experimental verification of nuclear cross-sections leads to the derivation of the closed-form representation of the thermonuclear reaction rates. These rates are expressed in terms of the four astrophysical thermonuclear functions (given in the notation chosen by Anderson et al., 1994)
$`I_1(z,\nu ):={\displaystyle _0^{\mathrm{}}}y^\nu \mathrm{exp}(y{\displaystyle \frac{z}{\sqrt{y}}})𝑑y,`$ (2.1)
$`I_2(z,d,\nu ):={\displaystyle _0^d}y^\nu \mathrm{exp}(y{\displaystyle \frac{z}{\sqrt{y}}})𝑑y,`$ (2.2)
$`I_3(z,t,\nu ):={\displaystyle _0^{\mathrm{}}}y^\nu \mathrm{exp}(y{\displaystyle \frac{z}{\sqrt{y+t}}})𝑑y,`$ (2.3)
$`I_4(z,\delta ,b,\nu ):={\displaystyle _0^{\mathrm{}}}y^\nu \mathrm{exp}(yby^\delta {\displaystyle \frac{z}{\sqrt{y}}})𝑑y,`$ (2.4)
where $`y=\frac{E}{kT}=\frac{\mu v^2}{2kT}`$ relates $`E`$ or $`v`$, respectively, to the mean thermal velocity and $`z=2\pi \left(\frac{\mu }{2kT}\right)^{1/2}\frac{z_iz_je^2}{\mathrm{}}=2\pi \left(\frac{\mu c^2}{2kT}\right)^{1/2}\alpha z_iz_j`$, were the velocity of light was introduced to make the dimension more apparent and to show the dependence on Sommerfeld’s fine structure constant $`\alpha `$. The closed-form representations of these integral functions, in terms of G- and H-functions, asymptotic values and, numerical results are discussed in Anderson et al. (1994) and Chaudhry (1999). Accounts about the developments of Meijer’s G-function and of Fox’s H-function are given in Mathai and Saxena (1973) and Mathai and Saxena (1978), respectively. Originally the investigations on these generalized hypergeometric functions were confined to theoretical results such as their analytical properties, integral representations, and asymptotic expansions, and then to the study of symmetric Fourier kernels, the solution of certain functional equations and other mathematical topics. Later, both the G-function and the H-function have been also used in the fields of statistical and astrophysical sciences, including extensive studies of the thermonuclear functions in eqs. (2.1)-(2.4) (Mathai, 1993).
3. The Extended Gamma Functions
The study of the astrophysical thermonuclear functions led to the development of a new class of special functions (Chaudhry and Zubair, 1998). In particular, we define the extended gamma functions by
$$\mathrm{\Gamma }(\alpha ,x;b;\beta ):=_x^{\mathrm{}}t^{\alpha 1}e^{tb/t^\beta }𝑑t,$$
(3.1)
and
$$\gamma (\alpha ,x;b;\beta ):=_0^xt^{\alpha 1}e^{tb/t^\beta }𝑑t.$$
(3.2)
It is to be noted that the functions (3.1) and (3.2) are special cases of the general class of extended gamma functions introduced in Chaudhry and Zubair (1998).
In fact we have
$`\mathrm{\Gamma }(\alpha ,x;b;\beta )=\mathrm{\Gamma }_{0,2}^{2,0}\left[(b,x)|\begin{array}{cc},& \\ ((0,1),& (\alpha ,\beta ))\end{array}\right]`$ (3.5)
$`:={\displaystyle \frac{1}{2\pi }}{\displaystyle _{ci\mathrm{}}^{c+i\mathrm{}}}\mathrm{\Gamma }(s)\mathrm{\Gamma }(\alpha +\beta s,x)b^s𝑑s,`$ (3.6)
and
$`\gamma (\alpha ,x;b;\beta )=\gamma _{0,2}^{2,0}\left[(b,x)|\begin{array}{cc},& \\ ((0,1),& (\alpha ,\beta ))\end{array}\right]`$ (3.9)
$`:={\displaystyle \frac{1}{2\pi }}{\displaystyle _{ci\mathrm{}}^{c+i\mathrm{}}}\gamma (s)\mathrm{\Gamma }(\alpha +\beta s,x)b^s𝑑s.`$ (3.10)
These functions satisfy the decomposition formula (Chaudhry and Zubair, 1998).
$$\gamma (\alpha ,x;b,\beta )+\mathrm{\Gamma }(\alpha ,x;b;\beta )=H_{0,2}^{2,0}\left[(b,x)|\begin{array}{cc},& \\ (0,1),& (\alpha ,\beta ))\end{array}\right].$$
(3.11)
The astrophysical thermonuclear functions (2.1) – (2.4) are special cases of the extended gamma functions (3.1) and (3.2) (Chaudhry, 1999)
$`I_1(z,\nu )=\mathrm{\Gamma }(\nu +1,0;z;{\displaystyle \frac{1}{2}})`$ (3.12)
$`I_2(z,d,\nu )=\gamma (\nu +1,d;z;{\displaystyle \frac{1}{2}})`$ (3.13)
$`I_3(z,t,\nu )=e^t{\displaystyle \underset{r=0}{\overset{\nu }{}}}\left(\begin{array}{c}\nu \\ r\end{array}\right)(t)^{\nu t}\mathrm{\Gamma }(\nu +1,t;z;{\displaystyle \frac{1}{2}})`$ (3.16)
$`I_4(z,\delta ,b,\nu )={\displaystyle \underset{r=0}{\overset{\nu }{}}}{\displaystyle \frac{(b)^r}{r!}}\mathrm{\Gamma }(\nu +r\delta +1,0;z;{\displaystyle \frac{1}{2}}).`$ (3.17)
It is straightforward to note that the transformation theorem
$$\mathrm{\Gamma }(\alpha ,x;b;\beta )=\frac{1}{\beta }\mathrm{\Gamma }_{0,2}^{2,0}\left[(b^{1/\beta },x)|\begin{array}{cc}& \\ ((0,\frac{1}{\beta }),& (\alpha ,1))\end{array}\right],(x0,b0,\beta 0),$$
(3.18)
for the extended gamma function reveals as a special case that
$`\mathrm{\Gamma }(\alpha ,0;b;{\displaystyle \frac{1}{n}})=H_{0,2}^{2,0}\left[b|\begin{array}{cc},& \\ (0,1),& (\alpha ,n)\end{array}\right]=(2\pi )^{(1n)/2}\sqrt{n}`$ (3.21)
$`\times G_{0,n+1}^{n+1,0}\left[\left({\displaystyle \frac{b}{n}}\right)^n|\begin{array}{ccc},& ,\mathrm{},& \\ 0,& {\displaystyle \frac{1}{n}},{\displaystyle \frac{2}{n}},\mathrm{},{\displaystyle \frac{n1}{n}},& \alpha \end{array}\right].`$ (3.24)
The relation (3.11) yields directly the closed form representation (Anderson et al., 1994)
$`I_1(z,\nu )`$ $`=`$ $`\mathrm{\Gamma }(\nu +1,0;z;{\displaystyle \frac{1}{2}})`$ (3.25)
$`=`$ $`\pi ^{1/2}G_{0,3}^{3,0}\left[{\displaystyle \frac{z^2}{4}}|0,{\displaystyle \frac{1}{2}},1+\nu \right].`$
The asymptotic representation of the extended gamma function $`\mathrm{\Gamma }(\alpha ,0;b;\beta )`$ for small and large values of $`b`$ are given as follows (Chaudhry and Zubair, 1998)
$$\mathrm{\Gamma }(\alpha ,0;b;\beta )\{\begin{array}{cc}\mathrm{\Gamma }(\alpha )+\frac{1}{\beta }\mathrm{\Gamma }\left(\frac{\alpha }{\beta }\right)b^{\alpha /\beta }+o(b),\hfill & \text{ for small }b,\alpha 0,\hfill \\ \frac{1}{\beta }\mathrm{ln}b,\hfill & \text{ for small }b,\alpha =0,\hfill \end{array}$$
(3.26)
and
$`\mathrm{\Gamma }(\alpha ,0;b;\beta ){\displaystyle \frac{1}{\beta }}\left({\displaystyle \frac{2\pi \beta }{1+\beta }}\right)^{1/2}\beta ^{(2\alpha +\beta )/2(1+\beta )}b^{(2\alpha 1)/2(\beta +1)}`$
$`\times \mathrm{exp}[(1+\beta )^{\beta /(1+\beta )}b^{1/(1+\beta )}],\text{ for large }b.`$ (3.27)
The representations (3.13) and (3.14) are useful in finding the asymptotic values of the thermonuclear reaction probability integrals.
4. Thermonuclear Reaction Probability Integrals
In this section we study the closed form representations of the thermonuclear reaction probability integrals corresponding to the various forms (1.9) – (1.13) of the residual function $`S(E)`$.
(4.1) Instantaneous Residual Function, $`S(E)=S_0\delta (EE_0)`$. The substitution of this value of $`S(E)`$ in (1.6) yields an elementary representation
$$\lambda _1=\left(\frac{8}{\mu \pi }\right)^{1/2}\frac{S_0}{(kT)^{3/2}}\mathrm{exp}\left(\frac{E_0}{kT}\frac{b}{\sqrt{E_0}}\right)$$
(4.1)
of the corresponding thermonuclear probability integral.
(4.b) Power Type Residual Function, $`S(E)=S_0(E/E_0)^{\alpha 1}`$. The substitution of the power type residual function in (1.6) yields
$$\lambda _2=\left(\frac{8}{\mu \pi }\right)^{1/2}\frac{S_0}{(kT)^{3/2}}_0^{\mathrm{}}(E/E_0)^{\alpha 1}\mathrm{exp}\left(\frac{E}{kT}\frac{b}{\sqrt{E}}\right)𝑑E.$$
(4.2)
The transformation $`t=E/kT`$ in (4.2) yields
$$\lambda _2=\left(\frac{8}{\mu \pi }\right)^{1/2}\frac{S_0}{\sqrt{E_0}}\left(\frac{kT}{E_0}\right)^{\alpha \frac{3}{2}}_0^{\mathrm{}}t^{\alpha 1}\mathrm{exp}\left(t\frac{b}{\sqrt{kT}\sqrt{t}}\right)𝑑t,$$
(4.3)
which is available from Haubold and Mathai (1998) and can be written in terms of the extended gamma function, i.e..
$$\lambda _2=\left(\frac{8}{\mu \pi }\right)^{1/2}\frac{S_0}{\sqrt{E_0}}\left(\frac{kT}{E_0}\right)^{\alpha \frac{3}{2}}\mathrm{\Gamma }(\alpha ,0;\frac{b}{\sqrt{kT}};\frac{1}{2}).$$
(4.4)
In view of the decomposition relation (3.5) we can simplify (4.4) to obtain
$$\lambda _2=\left(\frac{8}{\mu \pi }\right)^{1/2}\frac{S_0}{\sqrt{E_0}}\left(\frac{kT}{E_0}\right)^{\alpha \frac{3}{2}}H_{0,2}^{2,0}\left[\frac{b}{\sqrt{kT}}|\begin{array}{cc},& \\ (0,1),& (\alpha ,2)\end{array}\right]$$
(4.5)
that can further be simplified to give
$$\lambda _2=\left(\frac{8}{\mu \pi }\right)^{1/2}\frac{S_0}{\sqrt{E_0}}\left(\frac{kT}{E_0}\right)^{\alpha \frac{3}{2}}G_{0,3}^{3,0}\left[\frac{b^2}{kT}|0,\frac{1}{2},\alpha \right].$$
(4.6)
(4.c) Power Type Delayed Residual Function, $`S(E)=S_0(E/E_0)^{\alpha 1}H(EE_1)`$. The substitution of the power type delayed residual function in (1.6) gives
$$\lambda _3=\left(\frac{8}{\mu \pi }\right)^{1/2}\frac{S_0}{(kT)^{3/2}}_{E_1}^{\mathrm{}}(E/E_0)^{\alpha 1}\mathrm{exp}\left(\frac{E}{kT}\frac{b}{\sqrt{E}}\right)𝑑E.$$
(4.7)
The transformation $`t=E/kT`$ in (4.7) yields
$$\lambda _3=\left(\frac{8}{\mu \pi }\right)^{1/2}\frac{S_0}{\sqrt{E_0}}\left(\frac{kT}{E_0}\right)^{\alpha \frac{3}{2}}_{E_1/kT}^{\mathrm{}}t^{\alpha 1}\mathrm{exp}\left(t\frac{b}{\sqrt{kT}\sqrt{t}}\right)𝑑t,$$
(4.8)
or
$$\lambda _3=\left(\frac{8}{\mu \pi }\right)^{1/2}\frac{S_0}{\sqrt{E_0}}\left(\frac{kT}{E_0}\right)^{\alpha \frac{3}{2}}\mathrm{\Gamma }(\alpha ,\frac{E_1}{kT};\frac{b}{\sqrt{kT}};\frac{1}{2}).$$
(4.9)
The asymptotic representations of $`\lambda _3`$ for small and large values of $`b/\sqrt{kT}`$ can be found from (3.13) and (3.14). In fact for small value of $`b/\sqrt{kT}`$ and $`\alpha 0`$
$$\lambda _3\left(\frac{8}{\mu \pi }\right)^{1/2}\frac{S_0}{\sqrt{E_0}}\left(\frac{kT}{E_0}\right)^{\alpha \frac{3}{2}}\left\{\mathrm{\Gamma }(\alpha )+\frac{2\mathrm{\Gamma }(\alpha )b^2}{kT}\right\}$$
(4.10)
(4.d) Power Type Cut-off Residual Function, $`S(E)=S_0(E/E_0)^{\alpha 1}H(E_1E)`$. The substitution of the power type cut-off residual function in (1.6) gives
$$\lambda _4=\left(\frac{8}{\mu \pi }\right)^{1/2}\frac{S_0}{(kT)^{3/2}}_0^{E_1}(E/E_0)^{\alpha 1}\mathrm{exp}\left(\frac{E}{kT}\frac{b}{\sqrt{E}}\right)𝑑E$$
(4.11)
that can similarly be simplified in terms of the extended gamma function to give
$$\lambda _4=\left(\frac{8}{\mu \pi }\right)^{1/2}\frac{S_0}{\sqrt{E_0}}\left(\frac{kT}{E_0}\right)^{\alpha \frac{3}{2}}\gamma (\alpha ,\frac{E_1}{kT};\frac{b}{\sqrt{kT}};\frac{1}{2}).$$
(4.12)
(4.e) Exponential Type Residual Function, $`S(E)=S_0(E/E_0)^{\alpha 1}\mathrm{exp}(C(EE_1)),`$ $`\left(C>\frac{1}{kT}\right)`$. The substitution of the exponential type residual function in (1.6) yields
$$\lambda _5=\left(\frac{8}{\mu \pi }\right)^{1/2}\frac{S_0}{(kT)^{3/2}}\mathrm{exp}(CE_1)_0^{\mathrm{}}(E/E_0)^{\alpha 1}\left(\left(C+\frac{1}{kT}\right)E\frac{b}{\sqrt{E}}\right)𝑑E.$$
(4.13)
The transformation
$$t=\left(C+\frac{1}{kT}\right)E,\left(C>\frac{1}{kT}\right)$$
in (4.13) yields
$$\lambda _5=\left(\frac{8}{\mu \pi }\right)^{1/2}\frac{S_0}{(kT)^{3/2}}\left(\frac{CkT+1}{kTE_0}\right)^{\alpha 1}\mathrm{exp}(CE_1)_0^{\mathrm{}}t^{\alpha 1}\mathrm{exp}\left(t\frac{b}{\sqrt{\frac{kT}{CkT+1}}\sqrt{t}}\right)𝑑t.$$
(4.14)
The integral in (4.14) is solvable in terms of extended gamma function to give
$$\lambda _5=\left(\frac{8}{\mu \pi }\right)^{1/2}\frac{S_0}{(kT)^{3/2}}\left(\frac{CkT+1}{kTE_0}\right)^{\alpha 1}\mathrm{exp}(CE_1)\mathrm{\Gamma }(\alpha ,0;b\sqrt{\frac{kT}{CkT+1}};\frac{1}{2}).$$
(4.15)
The equation (4.15) can be simplified further in terms of the H-function to give
$$\lambda _5=\left(\frac{8}{\mu \pi }\right)^{1/2}\frac{S_0}{(kT)^{3/2}}\left(\frac{CkT+1}{kTE_0}\right)^{\alpha 1}\mathrm{exp}(CE_1)H_{0,2}^{2,0}\left[b\sqrt{\frac{kT}{CkT+1}}|\begin{array}{cc},& \\ (0,1),& (\alpha ,2)\end{array}\right].$$
(4.16)
5. Conclusion
Nuclear reactions govern major aspects of the chemical evolution of galaxies and stars. The analytic study of the reaction rates and reaction probability integrals is important in astrophysics. The development of the new class of extended gamma functions by Chaudhry and Zubair has facilitated in convenient notations for these analytic representations. We have considered various forms of the residual function occurring in the reaction probability integrals and have written them analytically in terms of the extended gamma functions.
Acknowledgment. The first author is indebted to King Fahd University of Petroleum and Minerals, Dhahran, Saudi Arabia for excellent research facilities.
References
Anderson, W.J., Haubold, H.J., and Mathai, A.M.: 1994, Astrophysical thermonuclear functions, Astrophys. Space Sci. 214, 49–70.
Bergstroem, L., Iguri, S., and Rubinstein, H.: 1999, Constraints on the variation of the fine structure constant from big bang nucleosynthesis, Phys. Rev. D60, 045005-1 - 045005-9.
Chaudhry, M.A. and Zubair, S.M.: 1998, Extended incomplete gamma functions with applications, Journal of the London Mathematical Society (submitted).
Chaudhry, M.A.: 1999, Transformation of the extended gamma function $`\mathrm{\Gamma }_{0,2}^{2,0}[(B,X)]`$ with applications to astrophysical thermonuclear functions, Astrophys. Space Sci. 262, 263–270.
Clayton, D.D.: 1983, Principles of Stellar Evolution and Nucleosynthesis, Second Edition, The University of Chicago Press, Chicago and London.
Critchfield, C.L: 1972, in: F. Reines (ed.), Cosmology, Fusion and Other Matters, George Gamow Memorial Volume, Colorado, Colorado: Associated University Press 1972.
Fowler, W.A.: 1984, Experimental and theoretical nuclear astrophysics: the quest for the origin of the elements, Rev. Mod. Phys. 56, 149-179.
Haubold, H.J. and Mathai, A.M.: 1998, On thermonuclear reaction rates, Astrophys. Space Sci. 258, 185–199.
Imshennik, V.S.: 1990, Rate of nuclear reactions for an anisotropic distribution of interacting particles, Sov. J. Plasma Phys. 16, 379-383.
Jennings, B.K. and Karataglidid, S.: 1998, $`S_{eff}(E)`$ and the $`{}_{}{}^{7}Be(p,\gamma )^8B`$ reaction,
http://xxx.lanl.gov/abs/nucl-th/9807007.
Lang, K.R.: 1999, Astrophysical Formulae, Vol.I: Radiation, Gas Processes and High Energy Astrophysics, Third Enlarged and Revised Edition, Springer-Verlag, Berlin Heidelberg New York.
Lapenta, G. and Quarati, P.: 1993, Analysis of non-Maxwellian fusion reaction rates with electron screening, Z. Phys. A346, 243-250.
Mathai, A.M.: 1993, A Handbook of Generalized Spacial Functions for Statistical and Physical Sciences, Clarendon Press, Oxford.
Mathai, A.M. and Saxena, R.K.: 1973, Generalized Hypergeometric Functions with Applications in Statistics and Physical Sciences, Springer-Verlag, Lecture Notes in Mathematics Vol. 348, Berlin Heidelberg New York.
Mathai, A.M. and Saxena, R.K.: 1978, The H-function with Applications in Statistics and Other Disciplines, John Wiley and Sons, New Delhi.
Morel, P., Pichon, B., Provost, J., and Berthomieu, G.: Solar models and NACRE thermonuclear reaction rates, Astron. Astrophys. 350, 275-285.
Rowley, N. and Merchant, A.C.: 1991, Barrier penetration at astrophysical energies, Astrophys. J. 381, 591-596;
(http://pntpm.ulb.ac.be/Nacre/nacre.htm). |
warning/0001/hep-ph0001285.html | ar5iv | text | # COSMOLOGICAL EXPANSION IN THE RANDALL-SUNDRUM WARPED COMPACTIFICATION
## 1 Large Versus Small Extra Dimensions
In the last few years there has been a revival of interest in the idea of extra dimensions, first proposed by Kaluza and Klein. The new realization of Arkani-Hamed, Dvali and Dimopoulos (ADD) was that the extra dimensions could be macroscopically large if one assumed that our ($`3+1`$) dimensional universe is a slice (a 3-brane) of the higher dimensional bulk.$`^\mathrm{?}`$ The particles of the standard model should be restricted to the brane so that no light Kaluza-Klein (KK) excitations exist, which otherwise would have already been seen. Gravity, however, can propagate in the extra bulk dimensions (otherwise they would have no observable consequences whatsoever). The effect of the extra dimensions can only be seen on distance scales less than the order of their size. With $`N`$ extra compact dimensions of size $`R`$, Newton’s gravitational force law for two masses $`m_1`$ and $`m_2`$, separated by a distance $`r`$, is modified to
$$F=\frac{\mathrm{\Gamma }(\frac{3+N}{2})}{4\pi ^{(3+N)/2}}\left(\frac{m_1m_2}{M^{2+N}r^N}\right),rR,$$
(1)
where $`M`$ is the new quantum gravity scale appearing in the Einstein-Hilbert action for gravity in $`4+N`$ dimensions,
$$S=\frac{1}{2}M^{2+N}d^4xd^Ny\sqrt{|g|},$$
(2)
and $`y_I`$ parametrize the extra dimensions. At larger separations, the gravitational flux is no longer diluted by spreading out in the bulk, so the force reverts to its usual form,
$$F=\frac{1}{8\pi }\left(\frac{m_1m_2}{M_p^2r^2}\right),rR,$$
(3)
involving the 4-D Planck mass $`M_p`$. Deviations of (3) from its $`1/r^2`$ form have only been tested at separations greater than a millimeter or so, showing that $`R`$ could be as large as 1 mm. This is obviously far bigger than the limit which exists if the standard model particles are allowed to propagate in the bulk, $`R\text{ }\stackrel{<}{}\text{ }10^3`$ fm.
The relationship between the $`4`$-D and the $`(4+N)`$-D gravity scales can easily be deduced by requiring that the action (2) reduce to the usual one after integrating over the extra dimensions. Let us suppose that the line element in the full spacetime has the simple form
$$ds^2=a^2(y_I)\left(dt^2\underset{i=1}{\overset{3}{}}dx_idx_i\right)b^2(y_I)\underset{I=1}{\overset{N}{}}dy_Idy_I.$$
(4)
Then, if the brane is located at $`y_I=0`$, the relationship is
$$M_p^2=M^{N+2}\left(\frac{a(y_I)}{a(0)}\right)^2b^N(y_I)d^Ny.$$
(5)
In ADD, the geometry is assumed to be factorizable, meaning that $`a`$ does not depend upon $`y_I`$; hence the integral just gives the volume of the extra dimensions: $`M_p^2=M^{N+2}V_N`$. Since $`V_N`$ can be quite large (mm<sup>3</sup>), this has the interesting consequence that $`M`$ could be at the TeV scale, yet be consistent with the much higher scale of $`M_p`$. This opens the mind-boggling possibility that all new particle physics, including quantum gravity, could become accessible at the LHC. Moreover we have a partial explanation of the weak scale hierarchy problem, the question of why $`M_p`$ is 16 orders of magnitude larger than the $`W`$ boson mass. It is not really a solution because one is left with the annoying question of why $`R`$ is so much larger than the natural scale, $`1/M`$. The exact size depends on the number of extra dimensions. If $`N=1`$ it is not possible to obtain $`M1`$ TeV because $`R`$ is too large; demanding that $`R=1`$ mm gives $`M10^8`$ GeV. But for $`N=2`$, the TeV scale emerges just as the experimental bound on $`R`$ is saturated, and for higher dimensions it can be attained with smaller sizes, $`R100`$ fm in the case of $`N=6`$.
The experimental constraints on large extra dimensions come from the effects of the KK excitations of the graviton, which can be very light, $`m_n=n/R`$ for integer $`n`$. The only thing which saves these particles from being easily discovered is their weak interactions; like the ordinary graviton, their couplings are suppressed by $`1/M_p^2`$ (as opposed to $`1/M^2`$). Consequently the bounds from accelerator physics are rather weak: $`M\text{ }\stackrel{>}{}\text{ }`$ several TeV.$`^\mathrm{?}`$ Astrophysics gives better constraints, at least for $`N=1`$ and 2. One such bound comes from requiring that supernova 1987A not cool too quickly by graviton emission,$`^\mathrm{?}`$ giving $`M\text{ }\stackrel{>}{}\text{ }100`$ TeV for $`N=1`$ and $`M\text{ }\stackrel{>}{}\text{ }5`$ TeV for $`N=2`$. In the early universe, KK gravitons can be produced by thermal processes, and decay slowly into photons that would distort the cosmic gamma ray background unless $`M`$ obeys bounds similar to the supernova ones.$`^\mathrm{?}`$
Actually, the last-mentioned bound is quite generous toward the ADD scenario because it assumes that, by some miracle, the universe is already free from primordially produced KK gravitons at temperatures near 1 MeV–a necessary condition since otherwise the gamma rays produced by their decays would destroy deuterium and consequently the successful predictions of big bang nucleosynthesis. It is quite difficult to justify this assumption. Benakli and Davidson showed that if $`M=1`$ TeV, the reheat temperature after inflation would have to be no greater than 0.1 GeV, for $`N=6`$; for smaller $`N`$ the bound is even more stringent.$`^\mathrm{?}`$ This is well below what is needed for electroweak baryogenesis, which is generally considered to be the lowest temperature mechanism available. Therefore baryogenesis presents a major challenge to the ADD idea.
Randall and Sundrum (RS) have suggested another way of solving the hierarchy problem with an extra dimension,$`^\mathrm{?}`$ which avoids the difficulties encountered by ADD. They considered just a single extra dimension, compactified on an orbifold $`S_1/Z_2`$, a circle modded by $`Z_2`$. The coordinate is in the range $`y[1,1]`$, with the endpoints identified and with $`yy`$ being the orbifold symmetry. One places a 3-brane at each of the orbifold fixed points, $`y=0`$ and $`y=1`$. They have equal and opposite tensions, $`\pm \sigma `$ (tension is the 4-D energy density, which has the same form as a 4-D cosmological constant). In addition there is a 5-D cosmological constant in the bulk, $`\mathrm{\Lambda }`$. The stress-energy tensor is therefore
$$T_{\mu \nu }=(g_{\mu \nu }n_\mu n_\nu )\sigma \left(\delta (y)\delta (y1)\right)/b+\mathrm{\Lambda }g_{\mu \nu },$$
(6)
where $`n_\mu `$ is the normal to the branes (hence the brane tensions make no contribution to $`T_{yy}`$). A static solution to the 5-D Einstein equations exists if
$$\mathrm{\Lambda }=\frac{\sigma ^2}{6M^3},$$
(7)
and it has the form of eq. (4) with
$$a(y)=e^{kby};k=|\mathrm{\Lambda }/\sigma |$$
(8)
and $`b`$, the size of the extra dimension, being undetermined. Using eq. (5), one finds that $`M_p`$ is related to the 5-D gravity scale by
$$M_p^2=\frac{M^3}{k}(1e^{2kb}).$$
(9)
The dramatic consequence of this solution is that if one considers the Lagrangian for a particle confined to the brane at $`y=1`$, it takes its canonical form only after a Weyl rescaling of the field by the “warp factor” $`a(1)=e^{kb}`$. If the Lagrangian originally had $`M_p`$ as the mass scale for the particle, it becomes rescaled by $`\mathrm{exp}(kb)`$. One can take all the parameters $`M`$, $`\mathrm{\Lambda }`$, $`\sigma `$ and $`k`$ to be of order $`M_p`$ to the appropriate power; then if $`b36/k36/M_p`$, one obtains TeV scale masses on the $`y=1`$ brane (henceforth called the TeV brane). Clearly $`bk36`$ is a much more moderate hierarchy than $`M_p/M10^{16}`$, so this constitutes an attractive possible explanation of the weak scale. Furthermore the extra dimension is still small, so the KK gravitons can be sufficiently heavy to present no difficulties in the early universe.
This solution to the hierarchy problem requires that we are living on the negative tension brane (taken to be at $`y=1`$). The positive tension brane at $`y=0`$ has no such suppression of its masses, so it is referred to as the Planck brane, and must constitute a kind of hidden sector.
## 2 Effect of Extra Dimensions on Cosmological Expansion
For a while it appeared that cosmology could provide an interesting constraint on large extra dimensions. Binétruy, Deffayet and Langlois (BDL) considered the cosmological expansion of 3-brane universes in a 5-D bulk and found solutions in which the Hubble expansion rate in the brane was related to the energy density $`\rho `$ on the brane by $`^\mathrm{?}`$
$$H=\frac{\dot{a}}{a}=\frac{\rho }{6M^3},$$
(10)
in contrast to the usual Friedmann equation, $`H\sqrt{\rho }`$. Although other authors had found inflationary solutions with this property,$`^\mathrm{?}`$ BDL were the first to point out that it would be a problem for later cosmology. Especially, such a modification to the expansion rate would probably drastically alter the predictions of big bang nucleosynthesis.
It is not difficult to see from the 5-D Einstein equations, $`G_{\mu \nu }=M^3T_{\mu \nu }`$, why one gets the unusual dependence of $`H`$ on $`\rho `$. Consider the $`00`$ component,
$$\frac{\dot{a}}{a}\left(\frac{\dot{a}}{a}+\frac{\dot{b}}{b}\right)=\frac{a^2}{b^2}\left(\frac{a^{\prime \prime }}{a}+\frac{a^{}}{a}\left(\frac{a^{}}{a}\frac{b^{}}{b}\right)\right)+\frac{1}{3M^3}T_{00}.$$
(11)
To obtain the delta functions in $`T_{00}`$, eq. (6), $`a^{}(y)`$ must be discontinuous at both branes, and the discontinuity is proportional to the total energy density on the branes. Moreover the orbifold symmetry (as well as common sense) requires that $`a(y)`$ be symmetric about either brane, so that $`a^{}(1ϵ)=a^{}(1+ϵ)`$. This implies that $`a^{}(y)`$ itself is linearly proportional to the brane tension $`\sigma `$. Now consider the $`yy`$ component, which has no delta functions:
$$\left(\frac{\dot{a}}{a}\right)^2+\frac{\ddot{a}}{a}=2\frac{a^2}{b^2}\left(\frac{a^{}}{a}\right)^2\frac{1}{3M^3}T_{yy}$$
(12)
Recalling that $`H=\frac{\dot{a}}{a}`$, clearly $`H^2`$ will get contributions proportional to $`(a^{}/a)^2\sigma ^2`$ as well as $`\mathrm{\Lambda }`$. In fact, if we allow for a cosmological constant in the bulk and extra energy densities $`\rho _P`$ and $`\rho _T`$, in addition to the respective tensions $`\sigma `$ and $`\sigma `$ on the Planck and TeV branes, the complete expression becomes<sup>a</sup><sup>a</sup>athe factor $`e^{2kb}`$ was first pointed out by ref. $`^\mathrm{?}`$
$$H^2=\frac{(\sigma +\rho _P)^2}{36M^6}+\frac{\mathrm{\Lambda }}{6M^3}=\frac{(\sigma +e^{2kb}\rho _T)^2}{36M^6}+\frac{\mathrm{\Lambda }}{6M^3}$$
(13)
It was noticed $`^{\mathrm{?},\mathrm{?}}`$ that by tuning $`\sigma `$ to cancel the contribution from $`\mathrm{\Lambda }`$ in the limit $`\rho _i=0`$, one obtains an expression for $`H`$ which at leading order in $`\rho `$ has the desired $`\sqrt{\rho }`$ form, plus small fractional corrections of order $`\rho /M^4`$. Not surprisingly, in retrospect, this tuning is precisely the same condition (7) required by RS to obtain their solution. (Any deviation from this condition results in an effective 4-D cosmological constant and therefore a nonstatic solution.) But at the time we first noticed this coincidence, it was striking to us, since we were unaware of RS and had thus come upon the condition (7) starting from a completely different motivation from that of RS.
However, all is not well with the cosmological solution leading to eq. (13). For one thing, the energy densities on the two branes are constrained, $`\rho _T=e^{2kb}\rho _P`$. Moreover, this implies that $`\rho _T<0`$, i.e., that the energy density of matter on the TeV brane is negative, a physically unacceptable situation. Thus, although cosmology appeared to be normal on the Planck brane (for densities $`\rho _PM^4`$), not so on the TeV brane, where the hierarchy problem is solved. This seemed to present a problem for the RS proposal.$`^{\mathrm{?},\mathrm{?}}`$
There were several attempts to solve this problem. In one it was observed that by decompactifying the orbifold,$`^\mathrm{?}`$ placing the TeV brane at the position required by the hierarchy problem ($`y36/kb`$), and giving it a tension between $`0`$ and $`\sigma /2`$, one could obtain the normal Hubble rate on the TeV brane with a positive energy density.$`^\mathrm{?}`$ However this solution involved simultaneous expansion of the bulk, which is unacceptable for late time cosmology because a growing $`b(t)`$ leads to a Planck mass which is increasing in time, according to eq. (5). Hence gravity would be getting weaker on the TeV brane, contrary to stringent constraints on the time variation of Newton’s constant. In another attempt it was pointed out that the normal Friedmann equation would ensue if the $`yy`$ component of $`T_{\mu \nu }`$ was allowed to have a rather complicated dependence on the bulk coordinate $`y`$, rather than being a constant ($`\mathrm{\Lambda }`$).$`^\mathrm{?}`$ The origin of such a dependence seemed obscure (but see below).
It was recently shown that both of the cosmological problems of brane universe models—the artificial relation between the energy densities on the two branes, and the generically “wrong” form for the Friedmann equation—can be solved by introducing a mechanism for insuring that the size of the extra dimension remains fixed while the branes expand.$`^\mathrm{?}`$ Recall that this degree of freedom was completely undetermined in the RS model, meaning that it corresponds to a modulus, i.e., a massless field, in this context called the radion. This is problematic in itself, because it implies a fifth force, as in scalar-tensor theories of gravity, which in the present case has couplings suppressed by the TeV rather than the Planck scale.$`^\mathrm{?}`$ In the absence of a mechanism for stabilizing this modulus (see ref. $`^\mathrm{?}`$ for such a mechanism), it is not surprising that a fine-tuning between the brane energies as in eq. (13) should be needed to insure that the bulk does not expand along with the branes.$`^\mathrm{?}`$
Somewhat less obvious is the fact that the normal Friedmann equation also results if $`b`$ (the size of the compact dimension) is stabilized. Heuristically, this occurs because the bulk cosmological constant $`\mathrm{\Lambda }`$ is now replaced by a potential for $`b`$, $`V(b)`$. Since the $`yy`$ component of Einstein’s equation comes from the variation of the action with respect to $`b`$, a new term appears in $`T_{yy}`$,
$$T_{yy}=b^2(V(b)+bV^{}(b)).$$
(14)
Therefore the $`G_{yy}`$ equation, which we used in the argument above to obtain $`H\rho `$, is no longer available for fixing the magnitude of $`H`$; rather it determines $`b`$,$`^\mathrm{?}`$ because $`b`$ no longer sits at the bottom of the potential during a period of cosmological expansion, but is slightly shifted. Moreover ref. $`^\mathrm{?}`$ showed that the $`y`$-dependent stress-energy needed in $`^\mathrm{?}`$ to get the correct Friedmann equation automatically arises from the stabilization of the radion.
In $`^\mathrm{?}`$ it is shown that, if the radion is stabilized, then as long as the excess energy densities $`\rho _i`$ (above and beyond those needed in eq. (7) to get a static solution) are small compared to the cutoff scales ($`M_p`$ on the Planck brane and 1 TeV on the TeV brane), the two branes expand at approximately the same rate, given by
$$H^2=\frac{8\pi G}{3}\left(\rho _P+e^{4kb}\rho _T+_0^1𝑑ybe^{4kby}\rho _{\mathrm{bulk}}(y)\right).$$
(15)
This expression can be derived from the effective 4-D Lagrangian obtained by integrating over the extra dimension in the background of the RS metric. Notice the factor of $`e^{4kb}`$ multiplying $`\rho _T`$. This is precisely the same redshifting of mass scales that applies to all masses on the TeV brane. Thus $`\rho _T`$ represents the bare value of the energy density, presumably of order $`M_p^4`$, while $`e^{4kb}\rho _T`$ is the physically observed value. No such suppression occurs for matter living on the Planck brane. Therefore, if the expansion of the universe is to be dominated by matter on our (TeV) brane, it is necessary to demand that the Planck brane (and the bulk) be devoid of matter. Fortunately this does not seem to pose a major challenge: inflation will effectively empty out the Planck brane, as long as it harbors no nearly massless particles and the reheat temperature is significantly below the cutoff. Although it is tempting to suggest that $`\rho _P`$ could be the dark matter of the universe, it is difficult to see how it could be made sufficiently small, if it is not zero.
## 3 Outlook
The Randall-Sundrum proposal for solving the hierarchy problem with a small extra dimension looks compatible with most cosmological requirements. Unlike the ADD scenario of large extra dimensions, it does not suffer from the problem of light KK gravitinos wreaking havoc with nucleosynthesis and the cosmic gamma ray background. Furthermore it might have a plausible string theoretic origin,$`^\mathrm{?}`$ perhaps being related to the 5D-anti-deSitter space/conformal field theory correspondence and the holographic principle.$`^\mathrm{?}`$ There remain a few puzzles. One is the apparent possibility of a “dark radiation” term in the Friedmann equation,$`^\mathrm{?}`$
$$H^2=\frac{8\pi G}{3}\left(\text{usual terms}+\frac{𝒞}{a^4}\right)$$
(16)
which arises as an initial condition, due to the fact that the solutions to 5-D gravity have an additional constant of integration relative to 4-D. Does this oddity also disappear when the extra dimension is stabilized? Another problem is inflation, which typically requires the presence of an intermediate scale, $`M_i10^{13}`$ GeV, to get the right magnitude of density perturbations, $`\delta \rho /\rho M_i/M_p`$. No such intermediate scale exists if the TeV scale is the true cutoff on our brane. A third interesting question is how to generalize the RS scenario to higher dimensions, which is just beginning to be explored.$`^{\mathrm{?},\mathrm{?}}`$
## References |
warning/0001/cond-mat0001131.html | ar5iv | text | # Magnetic relaxation in hard type-II superconductors
## 1 Introduction
In a hard type-II superconductor, flux-creep over energy barrier at a finite $`T`$ leads to magnetic relaxation over long time scales. The flux creep, occurring due to thermally activated hopping of the flux lines, tend to reduce the local field gradient $`\frac{dB}{dx}`$, and hence the current density $`J`$. In analyzing relaxation measurements, the initial magnetic field distribution is assumed to be that of Bean’s critical state wherein $`J`$ is replaced by the critical current density $`J_c`$. If the relaxation is close to $`J_c`$, the magnetisation decay $`M(t)`$ is then theoretically known to be logarithmic, as is the case in most low-$`T_c`$ superconductors . In this case, the effective pinning potential $`U(J)U_0(JJ_c)`$ is a good approximation.
In high-$`T_c`$ superconductors (HTSC), the large thermal energy available leads to rapid decay of $`M(t)`$. Experimentally the relaxation is observed to be non-logarithmic over several decades in time in these materials . This is interpreted as arising from a non-linear $`U(J)`$. The vortex-glass theory and collective-creep theory , which predicts a low temperature true superconducting state with finite $`J_c`$, expects a pinning potential of the form $`U(J)=U_0[(J_c/J)^\mu 1]`$ which diverges in the limit $`J0`$. An important experimental observation in HTSC is an apparent universal value of the normalized relaxation rate $`S(T)=\frac{1}{M}\frac{dM}{d\mathrm{ln}t}`$ around 0.02-0.035 over a wide range of $`T`$ . For the $`U(J)`$ given above, such a small $`T`$-independent $`S`$ requires $`\mu >2`$ which is beyond the range of existing theoretical models . Experiments have also indicated power law decay in HTSC which requires a logarithmically diverging $`U(J)=U_o\mathrm{ln}(J_c/J)`$ . This form of $`U(J)`$ cannot account for the experimentally observed plateau in $`S(T)`$.
It becomes then important to identify the relaxation behaviour of the magnetisation in a type-II superconductor in presence of a uniform background of pinning centers. Towards this end, we simulate the time decay of remanent magnetisation $`M(t)`$ at finite $`T`$ in a simple model of 2D Josephson junction array (JJA). The magnetisation in this model is studied by including the screening currents through the inductance $`L_{𝐑,𝐑^{}}`$ between the cells. The underlying discrete lattice of junctions provides an energy barrier for the vortex motion within the array and is the source of vortex pinning in JJA. The behaviour of JJA with screening is parameterized by $`\lambda _J^2=\frac{\mathrm{\Phi }_0}{2\pi L_0I_c}`$, where $`I_c`$ is the critical current of the junction, $`L_0`$ is the self-inductance of the cell, and $`\mathrm{\Phi }_0`$ is the quantum of flux. Detailed simulation at $`T=0`$ have shown that for $`\lambda _J^2<1`$, the magnetic response of this model is similar to a continuum hard type-II superconductor . In this paper, we show that the thermo-remanent relaxation of $`M(t)`$ of JJA captures essential features of magnetic relaxation in HTSC. At a crossover temperature $`T_{cr}`$, a new time scale $`\tau _\beta `$ emerges in an intermediate time window which along with the characteristic time scale $`\tau _\alpha `$ for long time behaviour governs the relaxation. The characteristic time scales $`\tau _\alpha `$ and $`\tau _\beta `$ diverges as a power law at a temperature $`T_{sc}`$ at which $`M(\tau )M_00`$ as $`\tau \mathrm{}`$.
## 2 The Model
We consider a 2D array of superconducting islands forming a homogeneous square lattice of $`N\times N`$ unit cells in the $`xy`$ plane. Tunneling of the macroscopic wavefunction $`\mathrm{\Psi }(𝐫)=\psi \mathrm{exp}[i\phi (𝐫)]`$ across neighbouring islands lead to Josephson coupling between them. In presence of an applied magnetic flux $`\mathrm{\Phi }_{ext}`$ (per cell) along $`\widehat{z}`$ direction, the junction behaviour is fully described by dynamics of the gauge-invariant phase difference $`\varphi _{𝐫,\delta }=\phi (𝐫)\phi (𝐫+\delta )\frac{2\pi }{\mathrm{\Phi }_0}_𝐫^{𝐫+\delta }𝐀𝑑𝐥`$ between neighbouring islands. Here, $`\delta `$ is a unit vector, and $`𝐀`$ is the vector potential corresponding to the total magnetic field. The inset of Fig.1 shows a schematic array of size $`N=5`$ along with variable $`\varphi _{𝐫,x}`$ and $`\varphi _{𝐫,y}`$.
The magnetic response of JJA is set by the induced flux $`\mathrm{\Phi }_{ind}(𝐑)`$ due to screening currents which is modeled by considering the geometrical inductance matrix $`L`$ of the array . This allows us to write the induced flux in a cell at $`𝐑`$ as $`\mathrm{\Phi }_{ind}(𝐑)=_𝐑^{}L(𝐑,𝐑^{})I(𝐑^{})`$ where $`I(𝐑^{})`$ is the cell current at $`𝐑^{}`$ and $`L(𝐑,𝐑^{})`$ is the mutual inductance between the cells at $`𝐑`$ and $`𝐑^{}`$. In order to ease the prohibitive computation cost involved when mutual inductance is considered , we consider the induced flux only due to self-inductance $`L_0=L(𝐑,𝐑)`$ of the cell. This approximation is equivalent to the case of a long (ideally infinite) superconductor parallel to an applied field for which demagnetisation factor $`N_D=0`$. The total flux at R then can be written as $`\mathrm{\Phi }_𝐑=\mathrm{\Phi }_{ext}+L_0I_𝐑`$ (the cell co-ordinates are used as subscripts). The cell current is only a convenient variable for introducing $`\mathrm{\Phi }_{ind}`$ as the current through the junction is $`I_{𝐫,\delta }=I_𝐑I_{𝐑\delta }`$ where r is the junction common to cells at R and $`𝐑\delta `$. Since, $`I_𝐑`$ is divergence-less, Kirchoff’s law is automatically satisfied at each node of the lattice.
The dynamical variable $`\varphi _𝐫`$ (subscript $`\delta `$ is implicit) is related to the cell current $`I_𝐑`$ through the flux-quantisation condition
$$\underset{𝐫𝐑}{}\varphi _𝐫=2\pi \frac{\mathrm{\Phi }_𝐑}{\mathrm{\Phi }_0}=2\pi \frac{\mathrm{\Phi }_{ext}}{\mathrm{\Phi }_0}\frac{2\pi }{\mathrm{\Phi }_0}L_0I_𝐑,$$
(1)
where the summation is taken around the cell in anti-clockwise direction. The equation of motion for $`\varphi _{𝐫,\delta }`$ then follows from
$$\frac{d\varphi _𝐫}{dt}=\mathrm{\Gamma }\frac{\delta }{\delta \varphi _𝐫},$$
(2)
where the Hamiltonian $``$ of the system is
$$=\underset{𝐫}{}E_J(1\mathrm{cos}\varphi _𝐫)+\frac{1}{2}\underset{𝐑}{}L_0I_𝐑^2,$$
(3)
where $`\mathrm{\Gamma }=(\frac{\mathrm{\Phi }_0}{2\pi })^2\frac{1}{R}`$ with $`R`$ and $`I_c`$ as the normal state resistance and critical current of the junction respectively. In the above equation, the first term is summed over all bonds and represent Josephson coupling energy for the junction, whereas the second term is the magnetic field energy due to screening currents and the summation is over all cells in the array.
The above equations can be written compactly by introducing the matrix M for the lattice curl operation . Eq.(1) then becomes $`𝖬\varphi =2\pi f\frac{2\pi }{\mathrm{\Phi }_0}L_0I_m`$ where $`\varphi `$ and $`I_m`$ are the column vectors formed by $`\varphi _𝐫`$ and $`I_𝐑`$ respectively, and $`f=\mathrm{\Phi }_{ext}/\mathrm{\Phi }_0`$. Also, the current through the junction $`I_b=𝖬^TI_m`$, where $`𝖬^T`$ is transpose of the matrix M. For finite temperature simulation, we couple the heat bath through the noise current in the shunt resistor $`R`$ across the junction. Langevin equation for the array then takes a simple form
$`{\displaystyle \frac{d\varphi }{d\tau }}`$ $`=`$ $`𝖬^TI_m\mathrm{sin}\varphi +X(\tau ),`$
$`𝖬\varphi `$ $`=`$ $`2\pi f{\displaystyle \frac{1}{\lambda _J^2}}I_m.`$ (4)
Here, the current $`I_m`$ is scaled by the critical current $`I_c`$ of the junction. The dimensionless time $`\tau =\frac{2\pi RI_c}{\mathrm{\Phi }_0}t`$, and $`\lambda _J^2`$ has been defined earlier. The random term $`X(\tau )`$ has zero mean and white noise correlation
$$X_𝐫(\tau )=0,\mathrm{and}X_𝐫(\tau )X_𝐫^{}(\tau ^{})=2T\delta (\tau \tau ^{})\delta _{𝐫,𝐫^{}},$$
(5)
where $`T`$ is the temperature of the bath (in units of $`I_c\mathrm{\Phi }_0/2\pi k_B`$). In this unit, the phase coherence between neighbouring islands is established at $`T=1`$ and is the superconducting transition. The set of equations in (4) is solved self-consistently at each time step with free-end boundary condition. For simplicity, $`I_c`$ is assumed to be independent of $`f`$ and $`T`$ (the $`f`$ dependence does not show any qualitative change in the results presented below). The magnetisation $`M`$ is obtained as $`M=(1/N^2)_𝐑(\mathrm{\Phi }_𝐑/\mathrm{\Phi }_0)f`$.
The simulations were performed for $`N=16(256`$ cells) with time step $`\mathrm{\Delta }\tau =0.05`$. The results presented below are for $`\lambda _J^2=0.1`$. Note that for $`\lambda _J^2<1`$, each cell accommodates more than a quantum of flux. The moderate size of the array allowed long time to be reached in the relaxation that is essential at low temperatures (the longest relaxation time run was of $`5\times 10^6\tau 10^8`$ iterations). The field cooling is done by quenching from a high temperature (typically $`T=2`$) to the temperature $`T`$ under consideration, followed by annealing for $`5\times 10^38\times 10^3\tau `$ before switching off the applied field at $`\tau =0`$ (thermoremanent magnetisation). We have also carried out simulation for slow cooling which shows no qualitative difference from that obtained after a sudden quench, though a small but discernible dependence on the cooling rate is observed. We defer such effects and further details to a future paper, and present here results which brings out generic features.
## 3 Results and discussions
The $`M(\tau )`$ is shown in Fig.1 over 7 decades of time after cooling to different temperatures in an applied field $`f=5`$ (curves for intermediate $`T`$ is not shown in order to avoid over crowding of the figure). The curves are scaled by $`M(\tau =0)`$. From the curves, two distinct temperatures $`T_{cr}0.26`$ and $`T_{sc}0.04`$ can be identified at which the dynamical behaviour changes remarkably. For $`T>T_{cr}`$, $`M(\tau )`$ can be characterized by a single time scale $`\tau _\alpha `$. At $`T=T_{cr}`$, $`M(\tau )`$ develops a kink in an intermediate time window which at lower temperatures evolve into a plateau. The $`\tau _\alpha `$ now characterizes the long time behaviour of $`M(\tau )`$. On the plateau, the magnetisation $`M(\tau )=M_0`$ is $`T`$-independent. The temporary freezing of dynamics indicates emergence of a new time scale, represented by $`\tau _\beta `$, which governs the relaxation for $`\tau <<\tau _\alpha `$. Both $`\tau _\alpha `$ and $`\tau _\beta `$ increases rapidly at lower temperatures, and at $`T_{sc}`$, $`M(\tau )`$ freezes asymptotically to $`M_0`$ for the time probed in simulation (note that for $`T_{sc}<T`$, $`M(\tau )`$ becomes zero in the limit $`\tau \mathrm{}`$). The temperature $`T_{sc}`$ thus marks the transition into a state with true persistent current. For $`T<T_{sc}`$, $`M(\tau )`$ relaxes towards $`M_0`$.
Also shown in the inset of Fig.1 is the magnetic susceptibility $`\chi =M/f`$ under field cooled (FC) and zero-field cooled (ZFC) conditions (for ZFC, $`f`$ is applied at $`T=0`$ and $`\chi `$ is calculated while increasing $`T`$). For $`T>T_{cr}`$, the $`\chi _{FC}`$ and $`\chi _{ZFC}`$ are equal and $`M(T)`$ is reversible. At $`T_{cr}`$, the difference $`\chi _{FC}\chi _{ZFC}>0`$, and the $`\chi _{ZFC}`$ below it depends strongly upon the thermal history . Similar behaviour in bulk superconductors allow us to identify $`T_{cr}`$ as the irreversibility temperature. Though $`\chi _{FC}`$ does not show any change at $`T_{sc}`$, $`\chi _{ZFC}`$ is $`T`$ independent below it. The $`T`$ independent value of $`\chi _{ZFC}0.6`$ gives $`M_{ZFC}3.0`$ which equals remanent magnetisation $`M_0`$ obtained from the relaxation for $`TT_{sc}`$ and is also found to hold for other values of $`f`$. This equality between the $`M_0`$ and $`M_{ZFC}`$ can be understood by invoking the relation $`M_{rem}(H)=M_{FC}(H)M_{ZFC}(H)`$ ($`H`$ is the applied field) which is experimentally known to be valid in bulk type-II superconductors . In case of strong flux pinning, $`M_{FC}0`$ and $`M_{rem}M_{ZFC}`$, which is the case observed here. Validity of the above relation is also an evidence for strong flux pinning in JJA arising due to discrete underlying lattice of junctions. Further, we analyze the relaxation occurring on different time scales.
The long time decay of $`M(\tau )`$, which we term as $`\alpha `$-relaxation, fits Kohlrausch-Williams-Watt (KWW) law $`M(\tau )=\mathrm{exp}[(\frac{\tau }{\tau _\alpha })^\alpha ]`$ over the temperature range $`T>T_{sc}`$ probed in the simulation. The $`\tau _\alpha (T)`$ shows 6 orders of increase between $`T_{sc}`$ and $`T_{cr}`$ implying that the system falls out of equilibrium on cooling through it. Though $`\tau _\alpha (T)`$ appears to fit Arrhenius law $`\tau _\alpha (T)e^{A/T}`$ (see Fig.2 inset (a)), notably, a distinct power law is observed at low temperatures as shown in Fig.2 where $`\tau _\alpha ^1`$ is plotted against $`TT_{sc}`$. Fit to $`\tau _\alpha (T)=\tau _0TT_{sc}^\gamma `$ gives $`T_{sc}=0.045,\gamma =3.74`$ for $`f=5`$, and $`T_{sc}=0.05,\gamma =4.735`$ for $`f=2`$ which is also included in the figure. The value of $`T_{sc}`$ obtained from the fit is not very different from that at which the relaxation is frozen asymptotically at $`M_0`$ but it must be treated with considerable reserve. It is probable that the actual value may be lower than that obtained here as the time scale probed at low temperatures put severe constraint on obtaining $`\tau _\alpha `$. In Fig.2, the deviation from linearity occurs at $`T0.24`$ which is close to $`T_{cr}`$ below which the relaxation shows a temporary frozen state. The $`T`$ dependence of the exponent $`\alpha `$ is also shown in Fig.2 (inset (b)). Due to small value of $`\alpha `$ at low temperatures, the $`M(\tau )`$ can be fit to $`\mathrm{log}\tau `$ over 1-2 decades in time, and is the regime where thermally activated flux-creep theory can be applied .
Appearance of a new time scale $`\tau _\beta `$ for $`T<T_{cr}`$ leads to two-step relaxation process : an initial part towards the plateau during $`\tau _0\tau \tau _\beta `$, and a later part away from the plateau in the interval $`\tau _\beta \tau \tau _\alpha `$ which at later time develops into $`\alpha `$-relaxation. Here, $`\tau _020\tau `$ is of the same order which appears in the power law fit for $`\tau _\alpha `$. We refer the relaxation during $`\tau _0\tau \tau _\alpha `$ as the $`\beta `$-relaxation regime due to apparent qualitative similarity with the $`\beta `$-relaxation seen in supercooled liquids . For $`T_{sc}<T`$, the relaxation during $`\tau _\beta \tau \tau _\alpha `$ fits a power law $`M(\tau )M_0c(\tau /\tau _\beta )^b`$ over 3 decades with exponent $`b`$ independent of $`T`$. This is shown in Fig.3 where $`M(\tau )`$ (for $`f=5`$) for different temperatures can be scaled on to a master curve when plotted against $`\tau /\tau _\alpha (T)`$. The thick line is the power law fit with $`b=0.85`$ and holds for $`\tau /\tau _\alpha (T)1`$ as evident from the fit. The scaling exponent $`b`$ allow us to obtain the $`T`$ dependence of $`\tau _\beta (T)`$ which is shown in Fig.3 (inset). $`\tau _\beta (T)`$ fits to a power law of the form $`\tau _\beta =\tau _0TT_{sc}^\psi `$ with $`\psi =1.58`$ and $`T_{sc}=0.033`$ ($`\tau _0`$ is of the same order as that obtained for $`\tau _\alpha `$). Since the $`\beta `$-relaxation at later time evolves into $`\alpha `$-relaxation, divergence in $`\tau _\beta `$ and $`\tau _\alpha `$ must occur at the same temperature. We attribute the small difference seen here to uncertainties in obtaining $`\tau _\alpha `$. Nevertheless, divergence of $`\tau _\beta `$ and $`\tau _\alpha `$ at $`T_{sc}`$ is significant as it implies absence of flux-creep (hence, divergent $`U(J)`$) as $`\tau \mathrm{}`$, and is consistent with the idea of true superconducting state below a transition temperature.
For $`\tau _0\tau \tau _\beta `$, the relaxation occur towards the plateau, and for $`T<T_{sc}`$ is asymptotically frozen on it. The relaxation can again be fit to a power law $`M(\tau )M_0c^{}(\tau /\tau _\beta )^a`$. The exponent $`a`$ is dependent on $`T`$ as shown in Fig.4. Also shown in the figure is $`S(T)`$ close to plateau which attains a $`T`$ independent value $`0.076`$ (on the plateau, $`S(T)`$ is an order of magnitude lower than this). This value is of the same order as that observed in HTSC for which it falls in the range $`0.020.035`$ . A plateau is $`S(T)`$ coupled with non-logarithmic decay of $`M(\tau )`$ over 4 decades is significant in view of similar experimental observation in HTSC. To understand the processes involved during the relaxation, the spatial distribution of $`\mathrm{\Phi }_𝐑`$ in different time windows is obtained and is shown in the inset of Fig.4. The slope of the flux profile $`(\mathrm{\Phi }/R_x)J`$ is observed to be $`T`$-independent on the plateau and equals the slope ($`J_c`$) in the remanent state at $`T=0`$ (note that $`I_c`$ of the junction is assumed to be $`T`$-independent in the simulation). At $`T=0`$, such a state have been shown to be the self-organized Bean’s critical state for a hard superconductor .
Thus, across the $`\beta `$-relaxation regime, the current density relaxes from the super-critical state $`J>J_c`$ towards the sub-critical state $`J<J_c`$. The power law behaviour in this regime can be attributed to the self-organization of magnetic flux around the critical current density $`J=J_c`$. The self-organization of vortices is driven by inter-vortex repulsion and vortex-pinning center attraction alone , and is independent of the temperature. This explains the plateau in $`S(T)`$ in Fig.4, and is also consistent with the apparent universality of $`S(T)`$ repeatedly observed in HTSC . Moreover, in HTSC the critical current density $`J_c`$ is 1-2 orders of magnitude less than that of low-$`T_c`$ superconductors over a wide region of $`HT`$ phase diagram. Since the relaxation towards $`M_0J_c`$ is a power law, the time required to set up critical state in HTSC is much longer compared to low-$`T_c`$ superconductors. This implies that relaxation measurement in HTSC is influenced by the self-organization for much longer time and could be the source of non-logarithmic behaviour. It is important to note that the $`\mathrm{log}t`$ relaxation due to thermal activation becomes a dominant process only in the sub-critical state.
The results can be summarized by a $`\tau T`$ plot for the magnetic relaxation in hard type-II superconductors which is shown in Fig.5. Emergence of a new time scale at a temperature at which history dependence sets in defines dynamically the irreversibility temperature for hard superconductors. The regime $`T_{cr}<T`$ is analogous to the vortex liquid (VL) phase with critical current density $`J_c=0`$. With decreasing $`T`$ below $`T_{cr}`$, the flux motion becomes rapidly viscous as evident from the power law increase in $`\tau _\alpha `$. The overall behaviour is analogous to the “supercooled” state in glass formers . In these system, relaxation of the correlation in density fluctuation shows scaling behaviour in the $`\beta `$-relaxation regime which is remarkably similar to Fig.3 here. Relaxation data in HTSC need to be reanalyzed to observe the $`T`$ dependence of the characteristic time scale across the irreversibility temperature and glass transition temperature as observed here. Also, in fabricated JJA, SQUID parameter $`\beta _J=\lambda _J^2`$ of 30 have been achieved which falls within the parameter range in which the simulation results can be applied.
In conclusion, simulation of the magnetic relaxation in a model of hard superconductor shows emergence of a new time scale below the irreversibility temperature. Self-organization of the magnetic flux around $`J_c`$ in this time scale leads to power law decay of the magnetisation. Divergence of time scales at a transition temperature leads to a state with finite remanent magnetisation (and hence, persistent current). The temperature dependence of the normalized relaxation rate is in good agreement with experiments on HTSC.
\***
Acknowledgements : The author is grateful to D. Dhar and A. K. Grover for useful discussions and critical reading of the manuscript.
## 4 Figure Captions
* Relaxation of the remanent magnetisation $`M(\tau )`$ on $`\mathrm{log}\tau `$ scale for $`\lambda _J^2=0.1,f=5`$ and array size $`16\times 16`$. The curves (from left to right) are for temperature $`T`$=0.9, 0.6, 0.4, 0.3, 0.24, and 0.20 to 0.02 in steps of 0.02. Inset: Plot of $`\chi (T)`$ for $`f=5`$ in FC and ZFC state. The irreversibility temperature is also the crossover temperature $`T_{cr}`$ (marked in the inset) at which $`M(\tau )`$ develops a kink. Also shown is a typical $`5\times 5`$ array with $`\varphi _{𝐫,x},\varphi _{𝐫,y}`$ and $`I_𝐑`$ for a single cell.
* Log-log plot of $`\tau _\alpha ^1(T)`$ against $`TT_{sc}`$ for all values of $`T`$ for which the simulation was performed. The symbols $``$ and $``$ are for $`f=2,f=5`$, respectively. The errors are less than the symbol size. The fitting parameter $`T_{sc}`$ and $`\gamma `$ is also given in the plot. Inset: (a) Arrhenius plot $`\tau _\alpha (T)`$ vs $`T^1`$. (b) The stretching parameter $`\alpha `$ as a function of $`T`$.
* $`M(\tau )`$ vs rescaled time $`\tau /\tau _\alpha (T)`$ for $`0.05T0.13`$ (in steps of 0.01) in the $`\beta `$-relaxation regime. The thick dashed line is a fit to $`M(\tau )\tau ^b`$ with $`b=0.85`$. Inset : $`\tau _\beta `$ as a function of $`TT_{sc}`$ for $`0.05T0.14`$ obtained by fitting the relaxation away from the plateau. The full line is a fit to $`\tau _\beta (TT_{sc})^\psi `$ with $`\psi =1.58\pm 0.2`$ and $`T_{sc}=0.033`$.
* The exponent $`aS`$ (normalized relaxation rate) for $`T0.12`$. Also shown is the $`S`$ obtained from $`M(\tau )`$ around the plateau in Fig.1. Inset : The flux distribution $`\mathrm{\Phi }_𝐑/\mathrm{\Phi }_0`$ across a central row of cells in the array for various $`\tau `$ at $`T=0.07`$ for $`f=5`$.
* The $`\tau T`$ plot showing various regimes obtained from the simulation. The dotted line is obtained by interpolating $`\tau _\beta (T)`$ to $`T_{cr}`$. A and B are the power law regimes with exponents $`a`$ and $`b`$, respectively. |
warning/0001/cond-mat0001068.html | ar5iv | text | # Holons on a meandering stripe: quantum numbers
## I Introduction
The nature of charge carriers in high-$`T_c`$ cuprate superconductors remains a subject of debate. The stoichiometric (“parent”) compounds are antiferromagnetic (AF) insulators well described by the Heisenberg model. Experiments show that, at least in some of the cuprates, doping with holes creates an intrinsically inhomogeneous state with periodically modulated charge and staggered spin densities. Neutron and X-ray scattering experiments indicate that staggered magnetization has a period of modulation that is twice as long as that of charge density. This is consistent with the notion of charged stripes separating AF domains with alternating Neel magnetization.
Stripes as domain walls in the ground state of a collinearly ordered antiferromagnet were predicted—prior to reliable experimental detection—on the basis of Hartree-Fock studies of the Hubbard model near half-filling. Mean-field calculations yield a linear density of $`\nu =1`$ doped hole per lattice site, which almost certainly means insulating stripes, in apparent contradiction with experiments. Besides, stripes observed in the cuprates tend to have a linear hole density $`\nu 1/2`$, at least when they are sufficiently well separated. To date, no reliable microscopic calculation yields the experimentally observed filling fraction $`\nu `$, let alone explains the transport and high-temperature superconductivity in the cuprates. Numerical simulations have so far been inconclusive.
In the absence of a reliable microscopic theory, attempts have been made to find a phenomenological description of the stripes. In one of the more popular routes, a stripe is modeled as a one-dimensional electron gas (1DEG) interacting with the surrounding environment and with the transverse motion of the stripe. As we argue below, this approach rests on the assumption that the physics of an isolated stripe is basically the same as in the limit $`\nu 1`$. In this limit, a stripe is almost completely filled with holes and can be described as an electron gas at low density $`1\nu `$. It is far from obvious, though certainly not implausible, that stripes with $`\nu 1`$ and $`\nu 1/2`$ should exhibit qualitatively similar behavior.
In this work we develop a qualitatively different (but not less plausible) phenomenology of a partially doped stripe. It is based on two different model calculations performed in the limit of low hole density on a stripe, $`\nu 1`$. Nominally, this is as far from the observed density $`\nu 1/2`$ as the electron-gas limit $`\nu 1`$. The quantum numbers of charge carriers (holons) in our model calculations are completely different from those of electrons. It likely means that the two limits are not adiabatically connected. It is clear then that the phase $`\nu 1/2`$ can resemble only one of the low-density limits: either $`\nu 1`$, or $`\nu 0`$—or possibly none of the above!
Building on our model calculations we conjecture that charge carriers of the “$`\nu 0`$” phase are holons (charge $`Q=1`$, spin $`S_3=0`$). In both models the loss of spin is compensated by the emergence of another spin-like degree of freedom, termed the transversal flavor by Zaanen et al. This happens because a holon always resides on a transverse kink or antikink of a domain wall. Thus holons are strongly coupled to transverse fluctuations of a stripe. Yet, the motion of such objects along the stripe is free, it does not produce any additional spin frustration. Such a holon gas is clearly very different from the electron gas of the “$`\nu 1`$” phase.
Of course, our approach should not be interpreted as a suggestion that a stripe with small linear hole density $`\nu 1`$ can be stable in any model relevant to high-$`T_c`$ materials. On the contrary, a domain wall in an undoped antiferromagnet is a highly excited texture. A finite linear density of holes is needed to stabilize a domain wall. In our analysis we always assume that the domain wall is supported externally (e.g., by the boundary conditions), while its untwisting is suppressed by a sufficiently strong anisotropy (strictly linear polarization in our Hartree-Fock analysis). In a model where partially filled stripes appear in the ground state, neither assumption would be necessary.
Indeed, we have reduced the symmetry of the problem in order to stabilize topologically the Ising-type domain walls. In any model with the O(3) symmetry of the Neel order parameter, there can be no topological arguments for their stability. For example, domain walls can be continuously untwisted in a broad class of Ginzburg-Landau models with a continuous O$`(N)`$ symmetry; such domain walls are locally unstable. However, local instabilities are not an issue for globally stable configurations. In practical terms, untwisting does not occur if domain walls appear in the ground state of a system. In the context of high-$`T_c`$ stripe phases, such nontopologically-stable domain walls were discussed in Ref. . Their global stability requires frustration of the AF order on some microscopic or intermediate length scale, e.g., as a result of doping.
Therefore, our model calculations should be viewed as an attempt to identify plausible ground states of an isolated stripe. In contrast to phenomenological approaches, we do not postulate effective one-dimensional models of a stripe. Instead, we derive them by starting with a two-dimensional model describing an antiferromagnet with a domain wall. Our 2D models may be unrealistic for the cuprates, but the resulting 1D effective theories have the set of elementary excitations consistent with the paradigm of a stripe as a doped fluctuating domain wall in an antiferromagnet. In essence, we rely on universality: if the number of qualitatively different ground states of a stripe (classified by quantum numbers and spectrum of low-lying excitations) is limited, all of them may be derived from simple 2D models.
From this perspective, our work adds the 1D holon gas to the list of potential stripe models. Note, however, that, in the presence of interactions, the description of a stripe in terms of holons, and a more conventional one in terms of electrons, are not necessarily incompatible. Both models can be viewed as Luttinger liquids with different collective modes: charge and spin in the case of electrons, charge and transverse fluctuations for holons. Thus, 1D electrons with a spin gap and 1D holons with a transverse gap may well represent one and the same phase. We intend to discuss the role of interactions among the holons in a future publication.
The paper is organized as follows. In Sec. II we analyze partially doped domain walls in a $`t`$$`J`$ model with large Ising anisotropy. The ground state and the spectrum of elementary excitations (spinons and holons) are found explicitly in the limit of small doping $`\nu `$. In Sec. III we present numerical evidence for holons in the Hubbard model, within the Hartree-Fock approach. The Hartree-Fock equations for the Hubbard model are further analyzed in Sec. IV, where we introduce an appropriate long-wavelength approximation and study the spectrum of midgap states induced by a domain wall. We give heuristic arguments for the existence of fermion zero modes around a transverse kink on a domain wall. Technical details are collected in the Appendixes.
## II Holons on a domain wall: $`t`$$`J`$ model with Ising anisotropy.
A simple yet very instructive example of holon gas on a domain wall is offered by the $`t`$$`J`$ model with Ising anisotropy, previously considered by Kivelson et al.,
$`H_{tJ_z}`$ $`=`$ $`{\displaystyle \underset{\mathrm{𝐫𝐫}^{}}{}}\{[ta_\sigma ^{}(𝐫^{})a_\sigma (𝐫)+{\displaystyle \frac{J_{}}{2}}s_+(𝐫^{})s_{}(𝐫)+\mathrm{h}.\mathrm{c}.]`$ (2)
$`+J_zs_z(𝐫^{})s_z(𝐫)+Vn(𝐫^{})n(𝐫)\},`$
with the usual constraint of no double occupancies; the sum is taken over pairs of nearest-neighbor sites. The $`t`$$`J`$ model proper is restored if we set $`J_{}=J_z=4V`$.
The analysis of the model (2) is greatly simplified in the strongly-anisotropic limit $`J_z\mathrm{}`$ (while the energy of an AF bond $`E_b=VJ_z/4`$ may remain finite). In this case, an isolated hole in the AF bulk is localized and has the energy $`ϵ=4E_b`$, the cost of 4 missing bonds. Two holes on adjacent sites share one missing bond, which is to say that their interaction energy is $`E_b`$. When severing a bond is costly ($`E_b<0`$ and has a large absolute value), hole-rich islands are formed in the otherwise unaltered antiferromagnet, with the energy $`2E_b`$ per doped hole (phase separation in the bulk). In the opposite limit, $`E_b|t|`$, doped holes strongly repel one another and stay apart. The energy per hole in this case is
$$\epsilon _{\mathrm{hole}}=4E_b=J_z4V.$$
(3)
Phase separation or not, doped holes are immobilized in the bulk of an antiferromagnet. As long as $`J_z`$ greatly exceeds both $`t`$ and $`J_{}`$, the cost of frustrated (ferromagnetic) bonds produced by a moving hole outweighs the reduction in kinetic energy.
This situation changes dramatically in the presence of a domain wall: doped holes become mobile.
### A Spinons on a domain wall
Elementary excitations of an undoped domain wall are kinks formed in pairs by flipping two nearest-neighbor spins \[Fig. 1 (a), (b)\]. The kinks are mobile, carry zero charge and spin $`s_z=\pm 1/2`$. We thus term them spinons. To make their relation to 1D spinons more explicit, we have integrated spin $`s_z`$ across the domain wall with a smooth envelope to obtain an effective 1D spin chain representing the domain wall \[Fig. 1 (a), (b), open symbols\]. By using Bloch states, one finds the spinon energy spectrum,
$$E_{\mathrm{spinon}}(k_x)=J_z/2+J_{}\mathrm{cos}2k_x+𝒪(J_{}^2/J_z)>0.$$
(4)
Clearly, for $`J_zJ_{}`$, such excitations are strongly gapped; an undoped domain wall is very stiff.
### B Holons on a domain wall
A single hole doped into a straight domain wall \[Fig. 1 (c), (d)\] cannot hop along the stripe in the limit $`J_z\mathrm{}`$: it must leave a spinon at the original location of the hole, which costs $`𝒪(J_z)`$ in magnetic energy. Therefore, a single hole can only oscillate across the stripe; the corresponding ground-state energy is
$$E_{\mathrm{hole}}=4E_b\frac{J_z}{2}|t|+𝒪(t^2/J_z).$$
(5)
It is important to realize that, once a hole starts moving along the domain wall, it does not create a string of ferromagnetic bonds, which localize a hole in the bulk. A moving charge is now associated with a kink in the transverse position of the stripe \[Fig. 1 (d)\]. This composite object has spin $`S_z=0`$ and charge $`Q=+1`$. Following the spin-chain convention it can be termed a holon. Moreover, holons can be created in pairs without spinons: two spinons created by two holes can annihilate each other. We can say that a localized doped hole decays virtually into a spinon and a holon. Another hole nearby can absorb the costly spinon and become a holon. The intermediate spinon is not needed if the two holes are on adjacent sites on the same side of the wall.
A holon with momentum $`k_x`$ along the stripe has the energy
$$E_{\mathrm{holon}}(k_x)=4E_b\frac{J_z}{2}2|t|\mathrm{cos}k_x+𝒪(t^2/J_z).$$
(6)
For small momentum $`k_x`$, this is smaller then the energy of a single hole (5). Therefore, a dilute gas of holons has a lower energy than a collection of holes similar to that in Fig. 1 (c). As shown in Appendix A, a dilute holon gas is stable against phase separation for $`E_b>0`$ when the interaction between the holons is strictly repulsive. When $`E_b<0`$, holes on adjacent sites attract. This attraction wins over an increase in kinetic energy for $`E_bt`$, causing phase separation: holes can lower their energy by forming a densely populated ($`\nu =1`$) stripe leaving the rest of the domain wall undoped (see Appendix A).
### C Preferred linear charge density
So far we have considered an antiferromagnet with a single domain wall maintained by the appropriate boundary conditions. Within this model, one can study arbitrary linear concentrations of holes $`\nu `$ on the wall. Particularly simple situations are the limit of dilute holes $`\nu 0`$ (gas of holons with a transverse flavor) and the opposite limit $`\nu 1`$ (1D electron gas).
Because $`J_z`$ is large, partially doped ($`\nu <1`$) stripes do not occur naturally in this model. The preferred linear density of charge $`\nu `$ can be found by using the usual Maxwell construction. To do so, one minimizes the energy per doped hole—including the cost of creating domain walls. Since partially doped domain walls contain costly ferromagnetic bonds, they will not occur if $`J_z\mathrm{}`$. A lower bound for the energy per doped hole is
$$\epsilon (\nu ^11)\frac{J_z}{2}4E_b+𝒪(t).$$
In the limit we consider, this expression is a strictly decreasing function of $`\nu `$; the optimal configuration of the stripe corresponds to $`\nu =1`$.
### D Implications
The anisotropic $`t`$$`J`$ model (2) dominated by the Ising term provides a good illustration to the strategy outlined in the Introduction. True, this simple model predicts insulating stripes with $`\nu =1`$, contrary to experimental observations. Nevertheless, it has enabled us to find two possible phases of conducting stripes that may arise in more realistic models: the 1D electron gas (the “$`\nu 1`$” phase) suggested previously and the 1D holon gas (the “$`\nu 0`$” phase). To do so, we have created a domain wall by fixing boundary conditions and doped it to any given hole density $`\nu `$. At this stage, the simplicity of the model turns into a virtue: quantum numbers of elementary excitations can be readily determined.
When macroscopic phase separation is absent, i.e., for $`E_b>0`$ ($`J_z<4V`$ but $`t,J_{}J_z`$), a gas of holons is formed on a weakly doped domain wall. Holons are mobile kinks of the domain wall with charge $`Q=+1`$ and no spin, $`S_z=0`$. Compared to their one-dimensional counterparts, domain-wall holons have an additional flavor, $`\rho =\pm 1/2`$ (isospin), which denotes the direction of the associated transverse kink \[Fig. 1 (d)\]. An effective model describing the holon gas has been previously outlined in Ref. .
Even though it does not yield the observed stripe filling $`\nu 1/2`$, the strongly anisotropic limit of the model (2) has certain appeal: it is simple enough to permit controlled calculations. In particular, spin waves are gapped and all associated dissipation effects are suppressed. In addition, holes are not allowed to leave the domain wall and therefore cannot go around each other. A domain wall in this limit is a strictly one-dimensional object, yet its transverse motion is fully accounted for.
## III Holons on a domain wall: Hubbard model
Could holons be generic to domain walls in an antiferromagnet or are they just a curiousity of the $`t`$$`J`$ model in the Ising limit? To answer this question, we have attempted to find similar excitations in the Hubbard model. Clearly, the problem is much more difficult because there is no controlled approximation in this case, certainly not in two spatial dimensions.
When the coupling strength $`U`$ is weak compared to the free-electron bandwidth of $`4t`$, one can hope to find some guidance in mean-field Hartree-Fock (HF) solutions. This approach has been successful—to a degree—in predicting the existence of stripes in the cuprates. Numerical HF calculations show that, away from half-filling, doped charges form stripes along directions (0,1) or (1,1) with $`\nu =1`$ doped hole per lattice site along a stripe. Such stripes are, indeed, AF domain walls. Because they are filled with holes to capacity, charged excitations are gapped and thus an individual stripe is an insulator. (The same problem arises in the $`t`$$`J`$ problem in the large-$`J_z`$ limit considered above.)
Holons described in Section II are solitons with a well-localized charge distribution. To find analogous excitations in the Hubbard model (at the HF level) let us recall the properties of midgap states induced by domain walls. It is well known that solitons with anomalous quantum numbers arise in connection with fermion zero modes induced by such topological defects (see Appendix B). A uniform domain wall in 2 dimensions confines a midgap state only in one direction, across the wall. Solitons of finite extension—both across and along the wall—can exist only if there is an inhomogeneity on the domain wall. In this section we shall show that a bond-centered domain wall with a wiggle may “bind” a holon, i.e., a soliton with quantized charge $`Q=1`$ and spin $`S_3=0`$.
Although the solitons are static at the HF level, it is merely an artifact of the mean-field approximation, in which the average spin and charge densities are assumed to be time-independent. Mean-field configurations with solitons at different positions $`x`$ along the wall should be viewed as degenerate minima of the action, i.e., as classical solutions with a soliton at $`x`$. Quantization of the soliton restores broken translational symmetry: plane waves are superpositions of states with a soliton at all possible sites $`x`$. At the semiclassical level, the energy of a soliton is given by
$$E=\sqrt{E_0^2+p_x^2v_{1\mathrm{D}}^2}E_0+p_x^2v_{1\mathrm{D}}^2/2E_0,$$
(7)
where $`E_0`$ is the mean-field energy of the soliton and $`v_{1\mathrm{D}}`$ is the 1D Fermi velocity calculated using the static HF wavefunctions. As discussed below, the holon energy spectrum is similar to that (6) of the $`t`$$`J_z`$ model, although both the inverse mass and velocity are substantially reduced (in the weak-coupling limit), reflecting the collective nature of the soliton.
Just as in the case of the anisotropic $`t`$$`J`$ model, we are not dealing with the ground state of the model. A stripe needs a finite linear density of charge in order to be stable. Because collective excitations tend to be large at weak coupling, an appreciable linear density of holons likely requires the coupling to be strong, or else their overlap will completely destroy their individual properties. With only a Hartree-Fock approach at our disposal, we cannot access the strong-coupling limit (although we have tried to mimick it in the anisotropic $`t`$$`J`$ problem). Instead, we maintain the domain wall by boundary conditions, and vary the linear charge density along the stripe by changing the total number of holes in the system.
The weak-doping expansion in the Hubbard model, however, has an additional problem, which was not present in our analysis of the $`t`$$`J_z`$ model. Namely, the undoped domain walls are always unstable. Indeed, the undoped system can be accurately described by a Heisenberg-like model, and here the energy of a domain wall can be continuously lowered by perturbing in the direction orthogonal to the original magnetization vector. In the static HF configurations presented below \[Figs. 26\], this untwisting instability was suppressed by imposing a constraint of linear polarization. Nevertheless, these textures are the (particular) solutions of the full set of HF equations.
The domain wall (antiphase stripe) is favored by the holes. A fully-doped stripe at $`\nu =1`$ is both locally and globally stable already at the HF level. For a partially-doped stripe, this approximation (plus the constraint of linear polarization) gives static localized holes. We have found that the untwisting in the full set of HF equations (arbitrary polarization) starts to develop on the undoped portions of the stripe, compressing the remaining holes into segments of a fully doped stripe ending with semi-vortices similar to those discussed in Ref. . We believe that this is an artefact of the used weak-coupling approximation. Namely, we expect that at sufficiently strong coupling $`U`$ the effective attraction between the holes (caused by the untwisting) will be compensated by their increased mobility along the stripe. In some sence this is similar to what happens in the anisotropic $`t`$$`J`$ model (2) in the region $`tE_b<0`$: even though holes can gain some potential energy by sitting next to each other, the associated loss of their kinetic energy prevents phase separation along the stripe.
In the remainder of this section, we present numerical results obtained in HF calculations. Further analysis of the results in a long-wavelength approximation is given in the next section.
### A Hubbard model: numerical results
We have solved self-consistently HF equations of the Hubbard model
$`t{\displaystyle \underset{\mathrm{\Delta }r}{}}\psi (𝐫+\mathrm{\Delta }𝐫){\displaystyle \frac{U}{2}}[\stackrel{}{\sigma }(𝐫)\stackrel{}{\sigma }\rho (𝐫)]\psi (𝐫)=E\psi (𝐫),`$
(the notation is explained in Sec. IV) for collinear spin configurations,
$$s_1(𝐫)=s_2(𝐫)=0,s_3(𝐫)0,$$
(8)
at small and intermediate interaction strengths $`U=2t\mathrm{}4t`$. We always started with two AF domains having opposite values of staggered magnetization $`m`$. On the lattice row separating the two domains, staggered magnetization was initially disordered. Such configurations could thus later converge into site-centered, bond-centered, or meandering stripes.
#### 1 Undoped stripe
At half-filling, the stripe always became bond-centered, as in the anisotropic $`t`$$`J`$ model. We have explicitly verified that a site-centered stripe always has a higher energy in the absence of doped charges.
In a few cases, an initially disordered stripe converged to a state with a higher energy, a bond-centered stripe with a defect where the domain wall shifts one lattice spacing sideways \[Fig. 2(a)\]. Ripples in staggered magnetization $`(1)^{x+y}s_3(𝐫)`$ around the wiggle are an interference effect between staggered spin, varying as $`(1)^{x+y}`$, and a smooth component of $`s_3`$. By averaging $`s_3(𝐫)`$ over four neighboring sites,
$`m_{00}(𝐫)=`$ $`[s_3(𝐫)+s_3(𝐫+\widehat{𝐱})+s_3(𝐫+\widehat{𝐲})`$ (10)
$`+s_3(𝐫+\widehat{𝐱}+\widehat{𝐲})]/4,`$
one can suppress the staggered component and uncover a spin soliton residing at the wiggle \[Fig. 2(b)\]. An even better view of the soliton is afforded when spurious long-range spin-density oscillations induced by the boundary are removed by plotting the symmetrized spin density $`s_3(𝐫)+s_3(𝐫)`$/2 \[Fig. 2(d)\]. Because particle density remains equal to one everywhere, the soliton has zero charge (see Appendix B).
Numerically, the soliton has a total spin $`S_3=\pm 1/2`$:
$$S_3(R)=\underset{|x|R}{}\underset{|y|R}{}s_3(𝐫)\pm 1/2\text{ as }R\mathrm{}.$$
(11)
We have checked that the spin is well localized: $`S_3(R)S_3(\mathrm{})`$ vanishes exponentially with $`R`$ \[Fig. 2(c)\]. This spin soliton is an exact analogue of the spinon found in the $`t`$$`J`$ problem \[Fig. 2(b)\]. Here we also find spinons of two flavors, those bound to transverse kinks and antikinks.
#### 2 Stripe with one doped hole
With one hole added, an initially disordered stripe typically converged to one of the two bond-centered configurations. A straight bond-centered stripe would eventually contain a polaron (Fig. 3), a nontopological soliton with the quantum numbers of a hole ($`Q=1`$, $`S_3=\pm 1/2`$). The presence of a nonzero spin density is manifested in the typical ripples of the staggered magnetization $`(1)^{x+y}s_3(𝐫)`$—the result of the interference between the staggered and smooth spin components.
The other typical configuration was a bond-centered wall with a transverse kink \[Fig. 4\], on which a charge $`Q=+1`$ is localized. The absence of interference fringes is a sign of zero spin density. Indeed, integrated spin $`S_3(R)`$ is numerically zero (of order $`10^{13}`$) for all values of $`R`$. The absence of spin can also be verified by plotting $`m_{00}(𝐫)`$, the smoothed spin density (10). Since $`Q=+1`$ and $`S_3=0`$, one immediately recognizes a holon in this soliton. As its counterpart in the anisotropic $`t`$$`J`$ problem \[Fig. 1(d)\], it binds to a transverse kink or antikink of the domain wall.
Why does a wiggle on a domain wall change the nature of a doped charge in such a dramatic way? Essentially, a bond-centered domain wall can be thought of as a 1D AF chain (roughly, two parallel spins across the domain wall create an excess spin 1/2). A straight domain wall corresponds to a spin chain with perfect AF order. If there is a transverse kink on the domain wall, the staggered magnetization of the effective chain changes sign at the transverse kink, i.e., staggered magnetization itself has a kink. In analogy with polyacetylene (as discussed, e.g., by Berciu and John), a doped charge becomes either a polaron (no AF kink), or a holon (an AF kink is present). To illustrate this, we plot the $`x`$-staggered magnetization $`m_{10}(𝐫)`$
$`m_{10}(𝐫)=(1)^x[s_3(𝐫)s_3(𝐫+\widehat{𝐱})`$
$`+s_3(𝐫+\widehat{𝐲})s_3(𝐫+\widehat{𝐱}+\widehat{𝐲})]/4`$
for a straight wall \[Fig. 3(c)\] and for a wall with a wiggle \[Fig. 4(c)\]. In both cases, the $`x`$-staggered magnetization is confined to a narrow strip, which can be identified with the effective chain of the $`t`$$`J`$ problem. Clearly, $`m_{10}`$ alters the sign at the center of a holon \[Fig. 4(c)\] but is only slightly depressed around a polaron \[Fig. 3(c)\]. Later, we will substantiate these qualitative arguments with an analysis of midgap states (particularly, fermion zero modes).
Although the question of stability of large solitons on a weakly doped stripe is purely academic (see the discussion in the Introduction), we have compared energies of isolated polarons and holons. At weak coupling, $`U2.5t`$, polarons have a slightly lower energy. Two or more polarons preferred to bind into spinless bipolarons. For $`U3t`$, holons had a lower energy. Moreover, two well-separated holons (Fig. 5) had a smaller energy than a bipolaron for $`U3t`$. The trend is clearly to favor holons as the coupling gets stronger.
## IV Hubbard model: a Hartree-Fock analysis
To interpret the numerical results discussed in the previous section, we have conducted a thorough analysis of electron states induced in the middle of the Hubbard gap by an AF domain wall. Midgap states of a straight domain wall have been previously explained in great detail by Schulz. Studying a domain wall with a wiggle on a lattice is a rather challenging task, therefore we first derive a long-wavelength approximation that could capture the essential physics (e.g., the difference between bond-centered and site-centered walls). We find that midgap states of an undoped domain wall resemble a half-filled 1D chain with the Fermi momentum $`k_F=\pi /2`$. The relevant low-energy states on a domain wall are composed of Bloch states with momenta near $`(\pm \pi /2,\pm \pi /2)`$. This leads to a theory of 8-component fermions (2 spin components $`\times `$ 4 Fermi points). Finally, we relate holons observed numerically to fermion zero modes induced by a wiggle, and trace the origin of holon transverse flavor to the doubling of fermion components (8 instead of the usual 4 in the 1D case).
### A Mean-field equations
The HF equations for the Hubbard model are
$$t\underset{\mathrm{\Delta }r}{}\psi (𝐫+\mathrm{\Delta }𝐫)\frac{U}{2}[\stackrel{}{\sigma }(𝐫)\stackrel{}{\sigma }\rho (𝐫)]\psi (𝐫)=E\psi (𝐫),$$
(12)
where $`\psi (𝐫)`$ is a 2-component spinor wavefunction, $`\stackrel{}{\sigma }`$ is the triplet of Pauli matrices and the sum is over vectors $`\mathrm{\Delta }𝐫`$ pointing to the four adjacent sites. The expectation values of spin $`\stackrel{}{\sigma }(𝐫)`$ and density $`\rho (𝐫)`$ should be calculated self-consistently. As discussed in the Introduction, we are specifically looking for collinear solutions, therefore we set $`\sigma _1=\sigma _2=0`$. It is customary to rotate the spin axes on one of the sublattices through $`\pi `$, which we do as follows:
$$\psi (𝐫)\sigma _1^{x+y}\psi (𝐫).$$
(13)
In the new basis, the HF equation reads
$$\sigma _1t\underset{\mathrm{\Delta }𝐫}{}\psi (𝐫+\mathrm{\Delta }𝐫)U[\sigma _3m(𝐫)n(𝐫)/2]\psi (𝐫)=E\psi (𝐫),$$
(14)
where mean-field parameters $`m(𝐫)`$ and $`n(𝐫)`$ are the staggered spin and the charge density.
To simplify further analysis, we neglect density fluctuations and set $`n(𝐫)=n`$ in Eq. (14), also shifting $`EE+Un/2`$. The resulting HF equation,
$$\sigma _1t\underset{\mathrm{\Delta }𝐫}{}\psi (𝐫+\mathrm{\Delta }𝐫)\sigma _3Um(𝐫)\psi (𝐫)=E\psi (𝐫),$$
(15)
acquires a charge conjugation symmetry
$$\psi (𝐫)\sigma _2\psi ^{}(𝐫).$$
(16)
In this particle–hole symmetric form, the mean-field equations resemble those in the theory of polyacetylene. The discrete symmetry (16) has important implications, among them the possibility of spin-charge separation.
When staggered magnetization is uniform, the HF Hamiltonian (15) can be readily diagonalized using the momentum basis:
$`E_𝐤=\pm \sqrt{ϵ_𝐤^2+U^2m^2},`$
where $`ϵ_𝐤=2t(\mathrm{cos}k_x+\mathrm{cos}k_y)`$. The one-electron energy spectrum has a gap $`2\mathrm{\Delta }_0=2U|m|`$, which makes the system an insulator at half-filling. Note that, thanks to the spin axis rotation (13), there is no doubling of the unit cell and the Brillouin zone is therefore not folded.
### B Straight domain wall along $`x`$
A domain wall implies a change of sign for staggered magnetization $`m(𝐫)`$ along a line on the lattice. Because the “local gap” $`U|m(𝐫)|`$ is reduced on this boundary, one expects to find electron states inside the Hubbard gap $`\mathrm{\Delta }_0`$. In the case of an isolated straight uniform domain wall, such states are localized across the wall and extended along it. In view of the charge conjugation symmetry (16), Eq. (15) has an even number of midgap states for every momentum $`\pi <k_x\pi `$, normally two.
Depending on the symmetry of the domain wall—it can be site or bond-centered—the two midgap bands intersect at $`k_x=\pm \pi /2`$ or are separated by a smaller gap $`2\mathrm{\Delta }_1`$, respectively (Fig. 6). Qualitatively, this can be understood by picturing electrons in the domain-wall bands as a 1D electron gas in an external magnetic field $`Um(x,y)=B(y)\mathrm{cos}(\pi x)`$, which is either antisymmetric (site-centered), or symmetric (bond-centered) in $`y`$. In the former case, $`B(y)`$ averaged across the wall vanishes and we deal with effectively free 1D electrons. In the latter case, the electrons feel a nonzero magnetic field staggered along $`x`$, which induces the smaller gap $`2\mathrm{\Delta }_1`$ near $`k_x=\pm \pi /2`$. A more comprehensive discussion can be found in Appendix C.
An undoped antiferromagnet is particle–hole symmetric, therefore only a half of the midgap states are filled. Since the total number of midgap states on a wall of length $`L`$ is $`2L`$, the maximum linear density of holes or electrons that can be doped is $`\nu =1`$. Once all the midgap states are filled, the system is again an insulator.
Stability of the “magic” filling $`\nu =1`$ can be explained by the following qualitative argument. There is only one energy scale $`\mathrm{\Delta }_0`$ in the limit of weak coupling. The preferred filling is determined by minimizing the total energy cost—including the energy of domain walls—per doped particle. When $`\nu <1`$, holes are doped into a midgap band, where one-particle energies are much less than the Hubbard gap $`\mathrm{\Delta }_0`$—at least in the weak-coupling limit. Therefore, the main part of the energy cost comes from creating a domain wall, which should be of order $`\mathrm{\Delta }_0`$ per unit length:
$`\epsilon (\nu ){\displaystyle \frac{E}{\nu L}}{\displaystyle \frac{\alpha \mathrm{\Delta }_0}{\nu }},\nu <1,`$
$`\alpha `$ is of order 1 (in 1D, $`\alpha =2/\pi `$). Doping beyond $`\nu =1`$ puts holes into the lower Hubbard band, separated by the gap $`\mathrm{\Delta }_0`$, at a cost
$`\epsilon (\nu )={\displaystyle \frac{\alpha \mathrm{\Delta }_0+(\nu 1)\mathrm{\Delta }_0}{\nu }},\nu >1.`$
At $`\nu =1`$, the energy per doped hole has a cusp, where the derivative $`\epsilon ^{}(\nu )`$ jumps from $`\alpha \mathrm{\Delta }_0`$ to $`(1\alpha )\mathrm{\Delta }_0`$. The cusp is actually a minimum: since stripes with $`\nu =1`$ are known to exist, they must have a lower energy per doped hole than the uniform AF state, $`\epsilon (1)<\mathrm{\Delta }_0`$, i.e., $`\alpha <1`$.
The minimum of $`\epsilon (\nu )`$ may shift to a lower filling if there are two energy scales for carriers on a stripe or if a fairly large gap (comparable to $`\mathrm{\Delta }_0)`$ opens up in the 1D band at some value of $`\nu <1`$. An example of the former scenario, described by Nayak and Wilczek, gives a smooth minimum at an incommensurate filling determined by the ratio of the two energy scales; this yields conducting stripes. In the latter case, a large gap can possibly be the result of a commensurate filling $`\nu _0<1`$, producing a cusp in $`\epsilon (\nu )`$; such stripes will likely be insulating. Neither argument provides a convincing explanation for experimentally observed conducting stripes with the “magic” filling fraction $`\nu =1/2`$.
### C Continuum formulation
In the weak-coupling limit, $`U4t`$, the characteristic distances are large, and a continuum approximation of some sort should provide a sufficiently accurate description of the system. The approximation must be intelligent enough to tell apart, say, bond-centered and site-centered domain walls, which, as we have seen, have quite different one-particle spectra. The difference comes from a change in the symmetry of the domain wall, and it should be possible to describe within a continuum theory.
To construct an effective continuum approximation for describing a weakly-deformed domain wall, we first have to identify the relevant modes. In a weakly-coupled Hubbard model near half-filling, all low-lying excitations are concentrated near the Fermi-lines $`|k_x\pm k_y|=\pi `$. A domain wall induces a 1D midgap electron band, which is half-filled ($`k_F=\pi /2`$) if the domain wall is not doped. The relevant modes for describing a weakly-deformed domain wall should be located near the intersection of the two Fermi-lines, which gives four “Fermi points” $`𝐤=(\pm \pi /2,\pm \pi /2)`$. Together with the spin index, this implies that the continuum description should be formulated in terms of $`8`$-component Fermion wavefunctions. An alternative, more quantitative way to reach the same conclusion is presented in Appendix C, where a straight domain wall on the lattice is analyzed.
We write an electron wavefunction as a sum of four terms with smoothly varying amplitudes:
$$\psi _s(𝐫)\underset{\alpha =\pm 1}{}\underset{\beta =\pm 1}{}\psi _{\alpha \beta s}(𝐫)e^{i\pi (\alpha x+\beta y)/2}.$$
(17)
\[As there is no folding of the Brillouin zone in our formalism, points $`(\pi /2,\pi /2)`$ and $`(\pi /2,\pi /2)`$ are not equivalent.\] In Eq. (17), we have added two more indices, $`\alpha =\mathrm{sgn}k_x`$ and $`\beta =\mathrm{sgn}k_y`$ to the staggered spin index $`s`$ \[Eq. (13)\] . Only those Fourier components of magnetization which connect the four Fermi patches are preserved:
$`s_3(𝐫){\displaystyle \underset{\alpha =0}{\overset{1}{}}}{\displaystyle \underset{\beta =0}{\overset{1}{}}}m_{\alpha \beta }(𝐫)e^{i\pi (\alpha x+\beta y)}.`$
Each index has its own set of Pauli matrices, $`\{\tau _i^x\}`$, $`\{\tau _i^y\}`$, and $`\{\sigma _i\}`$ for $`\alpha `$, $`\beta `$, and $`s`$, respectively ($`i=1,2,3`$). Any two operators from different sets commute with each other because they act on different indices. Some of these operators have a transparent physical meaning:
$`𝐤{\displaystyle \frac{\pi }{2}}(\tau _3^x,\tau _3^y),(1)^x=\tau _1^x,(1)^y=\tau _1^y,`$ (18)
$`s_1={\displaystyle \frac{\sigma _1}{2}},s_2={\displaystyle \frac{\sigma _2e^{i𝐐𝐫}}{2}},s_3={\displaystyle \frac{\sigma _3e^{i𝐐𝐫}}{2}}.`$ (19)
Using this notation we write the continuum approximation of the HF Hamiltonian (15) as
$`H_{\mathrm{HF}}=2ita\sigma _1(\tau _3^x_x+\tau _3^y_y)U\sigma _3m(𝐫),`$ (20)
$`mm_{11}+m_{01}\tau _1^x+m_{10}\tau _1^y+m_{00}\tau _1^x\tau _1^y,`$ (21)
where $`a`$ is the lattice constant. Only $`m_{11}(𝐫)`$, staggered magnetization proper, exists in the bulk of the antiferromagnet inducing the Hubbard gap $`\mathrm{\Delta }=U|m_{11}(\mathrm{})|`$. The energy spectrum near the 4 points $`(\pm \pi /2,\pm \pi /2)`$ has the form
$`E^2=4t^2a^2(p_x\pm p_y)^2+U^2m_{11}^2.`$
The other three components of $`m`$ (e.g., the average spin density $`m_{00}`$) can be induced around defects only. As we will show, states localized on defects are particularly sensitive to these components.
### D Midgap spectrum of a straight wall
To warm up, let us derive the midgap spectrum of a straight domain wall along the $`x`$ direction. Translational invariance requires that $`m_{01}=m_{00}=0`$. From lattice solutions (Appendix C) we know that the wall fermions have a gapless (gapped) energy spectrum for a site-centered (bond-centered) domain wall. The absence of a gap can be demonstrated by finding a fermion mode with zero energy at $`p_xi_x=0`$ (i.e., $`k_x=\pm \pi /2`$). The Schrödinger equation (20) for $`E=0`$ reads
$$\sigma _2\frac{d\psi }{dy}=\frac{U}{2ta}[\tau _3^ym_{11}(y)+i\tau _2^ym_{10}(y)]\psi (y).$$
(22)
If we neglect $`m_{10}`$ at first, solutions of Eq. (22) are eigenstates of $`\sigma _2`$, $`\tau _3^y`$ and, e.g., $`\tau _3^x`$ (then each zero mode comes from a single Fermi patch). Eq. (22) reduces to 8 uncoupled scalar equations, giving a total of 8 linearly independent solutions. As usual, only half of them \[those with eigenvalues $`\sigma _2\tau _3^ym_{11}(+\mathrm{})<0`$\] are localized on the wall, so that there are 4 zero modes.
Remarkably, in addition to the usual twofold spin degeneracy, there is another spin-like degree of freedom, which will prove to be the transverse flavor. The origin of isospin (at weak coupling) is thus exposed: compared to a 1D chain, there are twice as many “Fermi points” on a straight domain wall in 2D—see Fig. 6, right.
The difference between midgap spectra of site-centered and bond-centered walls arises already in the first order in $`m_{10}`$. Eq. (22) has four zero modes if $`m_{10}`$ is an odd function of $`y`$, i.e., for a site-centered wall. Otherwise, the middle band is split by a small gap (Fig. 6, left).
#### 1 Site-centered domain wall
In the presence of a nonvanishing $`m_{10}`$, the spectrum $`E(p_x)`$ of a straight site-centered stripe can be determined approximately by starting with $`p_x=0`$ \[Eq. (22)\] and treating the term containing $`_x`$ in the Hamiltonian (20) perturbatively. In the limit $`p_x0`$, states outside the main gap can be neglected, which reduces the Hilbert space to the four zero modes \[Eq. (22)\]. By using the degenerate perturbation theory, we find a Dirac spectrum (dashed lines in Fig. 6, left): as $`p_x0,`$
$$E(p_x)\pm vp_x,v=2ta\tau _1^y.$$
(23)
This compares well with a similar result (C9) obtained on a lattice.
#### 2 Bond-centered domain wall
Alternatively, one can explore the limit of a small $`x`$-staggered magnetization, $`m_{10}m_{11}(\mathrm{})`$. In that case, by starting with $`m_{10}=0`$, one finds 4 degenerate zero modes for any $`p_x`$. A nonzero $`m_{10}`$ induces a splitting of the zero modes. To lowest order in $`m_{10}`$, the energy at $`p_x=0`$ is
$`E(0)=\pm \mathrm{\Delta }_1=\pm Um_1\pm U{\displaystyle m_{10}(y)\psi ^{}(y)\psi (y)𝑑y}.`$
As claimed, the gap is proportional to the $`x`$-staggered magnetization felt by an electron on the domain wall.
### E Zero modes at a wiggle
One way to prove that a charged soliton at a wiggle indeed has zero spin is to show that the HF equations contain a doubly degenerate fermion zero mode. An elementary discussion of the connection between zero modes and separation of spin and charge is given in Appendix B. Because the problem is essentially two-dimensional (the domain wall is curved), it is much harder than its 1D analogs. In 1D, symmetry arguments are generally sufficient to prove the existence of zero modes in 1D—even on a lattice! (See Appendix C.) In contrast, we have not been able to find such a general proof for the 2D problem of a wall with a wiggle—neither on the lattice, nor in the long-wavelength approximation.
Instead, we offer a somewhat hand-waving argument in favor of zero modes in this case. Lack of rigor is compensated by an insight into the origin of the transverse flavor: it turns out that holons residing on transverse kinks and antikinks come from different points of the Brillouin zone, i.e., they are made of completely different stuff.
As illustrated in Fig. 7, magnetization on a bond-centered wall with a wiggle can be obtained by superimposing $`m(𝐫)`$ of a straight site-centered wall and that of a spin chain with a kink in $`x`$-staggered magnetization. Away from the wiggle, $`m_{00}(𝐫)=m_{01}(𝐫)=0`$. To simplify the discussion, we will neglect these components altogether (but this is nonessential and can be remedied). Decompose $`m(𝐫)`$ into an $`x`$-independent part and the rest:
$`m(𝐫)`$ $`=`$ $`m^{(0)}(y)+m^{(1)}(𝐫),`$ (24)
$`m^{(0)}(y)=m^{(0)}(y),`$ $`m^{(1)}(\pm \mathrm{},y)=m^{(1)}(\pm \mathrm{},y).`$ (25)
The Hamiltonian (20) can now be split in two parts:
$`H_{\mathrm{HF}}^{(0)}=2ita\sigma _1\tau _3^y_yU\sigma _3[m_{11}^{(0)}(y)+m_{10}^{(0)}(y)\tau _1^y],`$ (26)
$`H_{\mathrm{HF}}^{(1)}=2ita\sigma _1\tau _3^x_xU\sigma _3[m_{11}^{(1)}(𝐫)+m_{10}^{(1)}(𝐫)\tau _1^y].`$ (27)
As shown above, the “transverse part” (26) has 4 zero modes for each $`p_x`$. Within this Hilbert space, $`H_{\mathrm{HF}}^{(1)}`$ describes right and left-moving fermions with spin, which see a staggered magnetization
$$m_1(x)=𝑑yu^{}(y)[m_{10}(𝐫)+m_{11}(𝐫)\tau _1^y]u(y),$$
where $`u(y)`$ is a zero mode (22) of Eq. (26). The midgap fermion band acquires a gap of its own, $`\mathrm{\Delta }_1=U|m_1(\mathrm{})|<\mathrm{\Delta }`$, with two zero modes (one for each spin) inside this smaller gap. “Longitudinal” wavefunctions of the two zero modes satisfy the equation
$$\sigma _2\frac{d\psi (x)}{dx}=\frac{U}{2ta\tau _1^y}\tau _3^xm_1(x)\psi (x).$$
(28)
The existence of two holon flavors can now be deduced from Eqns. (22) and (28). The zero modes have a finite norm only if
$$\sigma _2\tau _3^ym_{11}^{(0)}(+\mathrm{})<0,\sigma _2\tau _3^xm_1(+\mathrm{})/\tau _1^y<0.$$
It follows then that the product of eigenvalues
$$\tau _3^x\tau _3^y=\mathrm{sgn}[m^{(0)}(y=+\mathrm{})m_1(x=+\mathrm{})\tau _1^y]$$
(29)
can be identified with the holon isospin $`2\rho `$. This can be seen by extrapolating Eq. (29) to larger values of $`U`$, which reduces the size of holons. We have $`\tau _1^y=(1)^y=(1)^{y_0}`$, where $`y_0`$ is the row number of the chain in Fig. 7. According to Eq. (29), if $`\tau _3^x\tau _3^y=+1`$, spins on the chain and to the right (left) of the wiggle are an extension of the upper (lower) AF domain, as for the $`\rho =+1/2`$ wiggle in Fig. 7. Thus, $`\rho =\tau _3^x\tau _3^y/2`$. This identification is consistent with numerical HF solutions (Fig. 4), where $`\tau _3^x\tau _3^y=\mathrm{sgn}k_x\mathrm{sgn}k_y`$ can be inferred from the orientation of a holon—the perfect nesting of the Fermi surface makes holon wavefunctions cigar-shaped and oriented along a lattice diagonal.
## V Conclusions
We have attempted to infer a set of plausible quantum numbers of low-lying excitations on a partially doped domain wall in a strongly-correlated antiferromagnet by analyzing artificially-created domain walls in simpler systems. Specifically, we have studied quantum numbers of well-separated holes doped into domain walls in the $`t`$$`J`$ model with Ising anisotropy and in the Hubbard model (in a Hartree-Fock approximation). In addition to a usual 1D electron gas, we have identified a new potential candidate: the 1D gas of holons. In this phase, which we have found at sufficiently small linear hole density $`\nu `$ in both models, charge carriers (holons) have spin $`S_3=0`$ and charge $`Q=+1`$. Each holon resides on a transverse kink of the domain wall, which leads to a strong interplay between charges and transverse fluctuations of a stripe.
We find it very encouraging that the charge carriers with identical quantum numbers result from two vastly different calculations. In the strongly coupled $`t`$$`J_z`$ model, the holons are small and immediately evident \[Fig. 1(d)\]. In the weakly-coupled Hubbard model, they are large and represent a collective effect (fermion zero modes). This indicates that a universality of some sort is at play, and, therefore, that the same “$`\nu 0`$” phase could result from more authentic models. Whether or not this phase is relevant for the cuprate stripes, which have $`\nu 1/2`$, remains an open question. In future, we intend to extend this work to intermediate linear hole densities.
## Acknowledgments
The authors thank C. Chamon, M. M. Fogler, S. A. Kivelson, F. Wilczek, and J. Zaanen for valuable discussions. Part of this work has been done at the Aspen Center for Physics. Financial support from DOE Grant DE-FG02-90ER40542 is gratefully acknowledged.
## A More on the $`t`$-$`J_z`$ model
In this appendix we give some estimates of the energies of domain walls in the strongly anisotropic $`t`$-$`J`$ model (2). In particular, we show that in this limit doped holes can lower their energy by forming (fully-packed) domain walls. This possibility has not been considered in Ref. . We also discuss stability of the dilute holon gas considered in Sec. II.
In the limit of infinite $`J_z`$, only fully doped domain walls, similar to those considered by Osman et al., have a finite energy (all frustrated bonds are covered with holes). A domain wall must therefore maintain continuity, i.e., adjacent holes must be nearest or next-nearest neighbors. This implies that a domain wall horizontal on average, can change its height $`y(x)`$ by at most one unit at a time. Such a domain wall can be fully described using a one-dimensional language, namely by specifying differences in the heights of neighboring holes:
$`y(x+1)y(x)=1,0,\text{ or }1.`$
The system is therefore equivalent to a spin-$`1`$ chain, with the Hamiltonian
$$H_{\mathrm{eff}}=\underset{n}{}[3E_bE_b\left(S_n^z\right)^2+\frac{t}{2}(S_n^+S_{n+1}^{}+\mathrm{h}.\mathrm{c}.)],$$
(A1)
where $`n=x+1/2`$ and $`E_b=VJ_z/4`$ is the energy of an AF bond. The first term in Eq. (A1) represents the energy of a straight segment of a wall (three broken bonds per hole), the second term counts the number of additional broken bonds due to kinks, while the last term describes the transverse hops of holes.
The properties of the effective spin Hamiltonian (A1) and its generalizations have been extensively studied. When $`E_bt`$, the AF ordering $`S_n^z=(1)^n`$ (a zigzag wall) is favored. For $`E_b`$ large and negative the ground state corresponds to $`S_n^z=0`$ (a flat domain wall). At $`|E_b|t`$, the system enters an intermediate critical phase, in which the density of domain wall kinks varies continuously.
Up to terms of higher order in $`t/|E_b|`$, the energies (per hole) of the two ordered phases are
$$\epsilon _{\mathrm{flat}}=3E_bt^2/|E_b|,\epsilon _{\mathrm{zigzag}}=4E_bt^2/E_b.$$
(A2)
For $`E_b<0`$ (two holes attract), phase separation in the bulk affords a lower energy, $`2E_b<\epsilon _{\mathrm{flat}}`$, thus hindering natural formation of stripes. For a strongly repulsive bond, $`E_b|t|`$, stripes win over a lump of immobilized holes: transverse fluctuations of a zigzag stripe lower its kinetic energy by an amount of order $`t`$ per hole \[cf. Eq. (3)\]. This possibility has not been considered in Ref. .
Let us now consider a domain wall created artificially (e.g., by boundary conditions at infinity) in order to study partially doped domain walls. In the absence of holes, such a wall is bond-centered and straight to minimize the number of broken bonds \[see Fig. 1a\]. The corresponding energy cost per unit length is $`J_z/2`$.
When holes are added to an antiferromagnet with a domain wall, they will necessarily bind to it (each hole on the domain wall reduces the number of frustrated spins; the energy gain is $`J_z/2`$ per hole.) As discussed in Sec. II, a single doped hole acquires mobility by riding a kink \[See Fig. 1d\]. The corresponding energy is given by Eq. (6). Assuming that holons are well separated, the energy per added charge is
$$\epsilon _{\mathrm{holon}\mathrm{gas}}=4E_b2t+𝒪(t^2/J_z,\nu t).$$
(A3)
Here we have not included the energy cost of creating a domain wall.
The assumption of large separation may be violated even at very small linear densities $`\nu 1`$ if there is an attractive interaction between holes. For example, for $`E_bt`$, the energy (A2) of the fully-packed flat phase can be smaller than that of the holon gas, (A3). The holes on the stripe will separate into a dense phase ($`\nu =1`$) leaving a portion of the domain wall comletely undoped ($`\nu =0`$).
To find a strict upper bound $`E_b^{\mathrm{min}}`$, such that for $`E_b<E_b^{\mathrm{min}}`$ the phase separation definitely happens, we can use a variational estimate for the ground state energy of the Hamiltonian (A1). The simplest estimate corresponds to all $`S_n^z=0`$, which immediately gives $`\epsilon <3E_b`$. This is smaller then the energy of a dilute holon gas (A3) for $`E_b<E_b^{\mathrm{min}}=2t`$. This estimate of the phase separation boundary $`E_b^{\mathrm{min}}`$ can be easily improved (increased) by using more sophisticated variational wavefunctions.
On the other hand, phase separation of this sort is not expected for $`E_b>0`$, when holes have uniformly repulsive interactions. This statement can be made more formal by evaluating a strict lower bound on the energy of any dense hole phase described by the Hamiltonian (A1). Estimating each term in the Hamiltonian independently, we have, for $`E_b>0`$,
$$H_{\mathrm{eff}}>\underset{n}{}\left(4E_b2|t|\right)=N_h\epsilon _{\mathrm{holon}\mathrm{gas}}.$$
The inequality is strict because the terms in the original Hamiltonian do not commute. It implies that a dilute holon gas is stable to phase separation to a completely doped region and a region of an undoped stripe, as expected on physical grounds for a repulsive interaction.
## B Zero modes and separation of spin and charge
We retrace the relation between fermion zero modes and separation of spin and charge. For completeness, we will use the lattice version of the HF equations (15),
$$\sigma _1t\underset{\mathrm{\Delta }𝐫}{}\psi (𝐫+𝚫𝐫)\sigma _3Um(𝐫)\psi (𝐫)=E\psi (𝐫).$$
(B1)
### 1 Symmetries
So long as we deal with collinear AF configurations, one component of spin—$`(\sigma _3/2)\mathrm{exp}(i𝐐𝐫)`$ in our notation—is a conserved quantity. The transformation
$$\psi (𝐫)\sigma _3e^{i𝐐𝐫}\psi (𝐫)$$
(B2)
is a symmetry of the mean-field Hamiltonian. Here $`𝐐=(\pi ,\pi )`$.
The unitary part of charge conjugation (16),
$$\psi (𝐫)\sigma _2\psi (𝐫),$$
(B3)
alters the sign of $`E`$ and as such is not a symmetry of the mean-field equations. Rather, it can be referred to as a symmetry of zero modes.
For the sake of convenience, we also want the system to be reasonably symmetric with respect to some sort of a parity transformation $`𝐫𝐫`$. If, however, the system has a domain wall, staggered magnetization will be antisymmetric under parity. To make it a symmetry, we combine parity with a spin flip:
$$\psi (𝐫)\sigma _1P\psi (𝐫)\sigma _1\psi (𝐫).$$
(B4)
This “combined parity” is a symmetry of HF equations, provided that $`m(𝐫)=m(𝐫)`$.
### 2 Fermion zero modes
Under some circumstances, Eq. (B1) has solutions with zero energy. Such solutions are normally localized on a topological defect, e.g., on a domain wall. Note that a straight domain wall confines a fermion zero mode in one direction only—across the wall. Therefore, a soliton of finite dimensions requires a domain wall with an inhomogeneity of some sort, e.g., a wiggle. Here we will assume that the wavefunction of a zero mode decays quickly enough in all directions.
Zero fermion modes always come in doublets. This is essentially a consequence of the charge conjugation symmetry (B3) at $`E=0`$. More formally, by starting with the symmetries of zero modes (B3) and (B4), we can construct a triplet of SU(2) generators
$$S_2=\frac{\sigma _2}{2},S_3=\frac{\sigma _3e^{i𝐐𝐫}}{2},S_1=i[S_2,S_3]=\frac{\sigma _1e^{i𝐐𝐫}}{2}.$$
(B5)
Components of $`\stackrel{}{S}`$ are conserved quantities for a zero mode. By inspection, $`\stackrel{}{S}\stackrel{}{S}=3/4`$, i.e., a zero mode is a doublet ($`S=1/2`$). Physically, $`S_3`$ is the component of the total spin of the system parallel to staggered magnetization and is therefore a conserved quantum number. $`S_1`$ and $`S_2`$ are components of the total staggered spin. They generate staggered rotations of spins and are conserved for zero mode fermions, but not for bulk states.
### 3 Spinons and holons
We consider a system with only one zero-mode doublet localized on a topological defect at the origin, $`𝐫=0`$. We assume that the system is symmetric, i.e., $`m(𝐫)=m(𝐫)`$.
Half-filled system. First let us show that a half-filled system has a uniform density. States with $`E<0`$ are filled, while those with $`E>0`$ are empty. One of the two zero modes is filled. Because the density contribution of the zero mode $`\psi _0^{}(𝐫)\psi _0(𝐫)`$ is invariant under rotations generated by $`\stackrel{}{S}`$, we can consider the situation when the mode with $`S_3=1/2`$ is occupied without loss of generality. Then the operator $`\sigma _2`$ toggles between occupied and unoccupied states. The expectation value for the density is
$`n(𝐫)`$ $`=`$ $`{\displaystyle \underset{\psi }{\overset{\mathrm{occ}}{}}}\psi ^{}(𝐫)\psi (𝐫)={\displaystyle \underset{\psi }{\overset{\mathrm{occ}}{}}}\psi ^{}(𝐫)\sigma _2^2\psi (𝐫)`$ (B6)
$`=`$ $`{\displaystyle \underset{\psi }{\overset{\mathrm{unocc}}{}}}\psi ^{}(𝐫)\psi (𝐫)={\displaystyle \frac{1}{2}}{\displaystyle \underset{\psi }{\overset{\mathrm{all}}{}}}\psi ^{}(𝐫)\psi (𝐫)=1.`$ (B7)
Thus a half-filled system has a uniform charge density. It will be shown below that it contains a charge-0, spin-$`1/2`$ soliton at its center.
Half-filled system $`\pm 1`$ electron. When a single electron or hole is added, the density distribution $`n(𝐫)`$ is determined by the profile of the zero mode $`\psi _0^{}(𝐫)\psi _0(𝐫)`$. Because the wavefunction is localized around a defect at the center, we find a soliton with charge $`\pm 1`$.
At the same time, the soliton has zero net spin $`S_3`$. Moreover, $`S_3=0`$ in any symmetric finite area around the defect. This happens because contributions to the total spin from $`𝐫`$ and $`𝐫`$ cancel each other:
$`s_3(𝐫)=(1)^{x+y}m(𝐫)=(1)^{x+y}m(𝐫)=s_3(𝐫).`$
Thus, at half-filling $`\pm 1`$ electron, the system has a soliton with charge $`Q=1`$ and spin $`S_3=0`$. Exactly at half-filling, the soliton has $`Q=0`$ and $`S_3=\pm 1/2`$. These are, respectively, a holon and a spinon.
## C Straight domain wall: a lattice analysis
Consider a straight domain wall along the $`x`$ axis. At a given lattice momentum $`k_x`$ along the wall, the mean-field Hamiltonian (15) reads
$`\sigma _1t[\psi (y+1)+\psi (y1)+2\mathrm{cos}k_x\psi (y)]`$ (C1)
$`\sigma _3Um(y)\psi (y)=E\psi (y).`$ (C2)
For definiteness, the domain wall is located on the line $`y=0`$, so that $`m(y)=m(y)`$. The site indices are integer ($`0,\pm 1,\mathrm{}`$) for a site-centered wall and half-integer ($`\pm 1/2,\pm 3/2,\mathrm{}`$) for a bond-centered wall.
We will first show that the smaller gap (separating the midgap bands) is absent in the case of a site-centered domain wall, for which $`m(y)=m(y)`$. It suffices to show that there exists a zero mode at $`k_x=\pm \pi /2`$.
### 1 Site-centered wall: zero modes
Setting $`E=0`$ and $`k_x=\pm \pi /2`$ converts Eq. (15) into
$$\psi (y+1)+\psi (y1)=i\sigma _2\mathrm{\Delta }(y)\psi (y),$$
(C3)
where $`\mathrm{\Delta }(y)=Um(y)/t`$. In the bulk,
$`\mathrm{\Delta }(+\mathrm{})=\mathrm{\Delta }(\mathrm{})=\mathrm{\Delta }_0>0.`$
We are looking for a finite-norm solution to Eq. (C3). Evidently, $`\sigma _2`$ can be immediately diagonalized: $`\sigma _2=\pm 1`$. The substitution
$`\psi (y)=\varphi (y)e^{i(\pi /2)\sigma _2y}\chi `$
where $`\chi `$ is an arbitrary constant spinor, yields a difference equation with real coefficients for a scalar $`\varphi (y)`$:
$$\varphi (y+1)\varphi (y1)=\mathrm{\Delta }(y)\varphi (y).$$
(C4)
Eq. (C4) has two linearly independent solutions with the asymptotic behaviors
$`\varphi _\pm (y)\mathrm{exp}(\pm \kappa y)\text{ as }y+\mathrm{}.`$
The spread of the wavefunction $`\kappa ^1`$ is given by the equation
$`2\mathrm{sinh}\kappa =\mathrm{\Delta }_0.`$
Clearly only $`\varphi _{}(y)`$ can be normalizable — provided that it is also well behaved as $`y\mathrm{}`$. It will be seen shortly that this is indeed the case.
Observe that one can write $`\varphi _{}(y)`$ as a superposition of another pair of linearly independent solutions, the eigenstates of parity,
$`\varphi _i(y)=\eta _i\varphi _i(y),`$
where $`i=1,2`$ and $`\eta _i=\pm 1`$. Ordinarily, $`\eta _1\eta _2`$, so that
$`\varphi _{}(y)=c_1\varphi _1(y)+c_2\varphi _2(y)`$
is not, in general, a parity eigenstate. Yet it must be if it is to have a finite norm. Indeed, from different asymptotic forms of the solutions $`\varphi _\pm (y)`$ we infer that $`\varphi _{}(y)=\eta \varphi _{}(y)`$, or else it will diverge at $`y=\mathrm{}`$. Luckily, all solutions of Eq. (C4) have the same (even) parity for a site-centered domain wall: $`\mathrm{\Delta }(0)=0`$ and therefore
$`\varphi (1)\varphi (1)=\mathrm{\Delta }(0)\varphi (0)=0.`$
Since $`\varphi _{}(y)`$ vanishes exponentially at both $`y=+\mathrm{}`$ and $`\mathrm{}`$, it has a finite norm.
We have thus proven that the HF Hamiltonian (C2) has 2 solutions of a finite norm with $`E=0`$ for $`|k_x|=\pi /2`$, one for each eigenvalue of $`\sigma _2`$, if the wall is site-centered. As a rule, there are no zero modes if the domain wall is bond-centered.
The asymptotic behavior of the zero modes as $`y\mathrm{}`$ is
$$\psi (𝐫)\mathrm{exp}\left(\pm i\frac{\pi }{2}xi\frac{\pi }{2}\sigma _2y\kappa |y|\right)\chi ,$$
(C5)
where $`2\mathrm{sinh}\kappa =\mathrm{\Delta }_0`$ and $`\chi `$ is an arbitrary constant spinor. We stress that, in the limit of weak coupling $`\kappa 1`$, the Fourier components of $`\psi (𝐫)`$ come from the 4 “Fermi patches” $`|k_x||k_y|\pi /2`$. This means twice as many Fermi points as in an ordinary 1D electron gas! It is for this reason that holons (and spinons) on a domain wall come in two flavors.
### 2 Site-centered wall: 1D electron band
For $`k_x\pi /2`$, the degeneracy of the two midgap states is lifted by the additional term $`2t\sigma _1\mathrm{cos}k_x`$ in the HF Hamiltonian (C2). Near $`k_x=\pm \pi /2`$, this term can be treated as a perturbation acting in the Hilbert space of the two zero modes, allowing us to determine their splitting to first order.
Parametrize the wavefunctions in the familiar way,
$`\psi (x,y)=e^{ik_xx}\varphi (y)e^{i(\pi /2)\sigma _2y}\chi ,`$
where $`\varphi (y)`$ is the solution of Eq. (C4) with the norm 1. The resulting $`\psi (x,y)`$ diagonalizes the mean-field Hamiltonian (C2) with $`\mathrm{cos}k_x=0`$. The two-fold degeneracy of the zero mode is due to the freedom in choosing the spinor $`\chi `$. In the framework of the degenerate perturbation theory, we compute the matrix element of the perturbation $`H_1=2t\sigma _1\mathrm{cos}k_x`$ between the states $`|k_{x1},\chi _1`$ and $`|k_{x2},\chi _2`$:
$`k_{x2},\chi _2|H_1|k_{x1},\chi _1=2t\mathrm{cos}k_{x1}`$ (C6)
$`\times 2\pi \delta (k_{x1}k_{x2})(\chi _2^{}\sigma _1\chi _1){\displaystyle \underset{y}{}}(1)^y\varphi ^2(y).`$ (C7)
This result is quite tangible: as far as the midgap states are concerned, “integrating out” the transversal degree of freedom $`y`$ yields electrons with nearest-neighbor hopping along the domain wall and no staggered magnetization — cf. Eq. (15):
$$\sigma _1t_{1\mathrm{D}}[\psi (x+1)+\psi (x1)]=E\psi (x).$$
(C8)
The effective 1D hopping amplitude is
$$t_{1\mathrm{D}}=t(1)^yt\underset{y}{}(1)^y\varphi ^2(y).$$
(C9)
Note that $`t_{1\mathrm{D}}t`$ in the limit of weak coupling. |
warning/0001/cond-mat0001366.html | ar5iv | text | # Anomalous low doping phase of the Hubbard model
## I Introduction
The 1-band Hubbard model on a two-dimensional square lattice has the Hamiltonian
$`H`$ $`=`$ $`t{\displaystyle \underset{i,j,\sigma }{}}(c_{i,\sigma }^{}c_{j,\sigma }+h.c.)+U{\displaystyle \underset{i}{}}(n_{i,}{\displaystyle \frac{1}{2}})(n_{i,}{\displaystyle \frac{1}{2}})`$ (2)
$`\mu {\displaystyle \underset{i}{}}(n_{i,}+n_{i,}).`$
Here, $`c_{i,\sigma }^{}`$ ($`c_{i,\sigma }`$) creates (annihilates) an electron with spin $`\sigma `$ in a Wannier orbital centered at lattice site $`i`$. The particle density at each site is given by $`n_{i,\sigma }=c_{i,\sigma }^{}c_{i,\sigma }`$. The first sum for the kinetic energy is restricted to include only the hopping matrix element $`t`$ between next-nearest neighbor sites $`i,j`$. Periodic boundary conditions are used throughout the following work. The second sum describes for $`U>0`$ an on-site Coulomb repulsion between particles of opposite spin that share the same lattice site. In the present work we restrict ourselves to $`U=8.0t`$. The chemical potential $`\mu `$ in the third sum controls the occupation of the finite lattice in the finite-temperature grand canonical Quantum-Monte-Carlo (QMC) simulation we preformed. At half-filling, particle-hole symmetry of the kinetic and $`U`$-term implies $`\mu =0`$. The analytic continuation of the dynamic imaginary-times QMC data to the real frequency axis is performed with state-of-the-art Maximum-Entropy (ME) techniques. For exhausting discussions concerning the QMC and ME methods we refer the reader to .
The 1-band Hubbard model exhibits several energy scales: in the repulsive case the high-energy scale $`U`$ is important in determining the insulating gap at half-filling, $`n=1.0`$. An important low-energy scale is set by the exchange interaction $`J=4t^2/U`$: in second order perturbation theory two particles with different spins on neighboring lattice sites can exchange via a virtual double occupation. This process is the source for the strong antiferromagnetic (AF) correlations found near and at half-filling.
With the exception of some known symmetry properties like invariance under global spin-rotation, whose generators form the $`SU(2)`$–algebra, and invariance under $`U(1)`$-transformation, i.e., charge-conservation, as well as the particle-hole transformation, no rigorous results are known for the 1-band Hubbard model in two-dimensions . The Mermin-Wagner theorem prevents a long-range ordered state in a two-dimensional system for finite temperatures, but it is commonly believed that the ground state of the spin-$`1/2`$ Heisenberg antiferromagnet, i.e. the large-$`U`$ limit of the repulsive half-filled Hubbard model , shows long-range Néel order in two dimensions. Néel order results also in the weak-coupling limit , where the gap $`\mathrm{\Delta }`$ is due to a spin-density-wave (SDW) instability related to perfect nesting.
It might seem that due to the Mermin-Wagner theorem the physics of the ordered phase is out of reach for our numerical technique, which is limited to finite temperatures. However, one may assume that the system ‘effectively orders’ as soon as the spin-correlation length becomes comparable to the system size. Due to the periodic boundary conditions the spin-correlation function $`\chi (𝒓)=(1/L^2)_i\stackrel{}{S}_i\stackrel{}{S}_{i+𝒓}`$ is a periodic function of $`𝒓`$ and if the value of this function at the maximum value $`|𝒓|=\sqrt{2L}`$ (with $`L`$ $``$ cluster size) is still appreciable, we may expect that the system is ‘effectively ordered’. A rough measure would be the spin-correlation length $`\zeta `$, obtained by fitting $`\chi (𝒓)`$ to the form $`a|𝐫|^be^{|𝐫|/\zeta }`$. Since the infinite system has the Néel temperature $`T=0`$, the spin-correlation length diverges as $`T0`$ and we may expect that in a finite system $`\zeta `$ becomes comparable to the cluster size $`L`$ at a finite temperature which depends on the lattice size. Below this temperature we then expect that the system resembles the ordered phase, although the spin-rotation symmetry persists even in this case. In that sense, the finite size of the system creates an artificial ‘Néel temperature’ which, however, depends on $`L`$ and $`U/t`$ etc and has no real counterpart in the infinite system. This has to be kept in mind when discussing the results.
We proceed by discussing some known approximations to the Hubbard model. The ‘classical’ approximation to the Hubbard model is the so-called Hubbard-I approximation. Its essence is the splitting of the electron annihilation operator into the two ‘eigenoperators’ of the interaction part:
$`c_{i,\sigma }`$ $`=`$ $`c_{i,\sigma }n_{i,\overline{\sigma }}+c_{i,\sigma }(1n_{i,\overline{\sigma }})`$ (3)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(d_{i,\sigma }+h_{i,\sigma }^{}),`$ (4)
whence
$`[d_{i,\sigma },H_U]`$ $`=`$ $`{\displaystyle \frac{U}{2}}d_{i,\sigma },`$ (5)
$`[h_{i,\sigma }^{},H_U]`$ $`=`$ $`{\displaystyle \frac{U}{2}}h_{i,\sigma }^{}.`$ (6)
The physical content of the Hubbard-I approximation, which neglects any spin-correlations, becomes clear by realizing that the equations of motion in this approximation are completely equivalent to an ‘effective Hamiltonian’ for double occupancy-like particles $`d_{i,\sigma }^{}=(1/\sqrt{2})c_{i,\sigma }^{}n_{i,\overline{\sigma }}`$ and hole-like particles $`h_{i,\sigma }^{}=(1/\sqrt{2})c_{i,\sigma }(1n_{i,\overline{\sigma }})`$:
$`H`$ $`=`$ $`{\displaystyle \frac{t}{2}}{\displaystyle \underset{i,j,\sigma }{}}(d_{i,\sigma }^{}d_{j,\sigma }h_{i,\sigma }^{}h_{j,\sigma })`$ (8)
$`{\displaystyle \frac{t}{2}}{\displaystyle \underset{i,j,\sigma }{}}(d_{i,\sigma }^{}h_{j,\overline{\sigma }}^{}+h.c.)`$
$`+`$ $`{\displaystyle \frac{U}{2}}{\displaystyle \underset{i,\sigma }{}}(d_{i,\sigma }^{}d_{i,\sigma }h_{i,\sigma }^{}h_{i,\sigma }).`$ (9)
This Hamiltonian contains terms which describe the pair creation of a hole and a double occupancy on nearest neighbors $`i,j`$, terms which describe the propagation of these effective particles, and an additional energy of formation of $`U`$ for the double occupancy. The matrix elements for the propagation are reduced by a factor of $`1/2`$, because in an uncorrelated spin-background there is a probability of $`1/2`$ for the spin on a nearest neighbor to have the proper direction to allow for the hopping of an electron. Solving (9) by Fourier and Bogoliubov transform:
$`\gamma _{1,𝒌,\sigma }`$ $`=`$ $`u_𝒌d_{𝒌,\sigma }+v_𝒌h_{𝒌,\overline{\sigma }}^{},`$ (10)
$`\gamma _{2,𝒌,\sigma }`$ $`=`$ $`v_𝒌d_{𝒌,\sigma }+u_𝒌h_{𝒌,\overline{\sigma }}^{},`$ (11)
yields the standard dispersion relation, which consists of the upper and lower Hubbard bands:
$$E_\pm ^{HubI}(𝐤)=\frac{1}{2}\left(ϵ_𝒌\pm \sqrt{ϵ_𝒌^2+U^2}\right).$$
(12)
Here $`ϵ_𝒌`$ denotes the free tight-binding dispersion $`ϵ_𝒌=2t(\mathrm{cos}(k_x)+\mathrm{cos}(k_y))`$. Using (4) we also obtain the correct spectral weights of the two Hubbard bands:
$`Z_\pm ^{HubI}(𝐤)`$ $`=`$ $`{\displaystyle \frac{1}{2}}(u_𝒌\pm v_𝒌)^2`$ (13)
$`=`$ $`{\displaystyle \frac{1}{2}}\left(1\pm {\displaystyle \frac{ϵ_𝒌}{\sqrt{ϵ_𝒌^2+U^2}}}\right).`$ (14)
As noted above, the key assumption in the Hubbard I approximation is the neglect of spin-correlations. Therefore, we expect that this approximation will become inaccurate as soon as spin-correlations become sufficiently strong so as to appreciably influence the propagation of holes and double occupancies. This effect will be strongest in the Néel ordered phase believed to be realized in the ground state. If we choose the ‘spin-background’, in which the holes and double occupancies propagate, to be the Néel state, the double occupancy $`d_i^{}`$ with spin $``$ can exist only on the $``$ sublattice and vice versa. Similarly, a hole $`h_i^{}`$ can be created only on the $``$ sublattice and vice versa. We, thus, expect that we have to modifiy the Hubbard-I Hamiltonian (9) into:
$`H`$ $`=`$ $`t{\displaystyle \underset{iA,jN(i)}{}}(d_{i,}^{}h_{j,}^{}+h.c.)`$ (16)
$`t{\displaystyle \underset{iB,jN(i)}{}}(d_{i,}^{}h_{j,}^{}+h.c.)`$
$`+`$ $`{\displaystyle \frac{U}{2}}{\displaystyle \underset{i,\sigma }{}}(d_{i,\sigma }^{}d_{i,\sigma }h_{i,\sigma }^{}h_{i,\sigma }),`$ (17)
where $`A`$ ($`B`$) denotes the $``$ ($``$) sublattice and $`N(i)`$ denotes the set of nearest neighbors of site $`i`$. Note that the $`d_{i,\sigma }^{}d_{j,\sigma }`$ and $`h_{i,\sigma }^{}h_{j,\sigma }`$ propagation terms drop out (because of the $`A`$ and $`B`$ sublattices), and the matrix element for pair creation are $`t`$ rather than $`t/2`$ \- this takes into account the fact that the spins on neighboring sites are antiparallel with probability $`1`$ rather than $`1/2`$ as was the case in the paramagnetic state. Fourier and Bogoliubov transformation of (17) yields the dispersion
$$E_\pm ^{SDW}(𝒌)=\pm \sqrt{ϵ_𝒌^2+\frac{U^2}{4}},$$
(18)
and using
$`c_{𝒌,}`$ $`=`$ $`h_{𝒌,}+d_{𝒌,}^{}`$ (19)
$`c_{𝒌+𝑸,}`$ $`=`$ $`h_{𝒌,}d_{𝒌,}^{}`$ (20)
(where $`𝒌`$ is within the AF Brillouin zone), we find the spectral weight of these bands:
$`Z_\pm `$ $`=`$ $`{\displaystyle \frac{1}{2}}(u_𝒌\pm v_𝒌)^2`$ (21)
$`=`$ $`{\displaystyle \frac{1}{2}}\left(1\pm {\displaystyle \frac{ϵ_𝒌}{\sqrt{ϵ_𝒌^2+\frac{U^2}{4}}}}\right).`$ (22)
This is precisely what is obtained from the SDW mean-field treatment of the Hubbard model by setting the staggered magnetization $`m`$ to a value of $`1`$ (which is a good approximation in the limit of large $`U`$). In general, the SDW mean-field approximation models the AF Néel state by assuming $`n_{i,\sigma }=\frac{1}{2}(1+\sigma me^{i\mathrm{𝐐𝐑}_i})`$ with AF nesting vector $`𝐐=(\pi ,\pi )`$ and staggered magnetization $`m=e^{i\mathrm{𝐐𝐑}_i}n_{i,}n_{i,}`$. This results in the two-band dispersion
$$E_\pm ^{SDW}(𝐤)=\pm \sqrt{ϵ_𝒌^2+\mathrm{\Delta }^2},$$
(23)
and the spectral weight
$$Z_\pm ^{SDW}(𝐤)=\frac{1}{2}\left(1\pm ϵ_𝒌/E_+^{SDW}(𝐤)\right).$$
(24)
The gap parameter is $`\mathrm{\Delta }=Um/2`$ and the staggered magnetization $`m`$ is determined self-consistently from the following equation:
$$m=\frac{2}{N}\underset{𝐤}{\overset{occupied}{}}\frac{Um}{\sqrt{(ϵ_𝒌ϵ(𝐤+𝐐))^2+m^2U^2}}.$$
(25)
The solution of this self-consistency equation yields the value $`\mathrm{\Delta }=3.56t`$ at $`U=8.0t`$, the value used in the present work. If we set $`m=1`$ on the other hand, as would be apropriate for $`U/t\mathrm{}`$, we obviously recover the results from our Hubbard-I-like Hamiltonian (17).
In the two limiting cases of no spin-correlations and of perfect Néel order we can thus treat the Hubbard model in a quite analogous fashion and, as will be seen below, the Hubbard-I results are indeed a good approximation to the actual spectral function in the limit of high temperature. The main problem then is how to describe the effect of spin-fluctuations and how to manage the crossover from the completely disordered to the Néel ordered phase. Below, we will adress this crossover by QMC simulations.
Thereby we want to take advantage of the possibility to calculate the spectra of specifically designed ‘diagnostic operators’. The first one of these is the ‘shadow’ operator:
$`\stackrel{~}{c}_{i,\sigma }`$ $`=`$ $`c_{i,\sigma }n_{i,\overline{\sigma }}c_{i,\sigma }(1n_{i,\overline{\sigma }})`$ (26)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(d_{i,\sigma }h_{i,\sigma }^{}).`$ (27)
We note first of all, that the Fourier transform of this operator, $`\stackrel{~}{c}_{𝒌,\sigma }`$ has precisely the same quantum numbers as the ordinary electron operator $`c_{𝒌,\sigma }`$: momentum $`𝒌`$, total spin $`1/2`$, $`z`$-spin $`\sigma `$, and identical point group symmetry at high symmetry momenta. It follows that the poles in its dynamical correlation functions
$`A_{\stackrel{~}{c}}(𝒌,\omega )`$ $`=`$ $`{\displaystyle \underset{\nu }{}}|\mathrm{\Psi }_\nu |\stackrel{~}{c}_{𝒌,\sigma }|\mathrm{\Psi }_\mu |^2`$ (29)
$`e^{\beta (E_\mu E_\nu )}\delta (\omega (E_\nu E_0))`$
originate from exactly the same final states $`|\mathrm{\Psi }_\nu `$ as those of the photoemission spectrum. It can happen only accidentally that a given state $`\nu `$ has an exactly vanishing weight in one of the correlation functions, but not the others. In this case, however, any arbitrarily small perturbation will remove the accidental vanishing of the peak. The only thing that can and will be different in the spectra of the diagnostic operator, are the weights of the peaks, $`|\mathrm{\Psi }_\nu |\stackrel{~}{c}_{𝒌,\sigma }|\mathrm{\Psi }_0|^2`$. In fact, comparison of equations (22) and (14) shows immediately that both for the Hubbard I approximation and for the SDW-approximation use of the operator $`\stackrel{~}{c}_{𝒌,\sigma }`$ instead of $`c_{𝒌,\sigma }`$ exchanges the weights of the two Hubbard/SDW bands. It follows that band portions which have a small spectral weight in the spectra of the ordinary electron operator will aquire a large spectral weight in the spectrum of $`\stackrel{~}{c}_{𝒌,\sigma }`$ and vice versa. Therefore, this diagnostic operator is a useful tool to map out the ‘shadowy’ parts of the spectra in the outer parts of the Brillouin zone. An additional benefit is that since the ME technique resolves peaks with large spectral weight more reliably than those with small weight, we can get more precise information about the dispersion of these bands with weak intensity. The second diagnostic operator that we will be using is the ‘spin-$`1/2`$ string’ operator
$`b_{i,}`$ $`=`$ $`{\displaystyle \underset{jN(i)}{}}(c_iS_j^z+c_iS_j^{}),`$ (30)
$`b_{i,}^{}`$ $`=`$ $`{\displaystyle \underset{jN(i)}{}}(c_i^{}S_j^z+c_i^{}S_j^+),`$ (31)
where the sum on the r.h.s. runs over the nearest neighbors $`j`$ of site $`i`$. This operator is the Clebsch-Gordan contraction of the adjoint rank-$`1`$ spinor $`c_{i,\sigma }`$ and the spin-$`1`$ vector operator $`\stackrel{}{S}_j`$ into yet another adjoint rank-$`1`$ spinor. It describes the annihilation of a ‘dressed’ electron, i.e., an electron with a spin-excitation on a nearest neighbor. Again, since the Fourier transforms of the diagnostic operator, $`b_{𝒌,\sigma }`$, agrees with the electron annihilation operator $`c_{𝒌,\sigma }`$ in all conceivable quantum numbers (momentum, total spin, $`z`$-component of the spin, point group operations in the case of high symmetry momenta), it obeys precisely the same selection rules, whence it is just as good as the $`c`$-operator itself to map out the band structure.
One further point, which is important from the technical point of view, is the following: under the particle-hole transformation $`c_{i,\sigma }e^{i𝑸𝑹_i}c_{i,\sigma }^{}`$ we have $`\stackrel{~}{c}_{i,}e^{i𝑸𝑹_i}\stackrel{~}{c}_{i,}^{}`$ and $`b_{i,}e^{i𝑸𝑹_i}b_{i,}^{}`$. This implies that at half-filling the spectra of the shadow and spin-$`1/2`$ string operator obey the same particle-hole-symmetry as those of the ordinary electron operator, i.e.
$$A(𝒌,\omega )=A(𝒌+𝑸,\omega ).$$
(32)
In our MaxEnt program particle-hole-symmetry is not implemented as an additional constraint, in other words: the MaxEnt procedure does not ‘know’ about this additional symmetry. The degree to which (32) is obeyed in the final spectra thus gives a good check for the accuracy of the reconstructed spectra. This is of particular importance in the case of the spin-$`1/2`$ string operator because the Wick-contraction of this operator on any given time slice produces a total of $`80`$ products of noninteracting Green’s functions. The computation is, therefore, much more prone to inaccuracies so that an additional check is desirable.
The present work is organized as follows: first, we compare the temperature dependent dynamic single-particle properties of the Hubbard model with the predicitions of the mean-field SDW and Hubbard-I approximations. In addition, we consider the temperature dependent two-particle excitations. Then, we will use our first diagnostic operator, the shadow operator $`\stackrel{~}{c}_{𝒌,\sigma }`$, to show the existence of a total of four bands in the photoemission spectrum and to shed light onto the temperature dependent crossover from the SDW to the Hubbard-I regime. Then we will investigate the 4-band structure in more detail: after a phenomenological fit, which is able to produce a total of four bands, we will consider the string picture which naturally leads to our second diagnostic operator, the spin-$`1/2`$ string operator $`b_{𝒌,\sigma }`$. This operator will be used to ultimately reveal the underlying mechanisms behind the generation of the 4-band structure, namely the dressing of the photoholes by clouds of AF spin-excitations. To resolve the ‘AF mirror image’ of the narrow quasiparticle spectral weight features between $`𝐤=(0,0)`$ and $`𝐤=(\pi ,0)`$ around momentum $`𝐤=(\pi ,\pi )`$ we introduce also a spin-$`3/2`$ string operator. Finally, we will concentrate on the doped regime, thereby showing the violation of the Luttinger theorem near half-filling, and discuss the hole concentration range in which these dressing effects dominante the low-energy physics.
## II Temperature dependent dynamics at half-filling
We start in figure 1 with the discussion of the angle-resolved single-particle spectral function $`A(𝐤,\omega )`$ for $`U=8.0t`$ and various temperatures in the range from $`T=4.00t`$ to $`T=0.10t`$. In the left column of the figure the spectral functions are shown as ‘grey-scale’ plots versus momentum $`𝐤`$ and energy $`\omega /t`$ with dark (light) areas corresponding to large (small) spectral weight. The same spectra are shown also in the right column, but now as line-plots at each momentum $`𝐤`$. The QMC data in the figure are compared to the renormalized results of the Hubbard-I approximation at high and medium temperatures and with the renormalized results of the mean-field SDW approximation at the lowest temperature, $`T=0.10t`$. ‘Renormalized’ means here that we have readjusted the parameters $`U`$ and $`t`$ in (12) and (23)so as to obtain an optimal fit to the ‘bands’ of high spectral weight in the spectra. These approximate dispersions are plotted as solid lines in the left column, while their spectra are shown in the right column as line-plots at each momentum $`𝐤`$. Thereby we have assumed a Lorentzian lineshape with a suitably chosen temperature dependent width (the Hubbard-I approximation does not provide any information about linewidths).
Starting at the highest calculated temperature, $`T=4.00t`$, we find that the Hubbard-I approximation with practically unrenormalized parameters ($`\stackrel{~}{t}=0.95t`$ and $`\stackrel{~}{U}=8.32t`$) fits the QMC spectral functions almost perfectly regarding both the general dispersion and the distribution of spectral weight (quadratic deviation per degree of freedom $`\chi ^2=0.05`$). This is not surprising since the Hubbard-I approximation is derived for the paramagnetic state thereby neglecting all effects of spin-correlations, i.e., the Hubbard-I approximation essentially describes the interplay between itinerant electrons and strong on-site repulsion. This should be a reasonable assumption at this high temperature since all relevant spin-degrees of freedom are thermally excited.
Lowering the temperature to $`T=1.00t`$ and further to $`T=0.33t`$, the data show that the Hubbard-I approximation increasingly fails to reproduce the entire spectrum and is only able to fit the peaks with maximal spectral weight reasonably well. Even then, in order to achieve these fits, one already has to renormalize the free parameters strongly. The values we found are $`\stackrel{~}{t}=1.38t`$ and $`\stackrel{~}{U}=5.57t`$ with $`\chi ^2=1.54`$ at $`T=1.00t`$ and $`\stackrel{~}{t}=1.40t`$ and $`\stackrel{~}{U}=3.96t`$ with $`\chi ^2=0.85`$ at $`T=0.33t`$. The peaks that are missed by the Hubbard-I approximation are the states that form the first ionization/affinity states around momentum $`𝐤=(0,0)/(\pi ,\pi )`$ on the photoemission/inverse photoemission side and two rather dispersionless bands at higher energies of $`\omega \pm 6.0t`$. The former states were previously resolved by Preuss and co-workers . Alltogether one can distinguish a total of four ‘bands’ in the single particle spectral density. As will be seen below, the temperature/doping regime where this 4-band structure is seen coincides with the regime where a collective low-energy mode with momentum $`(\pi ,\pi )`$ in the spin-response exists. Inspite of this, however, we stress that the 4-band structure cannot be explained by a backfolding of the band structure due to ordering effects since the spin-correlation length is $`1.5`$ lattice spacings at temperatures $`T0.33t`$.
For the lowest temperature, $`T=0.10t`$, the QMC data are compared with the results from the AF SDW approximation. As in the case of the Hubbard-I approximation at medium temperatures, the lowermost spectra of figure 1 show that the SDW approximation is only able to fit the peaks with large spectral weight. Again one has to renormalize the free parameters heavily to values of $`\stackrel{~}{t}=1.34t`$ and $`\stackrel{~}{\mathrm{\Delta }}=2.29t`$ with $`\chi ^2=0.83`$. Moreover, as was the case for the Hubbard-I approximation at higher temperatures, the SDW approximation neither explains the states that form the first ionization/affinity states around momentum $`𝐤=(0,0)/(\pi ,\pi )`$ on the photoemission/inverse photoemission side, nor the two dispersionless bands at higher excitation energies which can be seen rather clearly in the spectra.
All in all, the overall distribution of spectral weight is roughly reproduced by the Hubbard-I and SDW approximations as long as one forgets about the 4-band structure. In fact it is well known that the integrated photoemission or inverse photoemission weight (that means the electron momentum distribution) at each $`𝐤`$–point is reproduced quite well by the Hubbard-I approximation and the related $`2`$-pole approximation.
As already mentioned, the emerging of the 4-band structure in the photoemission somewhere in between $`T=1.00t`$ and $`T=0.33t`$ is closely related to a change in the spin-response: to illustrate this we consider figure 2, which shows the spin-correlation function, $`\chi _{sz}(𝐤,\omega )`$ (left column), and the charge-correlation function, $`\chi _{cc}(𝐤,\omega )`$ (right column), for different temperatures. Whereas the spin-response is entirely incoherent at $`T=1.00t`$, with decreasing temperature it can be fitted increasingly well by the spin-wave dispersion
$$E^{SW}(𝐤)=2J\sqrt{1\frac{1}{4}(\mathrm{cos}(k_x)+\mathrm{cos}(k_y))^2}.$$
(33)
This result is known from previous calculations , which demonstrated that the two-particle correlation functions like the spin-response can be described within the SDW approximation even for large values of the interaction $`U`$. The energy scale $`J`$ directly manifests itself in the spin-response since the spin-wave dispersion takes the value of $`E^{SW}(\pi ,0)=2J`$ at momentum $`𝐤=(\pi ,0)`$. The fit parameters are $`\stackrel{~}{J}=0.33t`$ with $`\chi ^2=0.01`$ at $`T=0.33t`$ and $`\stackrel{~}{J}=0.49t`$ with $`\chi ^2=0.11`$ at $`T=0.10t`$. The latter is already quite close to the strong coupling estimate $`J=4t^2/U=0.5t`$. Furthermore, the figure shows that with decreasing temperature the spin-response concentrates its weight more and more at the AF momentum $`𝐐=(\pi ,\pi )`$ (as it is the case in AF spin-wave theory) and at a characteristic energy $`\omega ^{}`$. The latter decreases with decreasing temperature, i.e., the spin-response comes closer and closer to the predictions of AF spin-wave theory (33). The spin-correlation length $`\zeta _{sz}(T)`$ can be derived from a real-space fit of the QMC equal-imaginary-times spin-correlation function $`\chi (𝒓)`$ to the form $`a|𝐫|^be^{|𝐫|/\zeta _{sz}(T)}`$ thereby incorporating the effects of the periodic boundary conditions. While this is the best one can do on a finite lattice with periodic boundary conditions, the fit will only lead to roughly correct values due to the relative small system size of $`8\times 8`$. The values obtained for the spin-correlation length then are $`\zeta _{sz}=0.30.5`$, $`\zeta _{sz}=1.01.3`$, $`\zeta _{sz}=1.61.9`$, $`\zeta _{sz}=2.12.8`$ and $`\zeta _{sz}>8`$ for $`T=1.00t`$, $`T=0.33t`$, $`T=0.25t`$, $`T=0.20t`$ and $`T=0.10t`$, respectively. We actually believe that the spin-correlation length reaches the system size already at a temperature of $`T0.20t`$, because at this temperature the fit results in values between $`\zeta _{sz}=2.1`$ and $`\zeta _{sz}=2.8`$ (with the exponent $`b`$ set to zero or not) but always with error bars of roughly the system size.
The charge-response $`\chi _{cc}(𝐤,\omega )`$, on the other hand, is rather broad in both momentum $`𝐤`$ and energy $`\omega `$ for all temperatures studied. Furthermore, the charge-response is gapped for temperatures below $`T1.00t`$ and, therefore, can certainly not be responsible for any low-energy features of the single particle spectrum.
It is then quite obvious that at roughly the same temperature where the two narrow dispersive quasiparticle-like bands (that cannot be interpreted within the framework of the Hubbard-I or SDW approximations) appear in the single particle spectrum, the spin-response develops a sharp collective low-energy mode. We conclude that the underlying mechanism behind the occurrence of the 4-band structure consists in dynamical magnetic correlation effects, which are beyond the scope of the Hubbard-I and SDW approximations.
In the following, we want to explore the single particle spectrum by means of our diagnostic operators. The first of them is the shadow operator $`\stackrel{~}{c}_{i,\sigma }`$ of equation (27) which will be used to transfer spectral weight from the inner parts of the Brillouin zone to the outer ones in case of normal photoemission $`\omega <\mu `$. As already discussed above, this also improves the resolution of the ME method in this region, since its resolution strongly depends on the spectral weight at a certain position. Nevertheless the spectrum of the shadow operator has to exhibit exactly the same peak positions as the normal photoemission spectrum.
Figure 3 shows the angle-resolved spectral function of the shadow operator for moderate and low temperatures. As expected, the shadow operator has its main spectral weight near $`𝐤=(\pi ,\pi )`$ on the photoemission side and near $`𝐤=(0,0)`$ on the inverse photoemission side. Furthermore, it’s spectrum supports the existence of a total of four bands, because it resolves a group of peaks forming dispersionless bands at energies of $`\omega \pm 6.0t`$, a region where the normal photoemission spectrum exhibits only some weak and smeared-out spectral weight. These two dispersionless bands at energies of $`\omega \pm 6.0t`$ are inconsistent with the dispersions of the Hubbard-I and SDW approximations of figure 1. We will further address this topic later in this work.
Next, we turn in more detail to the temperature dependence of the photoemission spectrum. Figure 4 shows some close-ups of the normal photoemission spectrum and of the spectrum of the shadow operator at momentum $`𝐤=(\pi ,\pi )`$ and different temperatures. For the normal photoemission operator these close-ups show a peak at $`\omega 1.5t`$, which would be consistent with Hubbard-I (see Figure 2). In the spectrum of the shadow operator this feature is visible as a single resolved peak only at the highest temperature, $`T=1.00t`$ whereas for temperatures down to $`T=0.20t`$ there is only some diffuse weight at this position. In the ordinary photoemisson spectrum the peak looses spectral weight with decreasing temperature. It disappears completely at $`T<0.20t`$ where the spin-correlation length $`\zeta _{sz}(T)`$ reaches the system size (see above). Thus, the temperature $`T0.20t`$ where we lose the Hubbard-I-like peak at $`\omega 1.5t`$ and $`𝐤=(\pi ,\pi )`$ coincides quite accurately with the temperature where ‘effective’ long-range order sets in. Moreover, we find that as the normal photoemission spectrum loses the peak at $`\omega 1.5t`$ and $`𝐤=(\pi ,\pi )`$, the spectrum of the shadow operator gains weight at $`\omega 6.0t`$. Thus, we expect that both features are closely related to the temperature development of the spin-correlation length $`\zeta _{sz}(T)`$. We note, however, that the crossover in the shape of the dispersion from Hubbard-I-like to SDW-like occurs in a quite unexpected way: the topmost band at $`(\pi ,\pi )`$ does not deform into the SDW form in any continuous way, but simply ‘fades away’ and eventually vanishes at the transition.
A further surprising result is the following: at $`T=0.10t`$ neither the ordinary electron operator nor the shadow operator pick up the ‘AF umklapp band’ corresponding to the narrow dispersive band seen for example at $`\omega 3t`$ for $`𝒌=(0,0)`$, i.e. there is no corresponding band at $`\omega 3t`$ and $`𝒌=(\pi ,\pi )`$. Note that in the framework if the SDW-approximation the shadow operator must reproduce this umklapp-band at $`(\pi ,\pi )`$ with the same weight as the original band at $`(0,0)`$ in the ordinary photoemission spectrum - see the discussion in the first section. That this is not the case shows that even at this lowest temperature, a simple SDW-like description of the band structure is invalid, in that the band structure cannot be understood by simple backfolding of the spectrum obtained without broken symmetry. As we will see in the following the AF SDW state provides only the ‘background’ for the dressing of the photoholes with AF spin-excitations which dynamically generate a total of four bands.
## III The 4-band structure
We return to the discrepancy between the 2-band dispersions of the Hubbard-I and SDW approximations and the 4-band structure actually observed for example in the spectrum of the shadow operator of figure 3.
In order to generate a ‘4-band structure’ out of the two bands of the Hubbard-I and SDW approximations we try as a phenomenological ansatz to mix the dispersions of the Hubbard-I/SDW approximation with two dispersionless bands at energies of $`E_\pm =\pm 3.0t`$. In other words, for both the photoemission and inverse photoemission spectrum we diagonalize an ‘effective’ $`2\times 2`$ Hamilton matrix:
$$H_\pm =\left(\begin{array}{cc}E_\pm ^{HubI/SDW}& V\\ V& E_\pm \end{array}\right),$$
(34)
and plot in figure 5 the four bands obtained in this way on top of the spectral density obtained from QMC at $`T=0.33t`$ and $`T=0.10t`$. For comparison the figure also shows the original (i.e. unhybridized) Hubbard-I bands plus the two phenomenological dispersionless bands at energies of $`\omega =\pm 3.0t`$. The figure shows that the overall agreement between the QMC peak positions and the four bands generated by diagonalizing $`H_\pm `$ of equation (34) is surprisingly good, particularly so in view of the fact that for both, Hubbard-I and SDW approximation, only unrenormalized parameters were used. In particular, the self-consistently determined value for the SDW gap $`\mathrm{\Delta }`$ of $`3.56t`$ was used at $`T=0.10t`$. The only ‘external’ parameter in this figure is the mixing matrix element $`V`$, which was set to a value of $`1.0t`$.
Thus, we find in contrast to previous works , that the introduction of the dispersionless bands reproduces the sinlge-particle gap and the width of the quasiparticle band correctly without any renormalization of parameters. Rather, the narrowing of the quasiparticle band and the reduction of the Hubbard gap as compared to the unrenormalized parameters is brought about by introduction of the dispersionless bands. This naturally raises the question as to their physical origin. In the present paper we restrict ourselves to a more phenomenological and ‘numerics based’ approach. A complementary and more mathematical discussion is given in Ref. , where an equation of motion approach similar to Hubbard’s original work is pursued.
We consider the commutator of the creation operator for hole-like particles, $`h_{i,\sigma }^{}=c_{i,\sigma }(1n_{i,\overline{\sigma }})`$, which annihilates a particle only on a singly occupied site, with the kinetic energy of the Hubbard model and find:
$`[h_{i,}^{},H_t]`$ $`=`$ $`t{\displaystyle \underset{jN(i)}{}}[(1{\displaystyle \frac{n}{2}})c_{j,}+(c_{j,}S_i^z+c_{j,}S_i^{})`$ (36)
$`{\displaystyle \frac{1}{2}}c_{j,}(n_in)+c_{j,}^{}c_{i,}c_{i,}.`$
Keeping only the first term in the square brackets of equation (36) reproduces the Hubbard-I approximation . The second term in the square brackets describes the dressing of the created hole by a spin-excitation and is closely related to the spin-$`1/2`$ string operator of equation (31). The third term describes in an analogous fashion the coupling of the hole to a density fluctuations, whereas the fourth term describes the coupling to the $`\eta `$-excitation. The two latter types of excitation are not important for a large positive $`U`$ near half-filling, $`n=1.0`$, and will be neglected. Therefore, the operator $`_{jN(i)}(c_{j,}S_i^z+c_{j,}S_i^{})`$ is in this case the most important correction over the Hubbard-I approximation. As already stated, it describes a hole dressed by a spin-excitation: this operator not only creates a hole on site $`j`$ but dresses this hole with a spin-excitation on a neighboring site, which is exactly the idea behind the spin-bag or spin-polaron pictures known in the literature.
Splitting this operator into eigenoperators of $`H_U`$:
$`\widehat{C}_i`$ $`=`$ $`{\displaystyle \underset{jN(i)}{}}(\widehat{c}_{j,}S_i^z+\widehat{c}_{j,}S_i^{}),`$ (37)
$`\widehat{D}_i`$ $`=`$ $`{\displaystyle \underset{jN(i)}{}}(\widehat{d}_{j,}S_i^z+\widehat{d}_{j,}S_i^{}),`$ (38)
we find $`[\widehat{D}_{i\sigma },H_U]=\frac{U}{2}\widehat{D}_{i,j,\sigma }`$ and $`[\widehat{C}_{i\sigma },H_U]=\frac{U}{2}\widehat{C}_{i,j,\sigma }`$. Assuming moreover that the mobility of these composite excitations is determined by the ‘heavy’ spin-excitation, it seems quite resonable to assume that these ‘particles’ are the source of the (more or less) dispersionless bands at $`\pm U/2`$.
Finally, the commutation relation (36) shows that the mixing matrix element between the $`h_{𝒌,\sigma }`$ and the new composite particles should be $`t`$. Based on these rough considerations we might thus expect that the two string-$`1`$ ‘effective particles’ (38) are excellent candidates for explaining the two dispersionless bands at $`\pm 3t`$ required to upgrade the Hubbard-I or SDW approximation so as to match the QMC data. However, so far the above considerations are pure speculation and in the following we will turn to QMC-results to back up this hypothesis by numerical evidence.
Before doing so, however, we want to illustrate the action of the ‘string-operator’ in two extreme cases: an ideal Néel state and a resonating valence-bond (RVB) state, i.e. a compact singlet covering of the plane (see figure 6). In the Néel state, a hole created initially on site $`i`$ can travel one place to a neighboring site $`j`$ thereby leaving behind a misaligned spin on the original site $`i`$. Exactly this process is described by the second term, $`c_{j,}S_i^{}`$: it creates a hole of opposite spin on a neighboring site $`j`$ and flips the spin on the original site $`i`$. Therefore, this process corresponds to the creation of a string of length $`1`$. In fact one might think about more sophisticated diagnostic operators incorporating the effects of longer-ranged strings . Indeed, Dagotto and Schrieffer and Eder and Ohta already measured the angle-resolved spectrum of a diagnostic operator containing strings with up to three lattice sites range by means of exact diagonalizations of the $`tJ`$ model . As already mentioned in the introduction, in the QMC method each observable has to be expressed in terms of free single-particle Greens functions on each time slice by the application of Wicks theorem . This results already in a quite large expression for the spin-$`1/2`$ string operator of equation (31) containing approximately $`80`$ contributions. The implementation of even longer-ranged string operators therefore was not possible.
Returning to the spin-$`1/2`$ string operator of equation (31) we note that the first term, $`c_{i,}S_j^z`$ will always annihilate a Néel state. This reflects the simple fact that spin-rotation symmetry is broken in the Néel state. This is not the case, however, in the fully rotationally invariant RVB state: again, creating a hole on site $`i`$ and allowing it to hop to site $`j`$ will produce a spin-excitation. However, in the case of the RVB state, it produces the superposition of two states: in one case the dotted ellipse stands for the $`S_z=1`$ component of the triplet, and this state would again be created by the term $`c_{j,}S_i^{}`$. There is, however, also a second state where the dotted ellipse corresponds to the superposition of a singlet and the $`S_z=0`$ component of the triplet. This second state then would be created by the term $`c_{i,}S_j^z`$, whereby the relative sign of the two terms in the string-1 operator makes sure that the two configurations are always produced with the proper phase. In both extreme cases, Néel state and ‘singlet soup’, the string-$`1`$ operator thus creates a hole dressed with the proper spin-excitation: this scan be a spin wave (i.e. a single inverted spin) in the case of a Néel state, or a singlet-triplet excitation in the case of an RVB state.
As a technical remark we still note that the excessive numerical effort which would have been necessary to compute spectra for the $`\widehat{C}_{i,\sigma }`$ and $`\widehat{D}_{i,\sigma }`$ (which are products of $`5`$ Fermion operators) has made it impossible to compute the spectra of these operators - instead we have been using the (Fourier transform of) the operator $`b_{i,}`$, defined in (31). Concerning the difference between this diagnostic operator and the operator $`_{jN(i)}(c_{j,}S_i^z+c_{j,}S_i^{})`$ obtained by commuting $`h_{i,}^{}`$ with the kinetic energy (see the second term on the r.h.s. of equation (36)), we note that their Fourier transforms differ only by phase factors of the form $`e^{i𝐤(𝐑_i𝐑_j)}`$. Both operators are indeed identical at momentum $`𝐤=(0,0)`$ and differ at momentum $`𝐤=(\pi ,\pi )`$ only by a factor of $`1`$ since $`𝐑_i`$ and $`𝐑_j`$ are next-nearest neighbor lattice sites. At all other momenta both operators have exactly the same peaks but will differ somewhat in their spectral weights.
After these remarks we discuss the angle-resolved spectral function of the spin-$`1/2`$ string operator $`b_{𝐤,}`$, shown in figure 7. The spectrum of the spin-$`1/2`$ string operator indeed picks up those band portions which, according to the rough hybridization scenario in figure 5, should have strong ‘flat-band character’, i.e. the part of the narrow low energy band between $`(0,0)`$ and $`(\pi /2,\pi /2)`$ and between $`(0,0)`$ and $`(\pi ,0)`$ at energies $`\omega 3t`$. We note that these are precisely the positions where the two quasiparticle-like dispersive narrow bands occurred in the normal photoemission spectrum at $`T=0.33t`$. In addition, the band portion at $`\omega 6t`$ for momentum $`(\pi ,\pi )`$ is also enhanced in the string-1 spectrum.
This, however, still leaves an important part of the band structure unexplained. Namely, the ‘AF umklapp band’ of the narrow quasiparticle band dispersing upward between $`(0,0)`$ and $`(\pi /2,\pi /2)`$ at $`\omega 3t`$ still is not seen in any of the spectra, not even at the lowest temperature studied. On the other hand, in a state with true antiferromagnetically broken symmetry we know that this mirror image must exist due to the backfolding of the Brillouin zone. To finally resolve this part, we now introduce the last diagnostic operator of this work, which we call the spin-$`3/2`$ string operator:
$$\stackrel{~}{b}_{i,}=\underset{jN(i)}{}(\mathrm{𝟐}c_{i,}S_j^zc_{i,}S_j^{}).$$
(39)
This describes again a composite object of a hole and a spin-excitation, but this time the two constituents are coupled to the total spin of $`3/2`$. We stress that this operator will detect states which can never be seen in an actual angle-resolved photoemission spectrum on a singlet ground state, because this is forbidden by the angular-momentum selection rule. The angle-resolved spectral function $`\stackrel{~}{A}(𝐤,\omega )`$ of the spin-$`3/2`$ string operator is plotted in figure 8 again for $`T=0.33t`$ (top) and $`T=0.10t`$ (bottom). It is then immediately obvious that it is this operator which resolves the ‘missing piece’ of the AF dispersion, i.e. the ‘AF mirror image’ of the narrow quasiparticle band. It should also be noted that the spectrum of $`\stackrel{~}{b}_{i,}`$ is remarkably independent of temperature, i.e. the states belonging to this ‘spin-$`3/2`$ band’ persist irrespectively of whether there is long-range order or not.
Combining the information obtained so far, suggests the following scenario for the crossover between the paramagnetic band structure at high temperature and the AF band structure at low temperature: in the paramagnetic state at high temperatures (such as $`T=0.33t`$), the spin is a good quantum number and the spin-$`3/2`$ ‘band’ does exist but cannot mix with any spin-$`1/2`$ band due to spin-conservation. The band of spin-$`3/2`$ quasiparticles thus plays no role whatsoever in the actual photoemission spectrum, which is presumably the reason why the outer part of the spectrum is so remarkably ‘invisible’ in actual ARPES experiments, leading to the idea of a ‘remnant Fermi surface’ in the insulator. Reaching this state by photoemission would only be possible if the photohole is created in a thermally admixed state of at least spin $`1`$. In an infinite system this situation changes discontinuously at the transition to the true broken symmetry state: there the total spin ceases to be a good quantum number, and the spin-$`1/2`$ band in the interior of the AF zone and the spin-$`3/2`$ band in the exterior now suddenly can mix with each other, thus leading to the familiar SDW dispersion. We note that another way to generate a coupling between these two bands would be application of a magnetic field - this also would break spin-rotation invariance and hence enable the hybridization of a spin-$`1/2`$ and a spin-$`3/2`$ band. Based on our results we thus believe that a magnetic field could enhance the spectral weight of the ‘shadow part’ of the band structure as seen in ARPES.
## IV Doping the Hubbard model at high temperatures – rigid bands
Summarizing our results so far, we may say that the Hubbard-I approximation, slightly improved by the introduction of new quasiparticles corresponding to dressed holes provides a very good description of the spectral function for the case of half-filling, $`n=1`$. In this section we want to proceed to the doped case $`n<1`$, which is of prime interest for cuprate superconductors. Here, an essential drawback of the QMC procedure is that reliable QMC simulations for lower temperatures are much more difficult or even impossible, since the absence of particle-hole symmetry away from half-filling introduces the notorious minus-sign problem into the algorithm. Truely low temperatures like $`T=0.10t`$, which in principle correspond to the physical temperature range, are therefore out of reach. On the other hand in the study of the half-filled case we have seen that a major change takes place as the spin correlation length reaches the system size, whence ‘effective long range order’ sets in. In the doped case the spin correlation length is expected to be short at any temperature, whence we may expect that the change of $`A(𝐤,\omega )`$ from high to small temperature is more smooth than at half-filling. In that sense, even $`A(𝐤,\omega )`$ data for the relatively high temperature $`T=0.33t`$ are interesting to study. Moreover, we can at least try to elucidate trends with decreasing temperature and thus construct a reasonably plausible scenario.
At half-filling, we have seen that the ‘approximation of choice’ for the paramagnetic case was the Hubbard-I approximation. This naturally poses the question as to how relevant the half-filled case is for the description of the doped case, i.e. how much of the Hubbard-I physics remains valid for finite doping. At half-filling the two ‘effective particles’ $`\widehat{d}_{i,\sigma }^{}`$ and $`\widehat{c}_{i,\sigma }^{}`$, form the two separate Hubbard bands. The effect of doping would now consist in the chemical potential cutting progressively into the top of the lower Hubbard band, in much the same fashion as in a doped band insulator. On the other hand, for finite $`U/t`$ the spectral weight along this band deviates from the free-particle value of $`1`$ per momentum and spin so that the Fermi surface volume (obtained from the requirement that the integrated spectral weight up to the Fermi energy be equal to the total number of electrons) is not in any ‘simple’ relationship to the number of electrons - the Luttinger theorem must be violated. This is the major reason why the Hubbard-I approximation has been dismissed by many authors as being unphysical.
We now wish to address the question as to what really happens if a paramagnetic (i.e. not magnetically ordered) insulator is doped away from half-filling, by QMC simulation. We therefore choose $`T=0.33t`$ and $`U=8.0t`$.
Figure 9 then shows the development of $`A(𝒌,\omega )`$ with doping. It is quite obvious from this Figure that initially the $`2`$ bands seen at half-filling in the photoemission spectrum (i.e. $`\omega <0`$) persist with an essentially unchanged dispersion. The chemical potential gradually cuts deeper and deeper into the topmost band, forming a hole-like Fermi surface centered on $`(\pi ,\pi )`$, the top of the lower Hubbard band. The only deviation from a rather simple rigid-band behavior is an additional transfer of spectral weight: the part of the topmost band near $`(\pi ,\pi )`$ gains in spectral weight, whereas the band with higher binding energy looses weight. In addition, there is a transfer of weight from the upper Hubbard band to the inverse photoemission part below the Hubbard gap. This effect is actually quite well understood. The band structure above the Hubbard gap becomes more diffuse upon hole doping in that the rather clear two-band structure visible near $`(\pi ,\pi )`$ at half-filling rapidly gives way to one broad ‘hump’ of weight. Apart from the spectral weight transfer, however, the band structure on the photoemission side is almost unaffected by the hole doping - the dispersion of the quasiparticle band becomes somewhat wider but does not change appreciably. In that sense we see at least qualitatively the behavior predicted by the Hubbard I approximation.
Next, we focus on the Fermi surface volume. Some care is necessary here: first, we cannot actually be sure that at the high temperature we are using there is still a well-defined Fermi surface. Second, the criterion we will be using is the crossing of the quasiparticle band through the chemical potential. It has to be kept in mind that this may be quite misleading, because band portions with tiny spectral weight are ignored in this approach (see for example Ref. for a discussion). When thinking of a Fermi surface as the constant energy contour of the chemical potential, we have to keep in mind that portions with low spectral weight may be overlooked. On the other hand the fact that a peak with appreciable weight crosses from photoemission to inverse photoemission at a certain momentum is independent of whether we call this a ‘Fermi surface’ in the usual sense, and should be reproduced by any theory which claims to describe the system. It therefore has to be kept in mind that in the following we are basically studying a ‘spectral weight Fermi surface’, i.e. the locus in $`𝐤`$ space where an apparent quasiparticle band with high spectral weight crosses the chemical potential. With these caveats in mind, Figures 10 and 11 show the low-energy peak structure of $`A(𝒌,\omega )`$ for all allowed momenta of the $`8\times 8`$ cluster in the irreducible wedge of the Brillouin zone, and for different hole concentrations. In all of these spectra there is a pronounced peak, whose position shows a smooth dispersion with momentum. Around $`(\pi ,\pi )`$ the peak is above $`\mu `$, whereas in the center of the Brillouin zone it is below. The locus in $`𝒌`$-space where the peak crosses $`\mu `$ forms a closed curve around $`(\pi ,\pi )`$ and it is obvious from the Figure that the ‘hole pocket’ around $`(\pi ,\pi )`$ increases very rapidly with $`\delta `$. To estimate the Fermi surface volume $`V_F`$ we assign a weight $`w_𝐤`$ of $`1`$ to momenta $`𝐤`$ where the peak is below $`\mu `$, $`0.5`$ if the peak is right at $`\mu `$ and $`0`$ if the peak is above $`\mu `$. Our assignments of these weights are given in Figures 10 and 11. The fractional Fermi surface volume then is $`V_F=\frac{1}{N}_𝐤w_𝐤`$, where $`N=64`$ is the number of momenta in the $`8\times 8`$ cluster. Of course, the assignment of the $`w_𝐤`$ involves a certain degree of arbitrariness. It can be seen from Figures 10 and 11, however, that our $`w_𝐤`$ would in any way tend to underestimate the Fermi surface volume, so that the obtained $`V_F`$ data points rather have the character of a lower bound to the true $`V_F`$. Even if we take into account some small variations of $`V_F`$ due to different assignments of the weight factors, however, the resulting $`V_F`$ versus $`\delta `$ curve never can be made consistent with the Luttinger volume, see Figure 12. The deviation from the Luttinger volume is quite pronounced at low doping. $`V_F`$ approaches the Luttinger volume for dopings $`20`$%, but due to our somewhat crude way of determining $`V_F`$ we cannot really decide when precisely the Luttinger theorem is obeyed. The Hubbard I approximation approaches the Luttinger volume for hole concentrations of $`50`$%, i.e. the steepness of the drop of $`V_F`$ is not reproduced quantitatively. The latter is somewhat improved in the so-called $`2`$-pole approximation. For example the Fermi surface given by Beenen and Edwards for $`n=0.94`$ obviously is very consistent with the spectrum in Figure 11 for $`n=0.95`$.
We return to Figure 9 and discuss the entire width of the spectra, in particular the question of the fate of the 4-band structure in the doped system. For $`n=0.95`$ the different features that are seen at $`n=1.0`$ are still rather clearly visible, but for $`n=0.86`$ the low energy quasiparticle band at $`k=(0,0)`$ starts to disappear, and at $`n=0.80`$ the dominant ‘band’ in the spectrum between $`\omega =4t`$ and $`\omega =2t`$ can be fitted by a sligtly renormalized free-electron band. As we have seen above, the Luttinger theorem also is valid in this case. This suggests to classify the doping as ‘underdoped’ for $`0<n<0.85`$, where the Luttinger theorem is invalid and the $`4`$-band structure known from half-filling persists, and ‘overdoped’ where the Luttinger theorem is valid and a renormalized free-electron band can be seen in the spectral function. Following the convention for cuprate superconductors, we call the doping where the crossover between the two regimes occurs the ‘optimal’ doping.
Next, the four plots of figure 13 show the spectra at selected $`𝒌`$-points. The system size is only $`6\times 6`$ in this case because this allows for smaller error bars. Closer inspection, especially, of the peaks at momentum $`𝐤=(0,0)`$ on the photoemission side (plot (a)) and at momentum $`𝐤=(\pi ,\pi )`$ on the inverse photoemission side (plot (d)) confirms, that with increasing hole concentration we are losing parts of the 4-band structure seen at half-filling.
To check the physics of the band structur in more detail, we again employ our diagnostic operators. Figure 14 shows the angle-resolved spectral functions $`\stackrel{~}{A}(𝐤,\omega )`$ of the spin-$`1/2`$ and spin-$`3/2`$ string operators, $`b_{i,}`$ and $`\stackrel{~}{b}_{i,}`$. As was the case at half-filling, the spectrum of the spin-$`1/2`$ string operator highlights exactly those peaks that we associate with the dipsersionless ‘dressed hole’ bands in Figure 5. The spectrum of the spin-$`3/2`$ string operator on the other hand has its peaks with maximal spectral weight around momentum $`𝐤=(\pi ,\pi )`$, indicating that also in the doped case there is an ‘antiferromagnetic mirror image’ of the quasiparticle band (which, however, consists of spin-$`3/2`$ states). Again, coupling of photoholes to thermally excited spin excitations may make these states visible in ARPES spectra, thus explaining the ‘shadow bands’ seen in photoemission experiments by Aebi et al.. Similarly as for half-filling one might speculate that a magnetic field, which would break spin symmetry and thus allow for a coupling of ‘bands’ with different total spin, would enhance the spectral weight of these shadow bands.
All in all we have ssen that the ‘band structure’ ($`4`$-band structure, dispersion of regions of large spectral weight, ‘character’ of the bands as measured by the diagnostic operators) stays pretty much unchanged as long as we are in the underdoped regime. At half-filling the $`4`$-band structure is closely related to the sharp low-energy mode in the dynamical spin correlation function, which naturally suggests to study the spin reponse also as a function of doping. Figure 15 shows the spin-correlation function, $`\chi _{sz}(𝐤,\omega )`$ (left column), and the charge-correlation function, $`\chi _{cc}(𝐤,\omega )`$ (right column), for $`T=0.33t`$ and densities $`n=0.95`$ (underdoped), $`n=0.90`$ (nearly optimally doped) and $`n=0.80`$ (overdoped). The spin-response is sharply confined in both momentum $`𝐤=(\pi ,\pi )`$ and energy $`\omega =\omega ^{}`$ only in the underdoped region i.e. the regime where we also observe the features associated with spin excitations in the single-particle spectra. As was the case at half-filling for temperatures below $`T0.33t`$, the spin-response can be fitted by the AF spin-wave dispersion (33) in the underdoped regime. On the other hand, as soon as the system enters the overdoped regime the spin-response is no longer sharply peaked at momentum $`𝐤=(\pi ,\pi )`$ and energy $`\omega =\omega ^{}`$: it broadens in momentum and spreads in energy by an order of magnitude with the scale changing from $`J=4t^2/U`$ to $`E_{kin}8.0t`$ accompanied by a similar change in the bandwidth of the single particle excitations. This result is already well known from previous QMC calculations and consistent with similar behaviour in the t-J model. The charge-response, $`\chi _{cc}(𝐤,\omega )`$, is always broad in both momentum $`𝐤`$ and energy $`\omega `$ for all densities studied. It merely decreases its width from $`12.0t`$ at $`n=0.95`$ to $`8.0t`$ at $`n=0.80`$.
Although the minus-sign problem of the QMC algorithm prevents reliable simulations of large systems at low temperatures in the doped regime, we nevertheless studied the temperature evolution of the angle-resolved spectral function $`A(𝐤,\omega )`$ at density $`n=0.93`$. This was possible due to the relative small system size of $`6\times 6`$, which alleviates the minus-sign problem as compared to $`8\times 8`$ at the $`T=0.25t`$. Figure 16 shows the results from this analysis: the uppermost plot (a) compares the angle-resolved spectral functions $`A(𝐤,\omega )`$ at density $`n=0.93`$ for $`T=0.50t`$, $`T=0.33t`$ and $`T=0.25t`$. We stress that the simulation at $`T=0.25t`$ suffers from minus-sign problems with a drastically reduced resolution. In the center plot (b), the quasiparticle peak weights around momentum $`𝐤=(\pi ,\pi )`$ of the $`n=0.93`$ simulation are compared with the quasiparticle peak weight at momentum $`𝐤=(\pi ,\pi )`$ of a half-filled, $`n=1.0`$, simulation for different temperatures. At half-filling, the Hubbard-I-like quasiparticle peak $`𝐤=(\pi ,\pi )`$ and $`\omega 1.5t`$ decreases in spectral weight with decreasing $`T`$ and disappears as the spin-correlation length (which increases with decreasing temperature) reaches the lattice size (at $`T0.20t`$). In the underdoped case for density $`n=0.93`$ the weights of the corresponding peaks around momentum $`𝐤=(\pi ,\pi )`$ located also decrease with decreasing $`T`$. Closer inspection of this peak (see the inset of the center plot (b)) reveals, that this peak even raises slightly in binding energy with decreasing temperature, very similar to the peak in the half-filled case. In a real photoemission experiment this peak would have dropped below the typical resolution of roughly $`10\%`$ spectral weight at a temperature of $`T0.25t`$. The spin-correlation length (again derived by a fit of the equal-imaginary-times spin-correlation function to a form $`a|𝐫|^be^{|𝐫|/\zeta _{sz}(T)}`$) also shows similar behavior in the underdoped and half-filled cases: the values for the spin-correlation length are $`\zeta _{sz}=0.50.8`$, $`\zeta _{sz}=0.81.0`$ and $`\zeta _{sz}=1.01.3`$ in the underdoped case and $`\zeta _{sz}=0.60.9`$, $`\zeta _{sz}=1.01.3`$ and $`\zeta _{sz}=1.61.9`$ at half-filling for $`T=0.50t`$, $`T=0.33t`$ and $`T=0.25t`$, respectively. The spin-susceptibility (shown in plot (c) of figure 16) also behaves very similar, but with changed magnitudes in the underdoped and in the half-filled cases. These data suggest a similar temperature evolution of the band-structure of the Hubbard model in the underdoped and half-filled cases driven by the temperature dependent spin-correlation length $`\zeta _{sz}(T)`$. Especially, we expect the Hubbard-I-like quasiparticle peaks at energies of $`\omega 1.0t`$ around momentum $`𝐤=(\pi ,\pi )`$ to vanish with decreasing $`T`$ in the underdoped case as the peaks at energies of $`\omega 1.5t`$ around momentum $`𝐤=(\pi ,\pi )`$ do in the case of half-filling.
The latter observation suggests a profound change of the Fermi surface with temperature: as seen above, it is precisely the Hubbard-I-like band near $`(\pi ,\pi )`$ which crosses the chemical potential and thus forms the Fermi surface in the doped case. It is then quite clear that the ‘disappearance’ of this band with decreasing temperature must affect the Fermi surface in some dramatic way. Studies at zero temperature are possible only by means of exact diagonalization. Analysis of the single particle spectrum shows the same ‘rigid-band’ behaviour as at high temperatures and analysis of the momentum distribution $`n(𝒌)`$ suggest that the doped holes accumulate at the surface of the magnetic zone (i.e. the line $`(\frac{\pi }{2},\frac{\pi }{2})(\pi ,0)`$) rather than around $`(\pi ,\pi )`$.
## V Summary
In the present work we have systemactically studied the temperature- and doping-dependent dynamics of the two-dimensional Hubbard model by finite-temperature QMC simulations. Comparing the QMC single particle spectral function, the dynamical spin response and the spectral functions of suitably chosen diagnostic operators, different physical regimes could be identified. In simplest terms there are two quantities, which basically determine the single particle spectrum: the hole concentration and the the spin-response function, whereby there is a certain relationship between the two.
At half-filling and high temeratures ($`Tt`$), the combined photoemission and and inverse photoemission spectrum $`A(𝐤,\omega )`$ displays two dispersive features, the upper and lower Hubbard band, roughly separated by $`U`$ ($`=8t`$, in our work). At these very high temperatures the system is in a spin disordered state. We have demonstrated here that the well-known Hubbard-I approximation gives an excellent description of the single particle spectrum in this state, reproducing quantitatively both the single-particle dispersion and the distribution of spectral weight. This is by no means trivial, since the Hubbard-I approximation is dynamically equivalent to a simplified effective Hamiltonian, which just contains hole-like ($`h_{i,\sigma }^{}`$) and double occupancy-like particles in a simple bi-quadratic form.
At lower temperatures ($`T0.33t`$), the Hubbard-I approximation needs to be improved; this is to be expected, because it neglects all effects of spin-correlations. In fact, the temperature where deviations from Hubbard-I become strong, coincides fairly well with the transition from a spin response $`\chi _{sz}(𝐤,\omega )`$ which is diffuse both in momentum and energy (with a spread of order $`t`$), to a more ‘spin-wave-like’ response. In this regime $`\chi _{sz}(𝐤,\omega )`$ displays the characteristic energy scale $`J=4t^2/U`$, with its spectral weight being concentrated at the AF wave vector $`𝐐=(\pi ,\pi )`$. It should be noted that this ‘spin-wave-like’ regime develops despite the fact that at $`T=0.33t`$ the spin-correlations length $`\zeta _{sz}(T)`$ is still short-ranged ($`2`$ lattice spacings). Only at the lower temperature $`T=0.10t`$ Néel order spreads over the entire QMC block, creating an effective (finite-size) Néel state.
It is well established by previous, in particular also QMC work, that in this temperature-regime ($`T0.33t`$) new spectral features appear. They have often been interpreted as four ‘bands’, two ‘coherent’ bands forming the topmost valence and the lowest conduction band in the insulator plus two ‘incoherent’ bands, i.e. the remaining upper and lower Hubbard band features (see for example ). Our present work not only definitively identifies these four bands but also clarifies their physical origin and their connection to the spin-excitations. In simplest terms the emerging spin waves at lower temperatures provide the excitations that can ‘dress’ the Hubbard quasiparticles, whence new bands corresponding to dressed holes/double occupancies appear in the single particle spectrum $`A(𝐤,\omega )`$. It has been shown in Ref. that the $`4`$-band structure which appears in $`A(𝐤,\omega )`$ at lower temperatures can be explained in this way, and our present numerical check by directly calculating the spectra of ‘dressed electrons’ supports this interpretation.
This physical picture at half-filling can be extended into the underdoped regime. This is most obvious in the single particle spectral function, which stays almost unchanged in the doped case (i.e. the $`4`$-band structure and the ‘character’ of the bands as measured by the diagnostic operators). The main change in fact consists in the chemical potential cutting gradually into the (top of the) lower Hubbard band, precisely as predicted by the Hubbard-I approximation. Contrary to widespread belief the ‘Fermi surface’, if determined by the Fermi surface crossings of the dominant band through the chemical potential, does not satisfy the Luttinger theorem. Rather, for small hole concentrations the Fermi surface volume is considerable larger than that for a slightly less than half-filled free-electron band. Very similar conclusions have in fact been reached by a calculation of the electron momentum distribution in the 2D t-J model by Puttika et al.. Their calculation actually was a high temperature series expansion plus a Padé extrapolation to lower temperature, and it is encouraging that this method gives similar results as our QMC results which are performed at relatively high temperatures. In its range of applicability, i.e. in the absence of strong magnetic correlations and close to half-filling, the Hubbard-I approximation thus works remarkably well, both at half-filling and in the doped case. We stress that this has profound implications for the theoretical treatment of the model: perturbation expansions in $`U`$ or partial and self-consistent resummations thereof, may not be expected to give any meaningful results in this strong-coupling/low doping regime.
An interesting question is the possibility to verify our results experimentally. As already mentioned above, a scan of the temperature development of $`A(𝒌,\omega )`$ shows that the part of the quasiparticle band near $`(\pi ,\pi )`$ (where the Fermi surface is located) is loosing weight with decreasing temperature. In fact in ARPES experiments on underdoped cuprate superconductors the ‘hole-pocket’ around $`(\pi ,\pi )`$ seen in our simulations (and the expansion of Puttika et al.) is not observed, but rather a small ‘Fermi arc’ near $`(\pi /2,\pi /2)`$, terminated by the ‘pseudo gap’ around $`(\pi ,0)`$ Although our simulations do not allow to make statements about the truely low temperatures in the experiments, we believe that this suggests a strong temperature dependence of the single particle spectrum, with the temperature scale being set by the exchange constant $`J`$ (which controls the degree of spin disorder). The latter is rather large in copper oxides, so that the temperature regime studied in our simulations probably cannot be accessed experimentally in these materials. We note, however, that an ARPES study for the 1D material $`Na_{0.96}V_2O_5`$ which has a smaller exchange constant, has indeed provided evidence for a strong $`T`$-dependence of $`A(𝒌,\omega )`$. Clearly, it would be interesting to study the Fermi surface evolution in a 2D material with lower exchange constant.
As was the case at half-filling, the dynamical spin-response plays an important role: throughout the Hubbard-I phase at low doping, the spin response shows the sharp low-energy mode at $`(\pi ,\pi )`$. The simultaneous disappearance of the $`4`$-band structure in $`A(𝒌,\omega )`$ and the low energy spin response with scale $`J`$ in the overdoped regime then show again the close relationship between the two. The dressing of holes by spin excitations apparently remains the most important correction over Hubbard-I. In the overdoped regime the spin response is spread out over an energy range of $`8t`$ and thus becomes more similar to the charge response. The single-particle spectral function is most consistent with a slightly renormalized free electron dispersion, and the Luttinger theorem appears to be satisfied even at the relatively high temperature $`T=0.33t`$. This is quite consistent with earlier results on the t-J model, which show that for hole concentrations $`25`$ % the spin and charge response can be approximated well by the self-convolution of the single particle spectral function. This is essentially what is to be expected for a system of weakly interacting Fermions, so that we conclude that in the overdoped regime we enter a new phase which most probably extends to the low-concentration limit where the Nagaoka $`T`$-matrix approximation becomes exact. Finally we note that exact diagonalization studies at finite temperatures also show some evidence for a ‘crossover’ between different physical regimes at a hole concentration of $`15`$%.
This work was supported by DFN Contract No. TK 598-VA/D03 and by BMBF (05SB8WWA1). Computations were performed at HLRS Stuttgart and HLRZ Jülich. One of us (W.H.) acknowledges hospitality of the Physics Department in Santa Barbara and many useful discussions with D. J. Scalapino. |
warning/0001/cond-mat0001220.html | ar5iv | text | # Reducing quasi-ergodicity in a double well potential by Tsallis Monte Carlo simulation
## 1 Introduction
The ergodic hypothesis is fundamental to statistical mechanics. This hypothesis states that the time average of an observable event equals the phase-space average. In practical application, however, problems can arise in various types of simulations if the system must overcome high energy barriers to reach other regions of phase space. In that case, the length of a simulation needed in order to obtain enough statistical samples of all regions of phase space may be extremely long. In the Monte Carlo simulation, the errors arise as a consequence of the finite length of the Monte Carlo walk. This error can be serious in canonical Monte Carlo sampling, especially at low temperatures. This problem, referred as ”quasi-ergodicity” by Valleau and Whittington , appears even in the simplest double well problem, where the two wells are separated by a large barrier .
Recently several authors have proposed a generalized simulated annealing method to locate the global minimum based on the generalized statistical mechanics proposed by Tsallis , which can overcome the slow convergence of the traditional simulated annealing method based on the standard Monte Carlo scheme. Subsequently, Andricioaei and Straub have pointed out that their generalized Monte Carlo (Tsallis Monte Carlo) method may be used to overcome the quasi-ergodicity appearing in the canonical average at constant temperatures.
In this short note we use a new Monte Carlo scheme based on Tsallis’s generalized statistical mechanics for the calculation of the thermal average of the potential energy of a simple double well potential at constant temperatures. We examine the performance of this new algorithm designed to overcome the quasi-ergodicity.
## 2 A generalized Monte Carlo scheme
In the generalized statistical mechanics proposed by Tsallis , a crucial role is played by the generalized entropy $`S_q`$ defined as
$$S_q=k\frac{1p_i^q}{q1}$$
(1)
where $`q`$ is a real number which characterizes the statistical mechanics, and $`p_i`$ is the probability of states $`i`$. This entropy $`S_q`$ becomes the usual Gibbs-Shannon entropy $`S_1=kp_i\mathrm{ln}p_i`$ when $`q1`$. By maximizing the generalized entropy with the constraints
$$p_i=1\text{and}p_i^qϵ_i=\text{const}$$
(2)
where $`ϵ_i`$ is the energy spectrum, the statistical weight for the generalized canonical ensemble characterized by the parameter $`q`$ is given by the Tsallis weight
$$p_i^q=\mathrm{exp}\left(\overline{\beta }\overline{ϵ}_i\right)/Z_q^q$$
(3)
with
$$\overline{ϵ}_i=\frac{q}{\overline{\beta }(q1)}\mathrm{ln}\left[1(1q)\overline{\beta }ϵ_i\right]$$
(4)
where $`\overline{ϵ}_i`$ is the generalized energy and $`\overline{\beta }`$ is the Lagrange multiplier, which plays the role of the generalized temperature in the Tsallis statistical mechanics. We should note here that the statistical weight is not given by the probability distribution $`p_i`$ but rather by $`p_i^q`$ in order to calculate thermal averages, as can be seen in (2). The generalized canonical partition function $`Z_q`$ in (3) is given by
$$Z_q=\mathrm{exp}\left(\overline{\beta }\overline{ϵ}_i/q\right)$$
(5)
In the limit $`q1`$, $`\overline{ϵ}_iϵ_i`$ and $`\overline{\beta }\beta =1/kT`$ where $`T`$ is the usual temperature , then the Tsallis weight (3) becomes the Boltamann weight $`p_i=\mathrm{exp}(\beta ϵ_i)/Z_1`$ of the usual canonical ensemble.
In this generalized statistical mechanics, Andricioaei and Straub have noted that since the average of an observable $`O`$ is defined by
$$<O>_q=p_i^qO_i,$$
(6)
the detailed balance condition should be written as
$$p_i^qW_{ij}=p_j^qW_{ji}$$
(7)
where $`W_{ij}`$ is the element of the transition matrix. Based on this observation, they pointed out that the generalized Monte Carlo algorithm originally proposed by Penna in his generalized simulated annealing method cannot satisfy the detailed balance. They then proposed a new generalized Monte Carlo scheme, where the acceptance probability $`p`$ of the Monte Carlo move is given by
$$p=\mathrm{min}[1,\mathrm{exp}(\overline{\beta }\mathrm{\Delta }\overline{ϵ})]$$
(8)
which has a similar form to the familiar Metropolis algorithm
$$p=\mathrm{min}[1,\mathrm{exp}(\beta \mathrm{\Delta }ϵ)]$$
(9)
where $`\mathrm{\Delta }\overline{ϵ}`$ is the increase of the generalized energy (4) while $`\mathrm{\Delta }ϵ`$ is that of the usual energy. Andricioaei and Straub combined this Monte Carlo algorithm with the usual simulated annealing technique to optimize the conformation of tetrapeptides. They showed that their algorithm is more effective than the standard simulated annealing based on the usual molecular dynamics or Monte Carlo methods.
Recently, Andricioaei and Straub pointed out that their generalized Monte Carlo scheme for simulated annealing can also be used to calculate the standard canonical ensemble averages at constant temperatures
$$<O>_1=\frac{O_i\mathrm{exp}(\beta ϵ_i)}{\mathrm{exp}(\beta ϵ_i)}$$
(10)
using the general rule of the importance sampling
$$<f>=f_i=\left(\frac{f_i}{g_i}\right)g_i$$
(11)
Then, equation (10) can be calculated using the generalized Monte Carlo scheme from
$`<O>_1`$ $`=`$ $`{\displaystyle \frac{O_i\mathrm{exp}\left(\left(\beta ϵ_i\overline{\beta }\overline{ϵ}_i\right)\right)\mathrm{exp}\left(\overline{\beta }\overline{ϵ}_i\right)}{\mathrm{exp}\left(\left(\beta ϵ_i\overline{\beta }\overline{ϵ}_i\right)\right)\mathrm{exp}(\overline{\beta }\overline{ϵ}_i)}}`$ (12)
$`=`$ $`{\displaystyle \frac{<O\mathrm{exp}\left(\left(\beta ϵ\overline{\beta }\overline{ϵ}\right)\right)>_q}{<\mathrm{exp}\left(\left(\beta ϵ\overline{\beta }\overline{ϵ}\right)\right)>_q}}`$ (13)
where the generalized average $`<>_q`$ is defined by (6). In practice, the equilibrium thermal average in the standard canonical ensemble can be calculated by conducting these ”generalized” Monte Carlo steps and accumulating the samples by
$$<O>_1=\frac{O_j\mathrm{exp}\left(\left(\beta ϵ_j\overline{\beta }\overline{ϵ}_j\right)\right)}{\mathrm{exp}\left(\left(\beta ϵ_j\overline{\beta }\overline{ϵ}_j\right)\right)}$$
(14)
where the sum $`j`$ runs over to the generalized Tsallis Monte Carlo steps produced by eq.(8) instead of the usual Metropolis algorithm (9). The equation (14) is a generalized expression, and we may choose $`\overline{\beta }=\beta `$. Then, the equation (14) reduces to
$$<O>_1=\frac{O_j\mathrm{exp}\left(\beta \left(ϵ_j\overline{ϵ}_j\right)\right)}{\mathrm{exp}\left(\beta \left(ϵ_j\overline{ϵ}_j\right)\right)},$$
(15)
which has been used in . From now on, the actual calculation will be done by choosing $`\overline{\beta }=\beta `$.
We would like to point out here that a similar reweighting technique is used in the histogram method and the multicanonical ensemble method . Andricioaei and Straub applied this Tsallis Monte Carlo algorithm to the two-dimensional Ising system , a classical one-dimensional potential and a 13-atom Lennard-Jones cluster , and showed that this new Monte Carlo algorithm is more effective than the standard Metropolis algorithm. They also proposed the generalized molecular dynamics based on Tsallis generalized statistical mechanics
## 3 Application to a classical double well potential
In order to check the performance of this new Monte Carlo algorithm in a classical system such as molecules and liquids, we look at the problem of a classical particle in a one-dimensional double well potential defined by
$$V(x)=\delta \left(3x^4+4(\alpha 1)x^36\alpha x^2\right)+1$$
(16)
where
$$\delta =\frac{1}{2\alpha +1}$$
(17)
which was examined by Frantz et al. to demonstrate the quasi-ergodicity of the standard Metropolis algorithm. The potential given by (16) has the double minimum at $`x=1`$ and $`x=\alpha `$; therefore, it represents a symmetrical double well when $`\alpha =1`$ and a single well when $`\alpha =0`$. This potential is more conveniently characterized by the parameter $`\gamma `$
$$\gamma =\frac{V(0)V(\alpha )}{V(0)V(1)}=\alpha ^3\left(\frac{\alpha +2}{2\alpha +1}\right)$$
(18)
In figure 1, we show the asymmetric double well potential $`V(x)`$ when $`\gamma =0.9`$ as well as the equilibrium Tsallis weight $`p_\mathrm{T}\mathrm{exp}(\beta \overline{V}(x))`$, which is renormalized, and Boltzmann weight $`p_\mathrm{B}=\mathrm{exp}(\beta V(x))/Z_1`$. It is obvious that the statistical weight $`p_\mathrm{T}`$ of the Tsallis generalized statistical mechanics has a greater chance of crossing the barrier and, therefore, a greater possibility of avoiding quasi-ergodicity.
Alternatively this Tsallis’s statistical mechanics can be regarded as a search-space smoothing method which deforms the rugged potential landscape by a smoother one. In figure 2, we show the smoothed potential $`\overline{V}(x)`$ calculated from (4), which has a smoother landscape and a lower energy barrier than the original landscape $`V(x)`$.
We use this double well potential to look at the quasi-ergodicity of the standard as well as the generalized Monte Carlo scheme by examining the classical average potential energy
$$<V>=\frac{V(x)\mathrm{exp}\left(\beta V(x)\right)𝑑x}{\mathrm{exp}\left(\beta V(x)\right)𝑑x}$$
(19)
which can be calculated rigorously by numerical quadrature as a function of the temperature $`\beta `$ and the asymmetry parameter $`\gamma `$. We employ the same Monte Carlo procedure of Frantz et al. and use the uniform sampling distribution from previous position $`x`$ to new position $`x^{^{}}`$ given by
$`T(x^{^{}}|x)`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Delta }}}\text{for}|xx^{^{}}|<{\displaystyle \frac{\mathrm{\Delta }}{2}}`$ (20)
$`=`$ $`0\text{otherwise}`$
The step size $`\mathrm{\Delta }`$ can be chosen to be large enough for the one-dimensional problem so that the quasi-ergodicity can be reduced . However, we follow Frantz et al. and use the scaling $`\mathrm{\Delta }=2.5/\sqrt{\beta }`$, which maintains an approximately 50% acceptance rate at all temperatures $`\beta `$, by considering the application to multidimensional systems. This 50% acceptance is commonly employed in the classical Monte Carlo community .
Since the error caused by the quasi-ergodicity depends upon the walk initialization, we start the Monte Carlo walk at the global minimum $`x=1`$ and at the metastable minimum $`x=\alpha `$. If the walk starts at the global minimum, then the average potential $`<V>`$ will be low if the walk is quasi-ergodic, as the distribution associated with the higher energy well around $`\alpha `$-minimum will be insufficiently sampled. While the walk is initialized at $`\alpha `$, the walk may be trapped at this metastable $`\alpha `$-well and the average potential $`<V>`$ will be too high.
In figure 3, we show the average potential energy $`<V>`$ calculated from (i) the standard Metropolis Monte Carlo and (ii$``$iv) the generalized Tsallis Monte Carlo algorithm as a function of temperature $`\beta `$. Similar curves can be found in figure 2 of reference ; however, those authors plotted the results from the usual Metropolis algorithm initialized at the global minimum $`x=1`$ and at the random positions. In order to check the quasi-ergodicity of the walk, we instead show two curves initialized at the global minimum $`x=1`$ and the metastable $`\alpha `$-minimum at $`x=\alpha `$ for each scheme. The average potential $`<V>`$ is obtained from 100 independently initialized walks, each consisting of 500 warm-up steps followed by $`10^4`$ steps with data accumulation. The average energy $`<V>`$ calculated from the standard Metropolis algorithm shows large errors due to the quasi-ergodicity. In particular, the walk initialized at the metastable $`\alpha `$-minimum greatly overestimates the potential energy. The walk starting from the global minimum also shows noticeable error and underestimates the average energy, as expected.
In contrast to the results from the standard Metropolis algorithm, the results from the generalized Monte Carlo algorithm show excellent agreement with the numerically exact result. The walks started from the global minimum $`x=1`$ produce the average energy which is almost indistinguishable from the exact result up to the lowest temperature (high $`\beta `$) examined. Even the walk started from the metastable $`\alpha `$ minimum produces the energies that are fairly accurate up to rather low temperature $`\beta 15`$ for $`q=1.5`$, which shows that this generalized algorithm does a much better job than the standard Monte Carlo scheme. The agreement with the exact result is further improved if we increase the value of parameter $`q`$ from $`1.5`$ (see Figure 3). It should be noted that the probability distribution associated with the parameter $`q=2`$ is given by the Cauchy distribution and with $`q>5/3`$ by the super-diffusive Lévy distribution . The theoretical upper boundary is $`q=3`$. The quasi-ergodicity can be largely circumvented by the use of Tsallis’s generalized statistical mechanics with $`q>1`$, as shown in figure 3.
In order to overcome the quasi-ergodicity, the Monte Carlo walk has to pass many times through the energy barrier. Figure 4 compares the time series of the Monte Carlo walk generated by the standard Monte Carlo scheme ($`q=1`$) and that of the generalized Monte Carlo scheme. Similar diagrams can be found in reference . Although the trajectory of the standard Monte Carlo scheme is mostly confined within the metastable $`\alpha `$ minimum, the number of barrier crossing of the Tsallis Monte Carlo scheme is significantly greater. This barrier crossing occurs more often when the parameter $`q`$ is larger. The quasi-ergodicity can be removed by Tsallis’s generalized Monte Carlo scheme by increasing the magnitude of the parameter $`q`$ from $`1`$.
In figure 5, we show the average energy $`<V>`$ as a function of the asymmetry parameter $`\gamma `$. Again a similar diagram can be found in reference . Here again we can clearly observe a large error for the initialization at the metastable $`\alpha `$ minimum when the standard Monte Carlo procedure is used. In this case, the superiority of Tsallis’s generalized Monte Carlo scheme over the traditional one is also obvious.
In table 1 we showed the effect of quasi-ergodicity in the convergence when $`\gamma =0.9`$ and $`\beta =10`$ for both the standard Monte Carlo and the generalized Monte Carlo algorithms with $`q=1.5`$, each originating at the metastable $`\alpha `$-minimum. We show the average energy $`<V>`$ and the standard deviation (STD) obtained from 100 independently initialized walks. We found that the problem of quasi-ergodicity obtained with the standard Metropolis algorithm is still present, even if we increase the number of Monte Carol steps. On the other hand the result from the generalized Monte Carlo simulation is significantly improved by increasing the number of Monte Carlo steps, and ergodicity can be recovered. Furthermore, STD decreases as $`1/\sqrt{N}`$ when the Monte Carlo step $`N`$ is increased for the generalized Monte Carlo simulation, while such a systemic decrease is not observed for the standard Monte Carlo simulation.
This convergence can be clearly visualized by plotting the visiting probability $`P(x)`$ of the Monte Carlo walks, which should be given by the Boltzmann weight for the standard Monte Carlo scheme and by the Tsallis weight (3) for the generalized Monte Carlo scheme. In figure 6, we compare $`P(x)`$ generated by the standard and generalized Monte Carlo walks with the Boltzmann and the Tsallis weights. The visiting probability $`P(x)`$ converges rapidly to the Tsallis weight as the number of the Mone Carlo walks is increased for the generalized Monte Carlo scheme. On the other hand, $`P(x)`$ recovers the Boltzmann weight only when the Monte Carlo steps are increased upto $`10^6`$ in the standard Monte Carlo scheme.
It is also interesting to discuss the convergence in the context of the ”generalized ergodic measure” of Thirumalai, Mountain, and Kirkpatrick . Using the energy metric, the ergodic measure $`d_V(n)`$ is defined by
$$d_V(n)=\left(V^a(n)V^b(n)\right)^2,$$
(21)
where $`V^a(n)`$ and $`V^b(n)`$ are the move average of the potential energy $`V`$ along the trajectories $`a`$ and $`b`$, which are defined from (15) by
$`V^a(n)`$ $`=`$ $`{\displaystyle \frac{_{j=0}^nV_j^a\mathrm{exp}\left(\beta \left(V_j^a\overline{V}_j^a\right)\right)}{_{j=0}^n\mathrm{exp}\left(\beta \left(V_j^a\overline{V}_j^a\right)\right)}}.`$ (22)
for the generalized Tsallis Monte Carlo scheme, and
$`V^a(n)`$ $`=`$ $`{\displaystyle \frac{_{j=0}^nV_j^a\mathrm{exp}\left(\beta V_j^a\right)}{_{j=0}^n1}}.`$ (23)
for the standard Monte Carlo scheme. For an ergodic system, Thirumalai et al suggested that the ergodic measure converges as $`d_V(n)\mathrm{}`$ if $`n\mathrm{}`$. They found that
$$d_V(n)d_V(0)\frac{1}{D_Vn}$$
(24)
where the ”diffusion” constant $`D_V`$ depends on temperature. In figure 7, we show $`d_V(0)/d_V(n)`$ for $`\gamma =0.9`$ potential obtained from 100 independent pairs ($`a`$, $`b`$) of Monte Carlo walks as a function of Monte Carlo steps $`n`$. We chose the trajectory starting from the metastable minimum $`x=\alpha `$ (=-0.9163) and that from the stable minimum $`x=1`$, respectively, as the two independent trajectories $`a`$ and $`b`$. It is clear from figure 7 that the normalized inverse of the ergodic measure $`d_V(0)/d_V(n)`$ grows linearly with the Monte Carlo steps $`n`$. Therefore, the ergodic measure $`d_V(n)`$ decreases according to (24), and the diffusion constant $`D_V`$ can be determined.
Figure 8 shows the diffusion constant $`D_V`$ as a function of temperature $`\beta `$. The diffusion constant is a decreasing function of the inverse temperature $`\beta `$; it is more difficult to recover the ergodicity when the temperature is lowered. However, at lower temperatures $`\beta >20`$ with $`\beta \mathrm{\Delta }V>1`$ where $`\mathrm{\Delta }V=V(\alpha )V(1)=0.1`$, the diffusion constant starts to increase again for the Tsallis Monte Carlo walks ($`q1`$). This is due to the fact that the equilibrium thermal distribution around the metastable minimum at $`x=\alpha `$ becomes negligibly small at such low temperatures, and the ergodicity of the trajectory among the stable and the metastable minima is irrelevant once trajectories $`a`$ and $`b`$ both fall into the stable minimum around $`x=1`$.
It is apparent from figure 8 that the diffusion constant $`D_V`$ obtained from the standard Monte Carlo algorithm obeys the usual activation form
$$D_V\mathrm{exp}\left(\beta E_b\right)$$
(25)
where the activation energy $`E_b`$ is estimated to be $`E_b0.4`$, which is the same order of magnitude as the barrier height ($`1`$) between two wells, as shown in figure 1. On the other hand, the diffusion constant $`D_V`$ obtained from the Tsallis generalized Monte Carlo scheme is characterized by a power-law form
$$D_V\beta ^\eta $$
(26)
with the exponent $`\eta `$. This result is conceivable, based on the acceptance probability (8) using the Tsallis weight, because
$`D_V`$ $``$ $`\mathrm{exp}(\beta \overline{E}_b)`$ (27)
$``$ $`\left[1(1q)\beta E_b\right]^{q/(q1)}`$
and $`D_V\beta ^{q/(q1)}`$. We note in figure 8 that $`\eta 1.6`$ for $`q=1.5`$ and $`\eta 0.85`$ for $`q=2.5`$. These values are close to $`q/(q1)=3`$ for $`q=1.5`$ and $`q/(q1)=1.66`$ if $`q=2.5`$. Therefore, the diffusion constant decreases more slowly as a function of the inverse temperature $`\beta `$ as $`\beta ^\eta `$ for the generalized Monte Carlo scheme, though it decreases exponentially faster for the standard Monte Carlo scheme. The exponent $`\eta `$ becomes smaller as $`q`$ is increased and the diffusion is possible even at low temperatures. This result can be anticipated based on data presented in figures 3 and 6. This power-law temperature dependence of the diffusion constant $`D_V`$ for $`q1`$ has been pointed out by Straub and Andricioaei .
## 4 Concluding remarks
In conclusion, we found that the so-called ”quasi-ergodicity” in a double well potential encountered in the standard Metropolis Monte Carlo algorithm is largely circumvented by the use of the generalized Monte Carlo algorithm of Andricioaei and Straub based on Tsallis’s generalized statistical mechanics. Therefore this algorithm will be useful in the Monte Carlo study of any system, in particular, when the first order phase transition occurs. We believe that this algorithm will be useful, for example, in the study of the melting transition of clusters for which the J-walk algorithm and the multiple histrogram method have been utilized. It is also interesting to apply this new algorithm to the problems where the slow dynamics due to randomness or frustration is serious, for example, the spin-glass problems. This generalized Monte Carlo algorithm based on Tsallis’s statistical mechanics use the reweighting technique similar to the multicanonical method . But this algorithm is much simpler than the multicanonical methods and appears to be more promising because it does not require an iterative determination of the weighting factor by preliminary Monte Carlo runs.
Recently, a new formulation of the Tsallis statistics has been proposed . With a new choice of energy constraint, unfamiliar consequences in the previous formulation have been erased. However, the Lagrange parameter used in the new formulation can be related to that in the previous formulation. Therefore, in performing Monte Carlo simulations with reweighting technique, we may use the previous formulation of the Tsallis statistics.
Acknowlegement
Authors are grateful to the reviewer for several useful comments which were incorporated into the final version of the manuscript. M. I. acknowledges support from the Hiroshima City University Grant for Special Academic Research. He is also grateful to the Department of Physics, Tokyo Metropolitan University, where this work was initiated, for its hospitality to him. |
warning/0001/astro-ph0001333.html | ar5iv | text | # Numerical Simulations of Globular Cluster Formation
## 1. Introduction
The globular clusters belong to the oldest populations in our galaxy. For the general reviews on globular clusters, see Meylan, & Heggie (1997) and Harris (1991). Their formation (Elmergreen et al. 1999) is closely related to the formation process of our galaxy. The formation of a globular cluster may take place in two stages: (1) the formation of a proto-globular cloud (PGC) and (2) the formation of a star cluster from the PGC. There are three scenarios for the formation of a PGC, denoted as primary, secondary, and tertiary model (Fall & Rees 1987), where the PGC forms in different stages, i.e., before, during, or after the collapse of the galaxy, respectively.
* The primary model was suggested by Peebles & Dicke (1968). They showed that the Jeans mass of the recombination stage of the universe is comparable to the observed masses of globular clusters and thus the PGC can form due to the gravitational instability. Globular clusters may be debris of these objects. The most serious problem with this primary model is that there have been few intergalactic globular clusters discovered. Almost all globular clusters ever discovered exist in galaxies.
* In the secondary model, on which we concentrate in the present paper, the PGC forms due to thermal instabilities (Fall & Rees 1985). Our detailed investigation and calculation are presented in the following sections.
* For the tertiary model, one example is that globular clusters form from large-scale unorganized motion of interstellar gas as in the Magellanic Clouds, which are now producing young clusters (Kumai, Basu, & Fujimoto 1993). Another example is that very young globular clusters are observed in nearby galaxy NGC 1705 and 1569 (Ho & Filippenko 1996).
Once the PGC forms, there might be many ways to form a globular cluster (GC). Murray & Lin (1993) summarized the scenarios of the formation of GCs from PGCs as follows. The PGC can be divided into two types depending on their masses. The cloud, whose mass exceeds the Jeans mass, is gravitationally unstable, thus spontaneously collapsing to form stars. The cloud, whose mass is smaller than the Jeans mass, is stable until some instabilities are introduced. A cloud-cloud collision or cloud-disk collision can trigger such instability. When such collisions occur, the cloud would become thermally unstable owing to the efficient cooling. This cooling would lead to the formation of a very dense region, thus inducing a burst of star formation.
In the present paper, we examine quantitatively the above still qualitative scenario of the globular cluster formation. We use our Smoothed Particle Hydrodynamics (SPH) code (Lucy 1977; Gingold & Monaghan 1977) to simulate the formation of a globular cluster from a PGC. Our code includes the following physical processes: radiative cooling, star formation, energy feedback from stars including stellar winds and supernovae, and chemical enrichment from stars. The SPH method which includes star formation processes like ours has been applied to many astrophysical problems. Such problems include the formation of isolated galaxies (Katz 1992; Steinmetz & Müller 1994), the evolution of galaxies (Friedli & Benz 1995), and the cosmological simulations (Navarro & White 1993). However, our study is the first attempt to apply this method to the globular cluster formation (preliminary results have been presented in Nakasato et al. (1996); Nakasato, Mori, & Nomoto (1999)).
Among the two triggers of the instability, collapse and collision, we concentrate on the collapse case in the present paper. The collision case for a wide range of parameters (masses of the two clouds, a relative velocity between the clouds etc.) will be discussed elsewhere. This paper consists of the following sections. In section 2 our SPH method is described. In section 3, we present our linear analysis for the thermal instabilities to form a PGC. In section 4, we present the results of the evolution of the PGC with two different initial compositions of gases. Sections 5 and 6 are devoted to discussions and concluding remarks, respectively.
## 2. Method
To simulate the formation of a star cluster from gases, we use our GRAPE-SPH code using Remote-GRAPE library (Nakasato, Mori, & Nomoto 1997). The GRAPE is a special purpose computer for the self gravity calculations in general N-body system (Sugimoto et al. 1990). Our GRAPE-SPH code includes various physical processes, e.g., radiative cooling, star formation, and feedback from formed stars. There are many different implementations of the SPH method (a recent comparison of various SPH implementations is presented in Thacker et al. (1998)). In the present paper, we describe an essential point of our SPH code. The SPH formulation that we use is the same as the described in Navarro & White (1993). We use the smoothing length that can vary spatially and evolve with time, and integrate equations of motion with a second order Runge-Kutta method as described in Navarro & White (1993). Details of the implementation of our code are described in Mori et al. (2000). In the following subsections, we describe the physical processes (cooling, star formation, and feedback) in some detail.
### 2.1. Radiative cooling
The radiative cooling rates depend on the temperature and ionization states of the gas. Also, the chemical composition of the gas affects the cooling rate drastically. We perform SPH simulations for the following three cases (A, B, C) of the gas with different physical state and chemical composition.
Case A: We assume that the chemical composition of the gas is primordial with no heavy elements and the gas is in ionization equilibrium. In this case, our treatment of radiative cooling is essentially the same as adopted by Katz, Weinberg, & Hernquist (1995); they computed the cooling rate using the two-body processes of H and He, and free-free emissions. The cooling function ($`\mathrm{\Lambda }(T)`$) is shown in Figure 1 with the solid line. In this case, the cooling rate decreases very rapidly as the temperature $`T`$ decreases below $`2\times 10^4`$ K so that the gas would not radiatively cool below $`T10^4`$ K.
Case B: We assume that the gas includes some heavy elements and is in ionization equilibrium. In this case, we use the cooling function with different chemical composition that is computed by MAPPINGS III software by R.S. Sutherland (MAPPINGS III is the successor of MAPPINGS II that is described in Sutherland & Dopita (1993)). We compute the cooling function of the ionization equilibrium gas for \[Fe/H\] $`=5.0`$ \- 0.0 with the solar abundance ratio (see Table 4 of Sutherland & Dopita (1993)) and present the results in Figure 1 with the dashed lines (each line corresponds to \[Fe/H\] = $`1,2,3`$ respectively from the top to the bottom). Existence of heavy elements significantly enhances the cooling rates. At $`T<10^4`$ K, the cooling due to the forbidden line emission of carbon and oxygen is efficient. For \[Fe/H\] = $`1`$, the cooling rates around $`T10^5`$K are 100 times larger than the cooling rates of primordial hydrogen and helium gas. These differences would make the evolution of the gas very different.
Case C: We concern about the non-equilibrium cooling. If the gas cools from high temperatures, ionization equilibrium is not realized. Figure 2 shows the ratio between the recombination time ($`t_{\mathrm{recom}}`$) of hydrogen and the cooling time ($`t_{\mathrm{cool}}`$) of primordial gas in ionization equilibrium; $`t_{\mathrm{recom}}`$/$`t_{\mathrm{cool}}`$. The $`t_{\mathrm{recom}}`$ and $`t_{\mathrm{cool}}`$ are defined as
$$t_{\mathrm{recom}}=\frac{1.0}{\alpha (T)f_e(T)}$$
(1)
and
$$t_{\mathrm{cool}}=\frac{3\mathrm{k}_\mathrm{b}T}{2\mathrm{\Lambda }(T)},$$
(2)
where the $`\alpha `$ is the recombination coefficient of hydrogen and electron, $`f_e`$ is the fraction of free electrons, and $`\mathrm{k}_\mathrm{b}`$ is the Bolzman constant. For $`\alpha `$, we use the value of Spitzer (1978). Clearly, the cooling time is much shorter than the recombination time for $`T2\times 10^4`$ K, so that ionization equilibrium is not realized when the gas cools from high temperatures. In the non-equilibrium case, the existence of electrons and ionized hydrogen at $`T<10^4`$ K makes it possible to form H<sub>2</sub> molecules through the creation of intermediaries H<sup>-</sup> and H$`{}_{}{}^{+}{}_{2}{}^{}`$ as:
$`\text{H}+\text{e}\text{H}^{}+\mathrm{photon}`$
$`\text{H}+\text{H}^{}\text{H}_2+\text{e}`$ (3)
and
$`\text{H}+\text{H}^+\text{H}_2^++\mathrm{photon}`$
$`\text{H}_2^++\text{H}\text{H}_2+\text{H}^+`$ (4)
(Shapiro & Kang 1987). These H<sub>2</sub> molecules cause further cooling down to $`T10^2`$ K owing to the rotational-vibrational line excitation.
To include the effect of such molecular cooling, we have to solve rate equations that determine the ionization states of H and He atoms, and the formation and destruction of H<sub>2</sub> molecules. In our case, the typical time step in solving the rate equations is shorter than the dynamical time step, which is determined mainly by the Courant condition. Integrating all equations with a shorter timestep than the dynamical time is very costed work. So we divide the dynamical timestep with dynamical variables (density etc.) being constant in solving the energy and rate equations, which is similar to the method adopted in Shapiro & Kang (1987). The rate-coefficients are also the same as used in Shapiro & Kang (1987). Included species are H<sup>0</sup>, H<sup>+</sup>, He<sup>0</sup>, He<sup>+</sup>, He$`{}_{}{}^{+}+`$, H<sup>-</sup>, H$`{}_{2}{}^{}{}_{}{}^{+}`$, H<sub>2</sub>, H$`{}_{}{}^{}{}_{2}{}^{}`$ and e, where H$`{}_{}{}^{}{}_{2}{}^{}`$ is the exited hydrogen molecule. In our SPH code, we solve the rate equations for 10 species in each SPH particle with a reasonable computing time. Solving the rate equations for the gas including heavy elements (over 200 species) is not feasible with a current resource so that we only concern the hydrogen and helium plasma in the present paper.
The cooling function for the non-equilibrium case is presented in Figure 1 with the dotted line. In computing these cooling rates, we follow the isobaric temperature evolution of the fully ionized gas. Initially, the temperature and hydrogen number density of the gas are 10<sup>7</sup> K and 0.01 cm<sup>-3</sup>, respectively. We assume optically thin plasmas so that the gas cools rapidly. Around $`T2\times 10^4`$ K, H<sub>2</sub> molecules begin to form. At $`T<10^4`$ K, the cooling rate due to H<sub>2</sub> molecules is comparable to the cooling rate for \[Fe/H\] = $`1`$ gas.
Finally, we summarize the treatment of the cooling rates in our SPH code. We can perform SPH simulations for these three cases A, B and C. In the present paper, we will study the evolution of a PGC with two different initial chemical compositions, e.g., a metal-free gas cloud and a metal-rich gas cloud. In the former case, we solve the energy and rate equations simultaneously for each SPH particle (case B) if the gas temperature is lower than $`3\times 10^4`$ K and the gas particle is not in the heating phase (see section 2.3). If these conditions are met, we use the pre-computed cooling table for the ionization equilibrium case (case A). In the latter case, we use the pre-computed cooling table to solve the energy equations (case C).
### 2.2. Star formation
Our treatment of star formation is the same as adopted by Katz (1992). Hereafter, “STAR” means “star particle”, which is not an individual star but an association of many stars. A STAR forms in the region where the flow is converging, cooling, and Jeans unstable. These conditions are expressed as
$$(𝒗)_i<0,$$
(5)
$$t_{\mathrm{cool}}<t_\mathrm{d},$$
(6)
$$t_\mathrm{d}<t_{\mathrm{sound}}.$$
(7)
Here $`t_{\mathrm{cool}}`$, $`t_\mathrm{d}`$ , and $`t_{\mathrm{sound}}`$ are the cooling time, dynamical time, and sound crossing time, respectively, and expressed as
$$t_{\mathrm{cool}}=\frac{\mu ^2u}{\rho \mathrm{\Lambda }},$$
(8)
$$t_\mathrm{d}=\frac{1}{\sqrt{4\pi G\rho }},$$
(9)
$$t_{\mathrm{sound}}=\frac{h_i}{c_s},$$
(10)
where $`\mu `$ is the mean molecular weight, $`u`$ is the specific thermal energy, and $`c_s`$ is the local sound speed. A STAR forms in the region where these all three conditions are satisfied. The star formation rate is given as
$$\frac{D\rho _{}}{Dt}=\frac{\rho }{t_{\mathrm{starform}}}=C\sqrt{4\pi G}\rho ^{\frac{3}{2}},$$
(11)
where $`\rho _{}`$ is the density of a star, and $`t_{\mathrm{starform}}=t_\mathrm{d}/C`$ with a star formation parameter $`C`$. In the present paper, we adopt $`C=1.0`$, i.e., a STAR forms in a dynamical time. The last term of equation 11 is similar to the Schmidt’s law (Schmidt 1959). Integrating equation (11) over one time step $`\delta t`$ and making calculations with the equations for SPH, we obtain the mass of a newly formed star in $`\delta t`$ as
$$m_{\mathrm{star}}=\left[1\frac{1}{1+0.5\frac{\delta t}{t_{\mathrm{starform}}}}\right]\pi h_i^3\rho _i.$$
(12)
The newly formed STAR is then treated as a collision-less particle. We introduce the minimum allowed mass ($``$ 10 $`M_{}`$) for the newly formed STAR to prevent the unphysical effects in the equation of motion and the treatment of the feedback processes.
We note the star formation recipes in our SPH code. In the regions with increasing densities, first two conditions for the star formation (Eq. 5 and 6) are almost satisfied. Thus whether STAR forms in some regions is determined mostly by the Jeans criterion (Eq. 7). The critical density for the star formation in our treatment is obtained as
$$\rho _{\mathrm{Jeans}}>\left(\frac{\gamma b}{4\pi G\mu m}\right)^3\frac{1}{m_i^2}T^3\left(\frac{1}{\beta ^6}\right).$$
(13)
Here, it is assumed that the neighbor radius is expressed as
$$h_i=\beta \left(\frac{m_i}{\rho _i}\right)^{1/3},$$
(14)
where $`\beta `$ is determined experimentally, because $`h_i`$ is determined in order to make the number of neighbor particles almost constant in some range (30 - 80 in our codes). Since the typical value of $`\beta `$ is 1.0 - 1.1, $`\beta ^6`$ ranges 1.0 - 2.0. Then if $`m_i`$ is constant, the critical density is determined almost only by the temperature. The typical calculations in the present paper use 5000 gas particles for a 10<sup>6</sup> $`M_{}`$ gas sphere. Thus the initial mass per particle ($`m_i`$ in Eq. 13) is $``$ 200 $`M_{}`$. For $`T=10^4`$ K and 10<sup>2</sup> K, the critical density for star formation is $`\rho _{\mathrm{Jeans}}10^{17}`$ and 10<sup>-23</sup> g cm<sup>-3</sup>, respectively. With the star formation recipes used in our SPH code, the STAR forms only in the very high density region if $`T10^4`$ K. There are, however, the maximum density ($`\rho _{\mathrm{max}}`$) that numerically reached in the SPH method; it is estimated
$$\rho _{\mathrm{max}}\frac{N_\mathrm{n}m_i}{ϵ^3},$$
(15)
where $`N_\mathrm{n}`$ is the number of neighbor particles and $`ϵ`$ is the gravitational softening length. In the present calculations, $`\rho _{\mathrm{max}}1.8\times 10^{20}`$ g cm<sup>-3</sup>, These arguments ensure that with the star formation recipes used in our SPH code and the initial conditions of the present calculations, the STAR forms in the region where the temperature is as low as 10<sup>2</sup> K.
### 2.3. Feedback
The formed stars eject gases and heavy elements in stellar winds and Type II supernova explosions and heats up, accelerate, and enrich a circumstellar and an interstellar medium. High energy explosions like supernova produce high temperature and low density regions in interstellar medium. In the SPH method, the numerical accuracy for high density regions is much better than mesh based methods but the accuracy for low density regions are poorer. A typical resolution of usual SPH simulations, e.g., Navarro & White (1993), including a star formation process (100 - 1000 pc) is larger than a typical size of supernova remnants ($`<100`$ pc). Thus, it is difficult to properly include the energy, momentum, and mass release from stars in the SPH method because of the nature of the SPH method and the poor resolutions in current computing resources. We must use some approximations to mimic real feedback processes in the SPH method.
One of the method has been proposed by Navarro & White (1993). In their method, the energy produced by a supernova explosion is distributed to neighbor gas particles of each STAR mostly as a thermal energy and the rest is distributed as a velocity perturbation to the gas particles; the fraction of the energy in a kinetic form is a free parameter (we note that Leitherer, Robert, & Drissen (1992) presented the population synthesis model of stellar feedback processes). In the present paper, we distribute the energy to neighbor particles in a pure thermal form as a zero-th order approximation, because the size and time scale of our model are much smaller than those of a galaxy formation model of Navarro & White (1993).
#### 2.3.1 Energy ejection
The energy ejection rate per STAR is given as
$$E_{\mathrm{eject}}=e_{\mathrm{SW}}R_{\mathrm{SW}}+e_{\mathrm{SNII}}R_{\mathrm{SNII}},$$
(16)
where $`e_{\mathrm{SW}}`$ is the total energy ejected by stellar winds during the stellar lifetime and $`e_{\mathrm{SNII}}`$ are the energy ejected by one Type II supernova explosion. The $`R_{\mathrm{SW}}`$ is the number of stars per unit time expelling their envelopes at the current epoch and $`R_{\mathrm{SNII}}`$ is the rate of Type II supernovae. We define the $`R_{\mathrm{SW}}`$ and $`R_{\mathrm{SNII}}`$ as follows
$$R_{\mathrm{SW}}=\frac{{\displaystyle _{M_{\mathrm{ms}}}^{M_{\mathrm{up}}}}\mathrm{\Phi }(m)𝑑m}{\tau (M_{\mathrm{ms}})}$$
(17)
$$R_{\mathrm{SNII}}=\frac{{\displaystyle _{M_{\mathrm{ms}}}^{M_{\mathrm{ma}}}}\mathrm{\Phi }(m)𝑑m}{\tau (M_{\mathrm{ms}})\tau (M_{\mathrm{ma}})},$$
(18)
where $`\mathrm{\Phi }(m)`$ is the initial mass function (IMF), namely $`\mathrm{\Phi }(m)dm`$ gives the number of stars in the mass range of ($`m`$, $`m+dm`$) and the $`\tau (m)`$ is the stellar lifetime (David, Forman, & Jones 1990). In the present study, we assume the power law type IMF as
$$\mathrm{\Phi }(m)=Am^{2.35},$$
(19)
where the $`A`$ is the constant. For the upper and lower limit masses in Eq. (19), $`M_{\mathrm{up}}=`$ 120 $`M_{}`$ and $`M_{\mathrm{lo}}=`$ 0.05 $`M_{}`$ are assumed. In Eq. (17) and (18), $`M_{\mathrm{ma}}`$ ($`=`$ 50.0 $`M_{}`$) and $`M_{\mathrm{ms}}`$ ($`=`$ 8.0 $`M_{}`$) are the upper and lower limit masses of the stars that explode as Type II supernovae.
For the supernova energy, we assume that $`e_{\mathrm{SNII}}=10^{51}`$ erg. For the stellar wind, $`e_{\mathrm{SW}}`$ is estimated to be 0.2$`\times 10^{51}`$ erg for solar metallicity stars from the observational data of OB associations (Abbot 1982). The Chemical abundance of a massive star significantly affects $`e_{\mathrm{SW}}`$ (Leitherer, Robert, & Drissen 1992), so that we use metallicity dependent e<sub>SW</sub> as $`e_{\mathrm{SW}}=0.2e_{\mathrm{SNII}}(Z/Z_{})^{0.8}`$, where $`Z`$ is the mass fraction of metal in the STAR.
#### 2.3.2 Mass ejection
In our code, the mass ejection due to stellar winds is combined with the mass ejection due to Type II supernova. Thus, the mass ejection rate per STAR is written as
$$M_{\mathrm{eject}}=m_{\mathrm{SNII}}R_{\mathrm{SNII}},$$
(20)
where $`m_{\mathrm{SNII}}`$ is the average mass ejected by stellar winds and Type II supernovae defined as
$$m_{\mathrm{SNII}}=\frac{{\displaystyle _{M_{\mathrm{ms}}}^{M_{\mathrm{ma}}}}m\mathrm{\Phi }(m)𝑑m}{{\displaystyle _{M_{\mathrm{ms}}}^{M_{\mathrm{ma}}}}\mathrm{\Phi }(m)𝑑m}m_{\mathrm{NS}},$$
(21)
Here $`m_{\mathrm{NS}}`$ is the mass that is locked up in the neutron star and assumed to be 1.4 $`M_{}`$. The fraction of heavy metal in $`M_{\mathrm{eject}}`$ is computed by the nucleosynthesis yield of Type II supernovae (Tsujimoto et al. 1996; Nomoto et al. 1997).
#### 2.3.3 Summary
We assume that the feedback phase continues for $`\tau (M_{\mathrm{ms}})=4.3\times 10^7`$ yr from the formation of each STAR and is divided into two phases: a stellar wind phase and a supernova phase. The stellar wind phase continues for $`\tau (M_{\mathrm{ma}})=5.4\times 10^6`$ yr, during which only the energy ejection from STARs is included; the ejected mass is included in the supernova phase for simplicity. The supernova phase begins at $`t=\tau (M_{\mathrm{ma}})`$ and ends at $`t=\tau (M_{\mathrm{ms}})`$. During the supernova phase, the energy ejection is given by Eq. (16). The mass ejection is the sum of the contributions by the stellar winds and Type II supernovae.
The thermal energies, gases, and heavy elements from stellar winds and supernovae are smoothed over neighbor particles of the STAR within a neighbor radius of $`R_f`$ (feedback radius). We treat $`R_f`$ as a parameter to meet the observational constraints. Such neighbor particles are called “in heating phase”. When the gas particles are in heating phase, we assume that the cooling is suppressed as proposed in Mori, Yoshii, & Nomoto (1999). This treatment produces the high temperature region around the STAR. Thus, the star formation is forbidden in the gas particles in heating phase.
## 3. Proto-Globular Cloud formation
We first examine the radiative condensations, which occur in a wide range of astrophysical circumstances from solar prominence to interstellar clouds (Meerson 1996). Radiative condensations in optically thin plasma have been considered by many authors since the pioneering work by Parker (1953) and Field (1965). The scale length of gravitational instability in a collapsing proto-galaxy is much larger than the radii of globular clusters (Lin & Murray 1996). As will be shown in the following sections, the scale length of radiative condensations is comparable to the radii of globular clusters. Thus, in a collapsing proto-galaxy, radiative condensations may be the only mechanism to form a PGC (Fall & Rees 1985; Lin & Murray 1996).
### 3.1. Radiative condensations in a collapsing proto-galaxy
The characteristic equation for the growth rate of radiative condensations, $`n`$, is obtained from the linearized equations for perturbations as
$$n^3+n^2c_s\left(k_T+\frac{k^2}{k_K}\right)+nc_s^2k^2+\frac{c_s^3k^2}{\gamma }\left(k_Tk_\rho +\frac{k^2}{k_K}\right)=0,$$
(22)
where $`c_s`$ is the sound speed, $`\gamma `$ is the ratio of the specific heats, $`k=2\pi /\lambda `$ is the wavenumber of the perturbation, $`k_\rho `$ and $`k_T`$ are the wavenumber of sound waves whose frequencies are equal to the growth rate of isothermal and isochoric perturbation, respectively, and $`k_K`$ is the inverse of the scale length of thermal conduction (Field 1965). $`k_\rho `$, $`k_T`$ and $`k_K`$ are expressed as
$$k_\rho =\frac{\mu (\gamma 1)\rho _0}{c\mathrm{k}_\mathrm{b}\mathrm{T}_0}\frac{\mathrm{\Lambda }(T_0)}{\mu ^2},$$
(23)
$$k_T=\frac{\mu (\gamma 1)}{c\mathrm{k}_\mathrm{b}}\frac{\rho _0}{\mu ^2}\frac{d\mathrm{\Lambda }}{dT},$$
(24)
$$k_K=\frac{c\mathrm{k}_\mathrm{b}\rho _0}{\mu (\gamma 1)\kappa },$$
(25)
where $`T_0`$ and $`\rho _0`$ are the equilibrium temperature and density, respectively, k<sub>b</sub> is the Bolzman constant, and $`\mathrm{\Lambda }(T)`$ is the cooling function (Figure 1). We assume $`\gamma =5/3`$ and $`\kappa =5.6\times 10^7T^{2.5}`$ erg s<sup>-1</sup> K<sup>-1</sup> cm<sup>-1</sup>. Solving Eq. (22) as a cubic equation of $`n`$ for different $`k`$, we obtain the dispersion relation between $`n`$ and $`k`$. In applying to our galaxy, we adopt $`T_01.0\times 10^6`$ K and $`\rho _01.7\times 10^{24}`$ g cm<sup>-3</sup> (Fall & Rees 1985). With these values, the dispersion relation has a peak at some $`k`$ . In the present paper, we assume that the scale for the maximum growth rate is typical scale of a PGC. To obtain the typical scale of a PGC for different $`T_0`$ and $`\rho _0`$, equation (22) is viewed as a quadratic in $`k^2`$ (see Section II (d) in Field (1965)).
The results for different $`T_0`$ and $`\rho _0`$ are presented in Figure 3. In applying to our galaxy, we obtain the scale length of $`600`$ pc. This scale length is consistent with the following simple estimate. For the adopted cooling function, the wavelengths for $`k_\rho `$ and $`k_K`$ with $`T_01.0\times 10^6`$ K and $`\rho _01.7\times 10^{24}`$ g cm<sup>-3</sup> are respectively obtained as
$$\lambda _\rho >10^3\text{pc}\text{and}\lambda _K<1\text{pc}.$$
(26)
This implies that the perturbation with a scale greater than $`10^3`$ pc is dumped owing to the limit of the sound speed, and the perturbation with a scale smaller than 1 pc is also dumped by thermal conduction. Thus the typical scale of a PGC in our galaxy is estimated to be several hundreds pc and the mass of a PGC ranges from 10<sup>7</sup> to 10<sup>8</sup> $`M_{}`$.
### 3.2. Density profile of a PGC
The estimated scale of a PGC is larger than the present radius of globular clusters (10 - 100 pc). This implies that during radiative condensations, a PGC shrinks before star formation begins or star formation in a PGC occurs in the central region of the cloud. If we consider the metal-free PGC where the ionization equilibrium is achieved, the cooling time, $`t_c`$, of the PGC is much shorter than the dynamical time, $`t_d`$ ($`t_c/t_d`$ 0.25 for $`T_01.0\times 10^6`$ K and $`\rho _01.7\times 10^{24}`$ g cm<sup>-3</sup>). After the short period of radiative condensations ($`t_c`$), the PGC becomes a warm dense cloud ($`T10^4`$ K) surrounded by a hot thin gas with $`T_0`$ (Fall & Rees 1985).
The qualitative discussion on the evolution of such a metal-free PGC has been presented by Lin & Murray (1996). They argued that in the first stage of collapsing galaxy, the assumption of ionization equilibrium is not valid because of little radiative emission. If the ionization equilibrium is not achieved, the temperature of high density regions can be as low as $`10^2`$ K due to hydrogen molecular cooling. In this case, the high density cold cloud with $`T=10^2`$ K is surrounded by a warm gas. The critical mass of the isothermal sphere confined by an external pressure depends on the cloud temperature (Ebert 1955; McCrea 1957). If the temperature of the cloud is $`10^2`$ K, the critical mass is $`10^2`$ $`M_{}`$. According to Lin & Murray (1996), such a small cloud as $`M10^2`$ $`M_{}`$ collapses to make the first stars in the proto-galaxy. These first generation stars are the source of radiative emission to maintain the ionization equilibrium. If such radiation continues long enough, a cold PGC will evolve almost isothermally with $`T10^4`$ K and form the isothermal profile of $`\rho r^2`$. The core size of the PGC will decrease from the initial size of several hundreds pc to a size comparable to the observed globular cluster. According to these arguments, we can assume that the density structure of a PGC has a form of $`\rho r^2`$.
For the initial conditions of the PGC, we use the King profile with fixed mass of $`M10^6`$ $`M_{}`$ which nearly equals to the Jeans mass of the current conditions ($`\rho _0`$ and $`T_0`$). We assume that this sphere represents the inner regions of the PGC. We examine three different cases for a initial radius of $`R_i=`$ 150, 200, and 300 pc. The initial gas density profiles of the three cases are presented in Figure 4. For the smaller $`R_i`$, the initial concentration of the inner region of the PGC is higher. The temperature of the sphere is assumed to be $`T10^4`$ K. The properties of the inner region of the PGC, which are the initial conditions of our calculations, are summarized as follows:
* The mass of the clouds is $`M10^6`$ $`M_{}`$.
* The radius of the clouds is $`R_i`$ 150 - 300 pc.
* The clouds is an isothermal sphere with $`T10^4`$ K.
* Initially, the velocity of the cloud is zero.
In the following calculations, the initial number of gas particles is $`5000`$ so that the initial mass of the one gas particle is $``$ 200 $`M_{}`$. The dependence of varying the initial number of gas particles is discussed in Section 4.3. In all cases, the gravitational softening length for all particles is set to be 1 pc and fixed during the calculations.
## 4. Results
First, we note on the feedback radius ($`R_f`$). We use the Strömgren radius ($`R_{\mathrm{St}}`$) of a typical OB star as $`R_f`$. The typical value is $`R_{\mathrm{St}}`$ 10 - 100 pc, where the density of the ISM is $``$ 1 cm<sup>-3</sup> (Osterbrock 1974). The radius depends on the density of the ISM as $`R_{\mathrm{St}}n^{2/3}`$. In the central region of our initial models, $`n`$ ranges from 100 cm<sup>-3</sup> to 1000 cm<sup>-3</sup>. Thus, we can estimate $`R_f`$ $`15`$ pc (using the largest value of $`R_{\mathrm{St}}`$). In the following two subsections, we choose $`R_f=`$ 3 pc for all gas particles. The results for both cases are summarized in Table 1. The results for different $`R_f`$ are described in Section 4.3.
### 4.1. Evolution of the metal-free PGC
In this section, we describe the evolution of the metal-free PGC for case C.
As presented in the previous sections, the metal-free PGC may evolve isothermally with $`T10^4`$ K. The isothermal evolution of PGC is terminated when the ionizing radiations from the first generation stars stop. Without ionizing radiation, the PGC cools efficiently by H<sub>2</sub> molecules. Our calculations start from that time when the PGC begins to cool. We assume $`M=10^6`$ $`M_{}`$, and the initial temperature of $`T=10^4`$ K and the King profile sphere. Since the size of the PGC at this stage is unknown, we examine three different cases for an initial radius of $`R_i=`$ 150, 200, and 300 pc.
In all three cases, the initial density of the central region is high enough for efficient cooling to occur, so that the temperature of the central region decreases rapidly to $`10^2`$ K. Due to the high density and low temperature ($`10^2`$ K), a burst of star formation occurs in the central region in all cases. After the first star formation, the central region becomes the heating region and the temperature of the central region increases gradually.
At $`t`$ 6 Myr, Type II supernovae start to occur and the temperature of the heating region increases rapidly to T $`10^6`$ K. The left panel of Figure 5 shows the evolutionary changes in the central temperature for $`R_i`$ = 150 pc. Such a high temperature region expands and the central density begins to decrease (see the right panel of Figure 5). The expansion of the central region leads to the formation of a shell structure as seen in the evolution of the gas density profile (Figure 6). The density profile at $`t=`$ 6 Myr, clearly shows the shell structure. After that time, the shell expands outwardly (see the right panel of Figure 6).
The star formation after the shell formation occurs not in the central region but in the shell. Figure 7 show the change in the radius of the star forming region as a function of time for $`R_i`$ = 150 pc. We can see that the star forming region moves outward. At $`t=`$ 10 Myr, where we stop the computation for $`R_i=`$ 150 pc, the stellar mass reaches $`1.3\times 10^5`$ $`M_{}`$. This means that the star formation efficiency is $``$ 13 %. The bound stellar mass at $`t=`$ 10 Myr is $`10^5`$ $`M_{}`$ and the mass is comparable to the typical mass of globular clusters ($`10^5`$ $`M_{}`$). At $`t=`$ 10 Myr, the gas is almost removed from the central region. Figure 8 shows the projected particle positions of STAR particles (left panel) and the stellar density profile (right panel) at $`t=`$ 10 Myr. The STAR particle shows elongated shape (bar like shape). Such a shape is caused by the radial orbit instability of the STARs formed at the large radius (Palmer & Papaloizou 1987). The central density of STARs is as high as $``$ 100 $`M_{}`$ pc<sup>-3</sup> and the central velocity dispersion of the STARs is $`3`$ km s<sup>-1</sup>. This value of velocity dispersion is smaller than the observed value (Dubath, Meylan & Mayor 1997). We need to follow further evolution of the star clusters for proper comparison, which is not feasible with present code. The number of STAR particles at $`t=`$ 10 Myr is $``$ 2500.
For $`R_i=`$ 200 pc and 300 pc, the overall evolution is similar to the case for $`R_i=`$ 150 pc and the stellar masses at $`t=`$ 10 Myr are $`8.5\times 10^4`$ $`M_{}`$ and $`5.6\times 10^4`$ $`M_{}`$, respectively. The bound stellar mass are $`5\times 10^4`$ $`M_{}`$ for both case. These results are summarized in Table 1. The initial concentration correlates with the final stellar mass and the final concentration of stellar system. In order for the stellar system as massive as 10<sup>5</sup> $`M_{}`$ to form, the PGC should be as compact as $`R_i<200`$ pc ($`\rho _\mathrm{c}>10^{22}`$ g cm<sup>-3</sup>) for the metal-free condition.
### 4.2. Evolution of the metal-rich PGC
The chemical composition of stars in globular clusters is one of the most crucial quantities to constrain the model for the globular cluster formation. The assumption that PGC initially has some heavy elements is quite reasonable. Actually, all globular clusters in our galaxy have some metals of \[Fe/H\] $`=2.25`$ \- 0. The metallicity distribution of the globular clusters shows bimodal distributions (Harris 1991). From this fact, we choose the initial metallicity of the cloud ranging from \[Fe/H\] $`=2`$ to 0. For comparison, we will evolve the lower metallicity (\[Fe/H\] $`=3`$) cloud. The results presented in this section are obtained for case B.
Using the same initial conditions presented in section 4.1, we evolve the PGC for different metallicity. For the initial radius of the cloud, we set $`R_i=`$ 300 pc. The first star formation occurs in the central region and the STARs heat up the surrounding matter to gradually increase the temperature of the central region. At $`t`$ 6 Myr, Type II supernovae begin to occur and the temperature of the heating region increases rapidly to $`T10^6`$ K. Figure 9 shows the evolutionary changes in the central temperature (left panel) and gas density (right panel) for \[Fe/H\] $`=2`$. The evolution for \[Fe/H\] $`=2`$ is very similar to the evolution of the metal-free cloud. After $`t`$ 6 Myr, the central high density region makes the shell structure as in the metal-free case. At $`t=`$ 10 Myr, the bound stellar mass reaches $`1.0\times 10^5`$ $`M_{}`$ for \[Fe/H\] $`=2`$.
For higher metallicity, i.e., \[Fe/H\] $`1`$, details of the evolution are somewhat different. We compare the central temperature evolution for different metallicity in Figure 10. Due to the different cooling rate, the initial decrease in the temperature is larger for higher metallicity. This difference leads to different star formation history as shown in Figure 11. For \[Fe/H\] $`=0`$, the initial star burst is very intense and then the SFR decrease sharply becoming lower than the low metallicity case after $`t=`$ 1 Myr. This is because heating due to the stellar winds much larger for higher metallicity (see Section 2.3.1). The rise in the SFR after $`t=`$ 6 Myr for \[Fe/H\] = 0 is caused by Type II supernovae.
The SFR for \[Fe/H\] $`=1`$ is almost constant during the evolution. For \[Fe/H\] $`=2`$, the first star formation occurs at $`t`$ 0.5 Myr because of the lower cooling rates. The SFR after $`t=`$ 2 Myr is almost constant but lower than for \[Fe/H\] $`=1`$.
The bound stellar mass at $`t=`$ 10 Myr is $`1.5\times 10^5`$ $`M_{}`$ and $`1.0\times 10^5`$ $`M_{}`$ for \[Fe/H\] $`=1`$ and 0, respectively. The reason for such a metallicity dependence is summarized below. The bound stellar mass, the central stellar density, and the central velocity dispersion are summarized in Table 1. The results depend on the metallicity as follows:
1. The initial metallicity significantly affects the star formation history. This difference makes the subsequent evolution different.
2. For the higher metallicity, the final stellar mass at $`t=`$ 10 Myr is larger because of the more efficient cooling rate.
3. The bound stellar mass is not an increasing function of the initial metallicity. This is because heating due to the stellar winds is larger for higher metallicity.
4. For lower metallicity (\[Fe/H\] $`=3`$), only $`3\times 10^3`$ $`M_{}`$ stars form. This implies that \[Fe/H\] $`2`$ is necessary to form globular clusters of $`10^5`$ $`M_{}`$ if $`R_i`$ 300 pc.
### 4.3. Parameter dependence
In this subsection, we describe the dependence of the results on various numerical parameters, e.g., the initial particle number ($`N_i`$) and the feedback radius ($`R_f`$). We use the metal-rich cloud of \[Fe/H\] $`=1`$ as a reference model, where $`N_i=`$ 5000 and $`R_f=`$ 3 pc.
First, we describe the dependence on the initial particle number. With larger (smaller) number of initial gas particles, the initial masses of the gas particles are smaller (larger). There is no strong dependence on the mass of the gas particles in our star formation recipes because Eq. (12) dose not include the mass of the gas particles. However, the maximum density ($`\rho _{\mathrm{max}}`$) that can be represented in the SPH method depends on the mass of the gas particles, which may produce weak dependence on $`N_i`$. To confirm this, we evolve the model with $`N_i=10^4`$. The overall evolution is almost indistinguishable to the reference model. The stellar mass at $`t=`$ 10 Myr becomes $`1.2\times 10^5`$ $`M_{}`$, which is slightly smaller than $`1.3\times 10^5`$ $`M_{}`$ in the reference model. For $`N_i=2500`$ (a half of the reference model), we obtain the stellar mass of $`1.5\times 10^5`$ $`M_{}`$. In Figure 12, we compare the gas density profiles at $`t=`$ 10 Myr for different $`N_i`$. The position of the shell for the model with $`N_i=10^4`$ is different from other models because of the smaller stellar mass at the same epoch. However, all three models show the similar density profiles so that we conclude that the dependence on the initial gas particle number is week. Thus, $`N_i`$ 5000 is sufficient for the current numerical model.
Secondary, we describe the dependence on the size of the feedback radius $`R_f`$. In our model, the cooling and the star formation are suppressed for the gas particles in heating phase. We therefore expect that the size of $`R_f`$ affects the star formation history. With larger $`R_f=`$ 5 pc, the star formation rates are smaller and the final stellar mass ($`0.9\times 10^5`$ $`M_{}`$) is smaller than $`1.3\times 10^5`$ $`M_{}`$ in the reference model ($`R_f=`$ 5 pc). On the other hand, we obtain the stellar mass of $`1.9\times 10^5`$ $`M_{}`$ with smaller $`R_f`$ (= 2 pc). Within the reasonable range of $`R_f`$ around 3pc, the final stellar mass changes by a factor of 0.7 - 1.5.
### 4.4. Summary of the numerical results
The results of our calculations are summarized as follows (see also Table 1) :
1. In all cases, the overall evolution is similar. Initially, the star burst occurs in the central region. The central star burst is halted by the heating due to Type II supernovae. The heating cause the expansion of the central region and forms a shell structure. The subsequent star formation occurs in the shell.
2. For the metal-free collapse case, the stellar mass at $`t=`$ 10 Myr correlates with initial concentrations. To obtain the star cluster as massive as globular clusters ($`10^5`$ $`M_{}`$), the initial concentration of the PGC must be large enough, i.e., $`R_i<200`$ pc.
3. For the metal-rich collapse, the initial metallicity significantly affects the evolution. To obtain the star cluster as massive as globular clusters, the initial metallicity must be as larger as \[Fe/H\] $`2`$. If the initial metallicity is low (\[Fe/H\] $`<2`$), very few stars form.
4. This suggests that during the initial phase of the galaxy formation, i.e., when the ISM contains little heavy element, the formation efficiency of the globular cluster is low.
5. The results is not strongly affected by the initial number of gas particles. $`N_i`$ 5000 is sufficient for the current numerical model. The dependence on the feedback radius is more evident.
## 5. Discussion
### 5.1. Star formation in the shell
In all our numerical models, a shell-like gaseous structure is formed. The formation of the shell-like structure of gases has also been reported in the simulations of the formation of dwarf elliptical galaxies (Mori et al. 1997; Mori, Yoshii, & Nomoto 1999). They argued that the difference in the density structure between normal ellipticals (de Vaucouleurs law) and dwarf ellipticals (exponential law) can be explained by the formation of such a shell-like structure and the star formation in the shell. The most crucial difference between the normal elliptical and the dwarf is their mass, and the less massive galaxy is more strongly affected by the energy inputs from stars. Similar argument may be applicable to our numerical models. In the case of the globular cluster formation, however, the mass is even smaller than the dwarf galaxies. Owing to this fact, the effect of the energy inputs from stars is more drastic so that the star formation in the shell is much less efficient than in the dwarf galaxies. Also, the stars formed in the shell is not gravitationally bound due to the outward velocity of the shell. Thus, the bound globular cluster consists of the stars formed before the shell formation. The stellar density of such stars is not affected by the shell formation and the star formation in the shell.
The chemical composition of stars in globular clusters is one of the most crucial quantities to constrain the model for the globular cluster formation. The stars in a globular cluster have almost the same heavy elements abundances (Suntzeff 1993). This small dispersion in metallicity ($`\sigma [\mathrm{Fe}/\mathrm{H}]`$) suggests that the formation period of globular clusters is so short that the stars in globular clusters can be regarded to form almost simultaneously as shown in our numerical models.
Brown, Burkert & Truran (1991, 1995) suggested that the second generation stars would form in the shell and the self-enrichment could occur there. In our model, the star formation takes place in the shell after $`t=`$ 5 Myr (see Figure 7). However, when the stars form in the shell, the shell has not been enriched with newly ejected heavy metal as shown in the solid line in Figure 7. This line shows the radius of the metal-enriched region defined as $`r=(\mathrm{\Sigma }r_im_i)/(\mathrm{\Sigma }m_i`$) by summing over the metal-enriched gas particles. The metal-enriched region expands outwardly but does not reach the star forming region, which implies that the self-enrichment dose not occur in our model (we note that the star formation is suppressed in the gas particles near the STARs in our numerical model as described in section 2.3.3). In the present model, no external medium outside the cloud is included, because of technical difficulties in the SPH method (see, however, Nagasawa & Miyama (1987) for a possible improvement). If there exists external medium outside the cloud, the density of the shell would be higher and the star formation history would be different; this possibility needs further study to confirm.
### 5.2. Failed Globular clusters
When a PGC is an initially metal-free cloud or the initial metallicity of a PGC is low, the resulted mass and the central stellar density are lower than the observed mass and density of the globular clusters. Such “failed” globular clusters might be the field halo stars. Another possibility is a open cluster. Typical age of open clusters in our galaxy is $`\tau <10^9`$ yr. However, there also exist such old open clusters as $`\tau >10`$ Gyr (Friel 1995). The age of the most old open cluster is comparable to the age of globular clusters. The central density of the “failed” globular cluster is as low as the central density of typical open clusters. Thus, the formation processes of such old open clusters may be the similar to our model of globular cluster formation but starting from lower initial concentration. Very old open clusters might be the debris of “failed” globular clusters and there might have existed many more open clusters at the formation of our galaxy.
## 6. Conclusions
For the processes of globular clusters formation, only the qualitative scenarios have been discussed previously (Fall & Rees 1985; Lin & Murray 1996). In this paper, we present the first attempt to simulate the globular cluster formation with three-dimensional hydrodynamical method, which includes the star formation and its feedback effects. We assume that, in the collapsing galaxy, isothermal cold clouds form through thermal condensations and become proto-globular clouds. We obtain the size of proto-globular clouds by means of the linear stability analysis (Figure 3) and compute the evolution of the inner region of the PGC starting from various initial radius $`R_i`$. The results of our calculations are summarized as follows:
1. In order for the globular cluster-like system to form from a metal-free PGC, the initial concentration of the PGC must be large enough.
2. It is required that the metallicity of a PGC is high enough to produce the globular cluster-like system. In our calculations, the required metallicity is estimated to be \[Fe/H\] $`2`$.
3. In all cases, the shell like structure of the gas forms. Although the star formation occurs in the shell, the self-enrichment is not seen to occur.
Based on the earlier qualitative works and the present quantitative results, the processes of globular clusters formation in the proto-galaxy can be understood as follows:
1. In the collapsing proto-galaxy, the first generation stars of $`M10^2`$ $`M_{}`$ form due to the efficient cooling by hydrogen molecules. Such population III stars eject the gas with heavy elements and chemically contaminate the proto-galaxy gases.
2. Such first generation stars radiate dissociative photons, and the entire proto-galaxy is settled down to ionization equilibrium.
3. With equilibrium cooling, the density perturbation grows due to thermal instability. and forms an isothermal cloud with a density structure of $`r^2`$. Such clouds are the proto-cloud of globular clusters.
4. When the density of the PGC becomes high enough, the burst of star formation occurs. Some high density clouds produce the globular clusters, and others may produce the field stars and/or the halo open clusters.
5. During the formation of the galaxy, the formation efficiency of the globular cluster become significantly large when the metallicity of the PGC become as large as \[Fe/H\] $`2`$.
As mentioned in section 5.1, we do not take account of the effect of the external medium in the present calculations. If we include the effect of the external medium, the structure of the shell may be different and further star formation may occur in that shell. Brown, Burkert & Truran (1995) suggested that only a few supernovae per Myr is sufficient to reverse the contraction of the 10<sup>6</sup> $`M_{}`$ cloud. Such an effect of the external medium needs to be studied.
The globular cluster formation is considered to occur not only in the proto-galaxies but also in the present galaxies, e.g., LMC, SMC, and interacting galaxy NGC4038/NGC4039. In such environment, different formation processes would take place, which correspond to the tertiary models (Kumai, Basu, & Fujimoto 1993). In the case of NGC4038/NGC4039 (the Antennae galaxies), many globular clusters are being produced by galaxy merging. Detail formation processes of globular clusters in such systems needs to be studied. Also, the effect of the formation of many globular clusters on the galaxy merging process needs a further study.
We would like to thank the anonymous referee for the valuable suggestions to improve the manuscript. Also, we would like to thank G. Mathews, J. Truran, T. Shigeyama and S. Zwart for useful discussion and comments. This work has been supported by JSPS Research Fellowships for Young Scientists (7664), and in part by the grant-in-Aid for COE Scientific Research (07CE2002) of Ministry of Education, Science, Culture, and Sports of Japan. |
warning/0001/math0001173.html | ar5iv | text | # How Many Turing Degrees are There?
## 1. Formulation of the problem
We denote by $`_T`$ the Turing equivalence relation on $`𝒫()=\{X:X\}`$, which we identify with $`2^{}`$, viewing sets as characteristic functions. (We use the standard set-theoretic convention that $`n=\{0,1,\mathrm{},n1\}`$ for all natural numbers $`n`$.) Then $`_T`$ is a Borel (in fact $`\mathrm{\Sigma }_3^0`$) equivalence relation on $`2^{}`$. We denote by $`𝒟`$ the quotient space $`2^{}/(_T)`$, i.e., the set of Turing degrees.
Now consider general Borel equivalence relations on $`2^{}`$ or even arbitrary Polish (separable completely metrizable) spaces. We measure their complexity by studying the following partial (pre)order of Borel reducibility: if $`E,F`$ are Borel equivalence relations on $`X,Y`$ respectively, then a Borel reduction of $`E`$ into $`F`$ is a Borel map $`f:XY`$ such that
$$xEyf(x)Ff(y).$$
If such an $`f`$ exists we say that $`E`$ is Borel reducible to $`F`$ and denote this by
$$E_BF.$$
Let also
$$E_BFE_BF\&F_BE$$
(this defines the concept of bi-reducibility) and
$$E<_BFE_BF\&F_BE.$$
Let us say that a function $`f_{}:X/EY/F`$ is Borel if it has a Borel lifting, i.e., there is a Borel function $`f:XY`$ such that $`f_{}([x]_E)=[f(x)]_F`$ for all $`xX`$. Then it is clear that $`E_BF`$ is equivalent to the assertion that there is a Borel injection from $`X/E`$ into $`Y/F`$, which we express by saying that the Borel cardinality, $`|E|_B`$, of $`E`$ is less than or equal to to that of $`F`$; in symbols,
$$|E|_B|F|_BE_BF.$$
Then define
$$|E|_B=|F|_BE_BF,$$
i.e., $`X/E,Y/F`$ have the same Borel cardinality, and
$$|E|_B<|F|_BE<_BF,$$
i.e., $`X/E`$ has (strictly) smaller Borel cardinality then $`Y/F`$.
We are now ready to formulate our problem as follows, where, by abusing notation, we write below $`|𝒟|_B`$ instead of $`|_T|_B`$ and call this the Borel cardinality of $`𝒟`$, instead of $`_T`$:
Question: What is the Borel cardinality, $`|𝒟|_B`$, of the set of Turing degrees $`𝒟`$?
If we denote the classical (Cantor) cardinality of $`𝒟`$ by $`|𝒟|`$, then we have $`|𝒟|=||`$. However, it is not hard to see that the Borel cardinality of $`𝒟`$ is bigger than that of the continuum. Let $`=_X`$ be the identity relation on the Polish space $`X`$. So $`|=_{}|_B`$ is the Borel cardinality which naturally represents the classical cardinality of the continuum.
Fact. $`(_T)>_B(=_{})`$.
###### Proof.
It is standard that there is a perfect set of pairwise Turing incomparable subsets of $``$, so $`(=_{})_B(_T)`$. If on the other hand $`f:2^{}`$ is Borel and Turing-invariant, i.e., $`x_Tyf(x)=f(y)`$, then for each Borel set $`A`$, $`f^1(A)`$ is a Turing-invariant Borel subset of $`2^{}`$, so it has measure 0 or 1. It follows that, for each $`n`$, the $`n`$th digit in the decimal expansion of $`f(x)`$ is fixed on a set of measure 1. So there is a Turing-invariant Borel set of measure 1 on which $`f`$ is constant, therefore $`f`$ cannot be a reduction of $`_T`$ into $`=_{}`$. Thus $`(_T)_B(=_{})`$.∎
We now have our question but it is not clear yet what kind of answer we should expect. In what sense can we hope to compute $`|𝒟|_B`$? To understand this, we have to dig a little deeper into the theory of Borel equivalence relations.
For our purposes, a crucial property of the Turing equivalence relation is that it has countable equivalence classes. In general, we call a Borel equivalence relation countable if every one of its classes is countable. We will next review some basic facts of the theory of countable Borel equivalence relations, for which we refer the reader to the papers Kechris \[K2\], Dougherty-Jackson-Kechris \[DJK\], Jackson-Kechris-Louveau \[JKL\], Kechris \[K1\], and Adams-Kechris \[AK\].
(i) (Feldman-Moore \[FM\]) Every countable Borel equivalence relation is generated by a Borel action of a countable group.
More precisely, given a countable Borel equivalence $`E`$ on a Polish space $`X`$, there is a countable group $`G`$ and a Borel action $`(g,x)gx`$ of $`G`$ on $`X`$ such that, if $`E_G^X`$ is defined by
$$xE_G^XygG(gx=y),$$
then $`E=E_G^X`$.
In particular, $`_T`$ is given by a Borel action of a countable group on $`2^{}`$. It seems like an interesting, but somewhat vague, question to find out whether one can obtain such a representation that has some recursion theoretic significance.
###### Remark 1.1.
Using the Feldman-Moore theorem and related facts, within a Schröder-Bernstein argument, one can show that, for countable Borel equivalence relations $`E`$ and $`F`$, $`E_BF`$ is equivalent to the existence of a Borel bijection of $`X/E`$ with $`Y/F`$.
(ii) There is a universal countable Borel equivalence relation, in the sense of $`_B`$.
That is, there is a countable Borel equivalence relation $`E`$ such that, for any countable Borel equivalence relation $`F`$, we have $`F_BE`$. This $`E`$ is clearly unique, up to $`_B`$, and denoted by $`E_{\mathrm{}}`$.
An example of a universal countable Borel equivalence is given by the orbit equivalence relation of the shift action of $`F_2`$, the free group on two generators, on $`2^{F_2}`$ given by
$$gx(h)=x(g^1h),g,hF_2,x2^{F_2}.$$
(iii) There is a smallest, in the sense of $`_B`$, countable Borel equivalence relation on uncountable Polish spaces, namely $`=_{}`$.
So for every countable Borel equivalence relation $`E`$ on an uncountable Polish space, we have $`(=_{})_BE`$. If $`(=_{})_BE`$, we say that $`E`$ is smooth. For example, $`_T`$ is not smooth. Another example of a non-smooth countable Borel equivalence is the following one, defined on $`2^{}`$:
$$xE_0ynmn(x(m)=y(m)).$$
This turns out to be the smallest, in the sense of $`_B`$, non-smooth countable Borel equivalence relation. This is a particular instance of the general Glimm-Effros Dichotomy proved in Harrington-Kechris-Louveau \[HKL\], but this special case can already be derived from Effros \[E\].
(iv) (Glimm-Effros Dichotomy) If $`E`$ is a countable Borel equivalence relation which is not smooth, then $`E_0_BE`$.
(v) $`E_0<_BE_{\mathrm{}}`$.
Thus we have
$$(=_{})<_BE_0<_BE_{\mathrm{}}$$
and every other countable Borel equivalence relation on an uncountable space is in the interval $`(E_0,E_{\mathrm{}})`$.
(vi) (Adams-Kechris \[AK\]) There are continuum many pairwise incomparable, under $`_B`$, countable Borel equivalence relations.
We now have all the ingredients to formulate a precise conjecture, in response to the question about the Borel cardinality of $`𝒟`$. This was originally formulated (as a question) in Kechris \[K2\] and listed (as a conjecture) in Slaman’s list of Questions in Recursion Theory, item 2.3, posted in http://math.berkeley.edu/$``$slaman/.
Conjecture: $`_T`$ is a universal countable Borel equivalence relation, i.e., $`(_T)_BE_{\mathrm{}}`$.
## 2. Known results and implications
There is some information already available about the complexity of $`_T`$.
###### Theorem 2.1.
(Slaman-Steel \[SS\]) $`E_0<_B(_T)`$.
This has been strengthened in Kechris \[K1\] to show that $`_T`$ is not amenable and in Jackson-Kechris-Louveau \[JKL\] to show that $`_T`$ is not treeable, all indications that $`_T`$ is quite complex.
One of the intriguing implications of the conjecture that $`_T`$ is universal concerns the existence of unusual functions on the Turing degrees. Recall that we call a function $`f:𝒟^n𝒟`$ Borel if there is a Borel function $`F:(2^{})^n2^{}`$ such that
$$f([x_1]_T,\mathrm{},[x_n]_T)=[F(x_1,\mathrm{},x_n)]_T$$
for all $`x_1,\mathrm{},x_n2^{}`$, where $`[x]_T`$ is the Turing degree of $`x2^{}`$. A pairing function on $`𝒟`$ is a bijection $`,:𝒟^2𝒟`$.
Fact. If $`_T`$ is universal, then there is a Borel pairing function on $`𝒟`$.
###### Proof.
If $`E,F`$ are Borel equivalence relations on $`X,Y`$ respectively, let $`E\times F`$ be the Borel equivalence relation on $`X\times Y`$ given by
$$(x,y)(E\times F)(x^{},y^{})xEx^{}\&yFy^{}.$$
Clearly $`E_{\mathrm{}}\times E_{\mathrm{}}_BE_{\mathrm{}}`$, so, since $`E_{\mathrm{}}`$ is universal, $`E_{\mathrm{}}\times E_{\mathrm{}}_BE_{\mathrm{}}`$. Hence, if $`(_T)_BE_{\mathrm{}}`$, we have
$$(_T)\times (_T)_B(_T),$$
which shows that there is a Borel pairing function on $`𝒟`$.∎
The well-known Martin Conjecture (or the 5th Victoria Delfino problem), see Kechris-Moschovakis, Eds. \[KM\] or Slaman’s list, item 2.2, seeks to classify definable functions on $`𝒟`$, asymptotically, i.e., up to identification on a cone of degrees. One part of the conjecture asserts, in particular, that if a Borel $`f:𝒟𝒟`$ is not constant on a cone, then $`f(d)d`$ on a cone. We can now easily see the following:
Fact. If $`_T`$ is universal, then Martin’s Conjecture fails.
###### Proof.
Fix $`d_0d_1`$ in $`𝒟`$ and let $`,`$ be a Borel pairing function on $`𝒟`$. Let $`f_0(d)=d_0,d`$ and $`f_1(d)=d_1,d`$. Then $`f_i:𝒟𝒟`$ is Borel for $`i=0,1`$ and, if $`A_i=\mathrm{rng}(f_i)`$, then $`A_0A_1=\mathrm{}`$. Since $`_T`$ is countable, one can show that the inverse of the pairing function $`,`$ is also Borel, so the sets $`A_i`$ are Borel.
Clearly $`f_0`$ and $`f_1`$ are injective, so they are not constant on a cone. Thus, if Martin’s Conjecture were true, we would have that $`f_i(d)d`$ on a cone for $`i=0,1`$. Then $`A_0`$ and $`A_1`$ would be cofinal in the Turing degrees, so, by Borel Determinacy, each would contain a cone, contradiction. ∎
## 3. Some more questions and answers
There are of course several other notions of equivalence and degree studied in recursion theory, and similar questions and conjecture can be considered for them too. We will concentrate here on one of the finest, recursive isomorphism, and one of the coarsest, arithmetic equivalence.
Let $`S_{\mathrm{}}`$ be the group of permutations of $``$, and let $`S_r`$ be the subgroup consisting of all recursive permutations. We let $`_r`$ denote recursive isomorphism for subsets of $``$. Via our identification of $`𝒫()`$ with $`2^{}`$, we have for $`x,y2^{}`$:
$$x_ry\pi S_r(x\pi =y).$$
For any $`n\{2,3,4,\mathrm{}\}\{\}`$ we also define recursive isomorphism on $`n^{}`$ by
$$x_r^ny\pi S_r(x\pi =y),$$
so that $`(_r^2)=(_r)`$.
It is well-known that $`(_T)_B(_r)`$, because $`x_Tyx^{}_ry^{}`$, where $`x^{}`$ is the Turing jump of $`x`$. Hence, if $`_T`$ is universal, then $`_r`$ is universal; and proving that $`_r`$ is universal could be viewed as providing additional evidence that $`_T`$ is universal.
Finally, we denote by $`_A`$ the notion of arithmetic equivalence on $`2^{}`$. So $`(_r)(_T)(_A)`$.
Again, one can conjecture that $`_r`$ and $`_A`$ are universal. Here, though, we have some answers.
###### Theorem 3.1.
(Slaman-Steel, unpublished). Arithmetic equivalence, $`_A`$, is universal, i.e., $`(_A)_BE_{\mathrm{}}`$.
So arithmetical equivalence has a Borel pairing function, and the arithmetical analogue of Martin’s Conjecture fails.
The problem for recursive equivalence is still open, but there has been a lot of progress.
###### Theorem 3.2.
(Dougherty-Kechris \[DK\]). Recursive isomorphism on $`^{}`$ is universal, i.e., $`(_r^{})_BE_{\mathrm{}}`$.
This was very recently improved to
###### Theorem 3.3.
(Andretta-Camerlo-Hjorth \[ACH\]). Recursive isomorphism on $`5^{}`$ is universal, i.e., $`(_r^5)_BE_{\mathrm{}}`$.
However, it is not yet clear how to reduce 5 to 2.
Actually, Theorems 3.2 and 3.3 are much more general. In each case, one actually shows that there is a fixed subgroup $`S_0`$ consisting of primitive recursive (in fact much simpler) permutations such that the result is true if $`S_r`$ is replaced by any countable group $`S`$ with $`S_0SS_{\mathrm{}}`$.
There is one last problem related to Theorem 3.2, that has further interesting implications.
First recall that an action of a group $`G`$ on a set $`X`$ is called free if $`gxx`$ for any $`xX`$ and $`g1_G`$. Also recall from §2 that every countable Borel equivalence relation is induced by a Borel action of a countable group $`G`$. From considerations in ergodic theory, it turns out that it is not always possible to find a free such action that induces it; see Adams \[A\]. It has been observed though that every known example of a countable Borel equivalence relation $`E`$, which cannot be induced by a free Borel action of a countable group, admits an invariant Borel probability measure (measure for short). (A measure is invariant for $`E`$ if it is invariant for any Borel action of a countable group that generates it.) It has in fact been conjectured that this is always the case. In other words, a countable Borel equivalence relation which does not admit an invariant measure can be induced by a free Borel action of a countable group.
By using the arguments in §2 of Dougherty-Jackson-Kechris \[DJK\] and a theorem of Nadkarni \[N\], it can be seen that this last assertion is equivalent to the following:
(†) There is a universal countable Borel equivalence relation, which is induced by a free Borel action of a countable group.
We return now to Theorem 3.2. We have that $`_r^{}`$ is induced by the following Borel action of $`S_r`$ on $`^{}`$:
$$\pi x=x\pi ^1.$$
This action is not free, but its restriction to
$$[]^{}=\{x^{}:x\text{ is one-to-one}\}$$
is. It is natural to conjecture that Theorem 3.2 can be strengthened to the statement that $`(_r)[]^{}`$ is universal. If this turns out to be the case, this will also prove (†).
## 4. Some proofs
We will give here our proof of Theorem 3.2 (and a related result). This comes from the unpublished Dougherty-Kechris \[DK\]. Although Theorem 3.2 has now been superseded by Theorem 3.3, our proof uses different methods and may find other applications in the future.
As we indicated in §3, one has in fact a stronger result. For any subgroup $`S`$ of $`S_{\mathrm{}}`$, and any $`X`$, let for $`x,yX^{}`$:
$$x_S^Xy\pi S(x\pi =y).$$
So $`(_r^{})=(_{S_r}^{})`$. We call $`S`$ primitive recursive if $`S=\{g_n:n\}`$, with $`g(n,m)=g_n(m)`$ primitive recursive. We now have:
###### Theorem 4.1.
There is a primitive recursive countable group $`S_0S_{\mathrm{}}`$ such that for any countable group $`S`$ with $`S_0SS_{\mathrm{}}`$, we have that $`_S^{}`$ is a universal countable Borel equivalence relation. In particular this is true for $`_r^{}`$.
###### Proof.
To explain the basic idea, consider a countable infinite group $`H`$ and fix a one-to-one enumeration $`H=\{h_n:n\}`$ of it. Then any $`h_aH`$ corresponds to a permutation $`\stackrel{~}{a}S_{\mathrm{}}`$ given by $`h_{\stackrel{~}{a}(n)}=h_nh_a`$ (the right regular representation). Fix also a bijection $`,:^2`$ and let $`\pi _aS_{\mathrm{}}`$ be defined by
$$\pi _a(n,m)=\stackrel{~}{a}(n),m.$$
Now given an action $`(h,x)hx`$ of $`H`$ into a space of the form $`X^{}`$ and the corresponding equivalence relation $`E_H`$, define the function $`f:X^{}X^{}`$ by
$$f(x)(n,m)=(h_nx)(m).$$
Then we have
$`f(h_ax)(n,m)`$ $`=(h_n(h_ax))(m)`$
$`=(h_{\stackrel{~}{a}(n)}x)(m)`$
$`=f(x)(\stackrel{~}{a}(n),m)`$
$`=(f(x)\pi _a)(n,m);`$
hence, $`f(h_ax)=f(x)\pi _a`$. It follows that if $`H_0=\{\pi _a:a\}`$ (a countable subgroup of $`S_{\mathrm{}}`$), then
$$xE_Hyf(x)_{H_0}^Xf(y).$$
(\*)
Unfortunately, if $`S_{\mathrm{}}H^{}H_0,H^{}`$ a countable group, then we cannot, in general, replace $`H_0`$ by $`H^{}`$ in (\*) since it could be that $`f(x)_H^{}^Xf(y)`$ via some $`\pi H^{}H_0`$. After appropriately choosing $`H`$, $`X`$, and the action of $`H`$ on $`X^{}`$ (so that at least $`E_H`$ is universal), we will modify $`f(x)`$ to $`f^{}(x)(X^{})^{}`$, for some $`X^{}`$, by encoding in it some further information, so that even if $`f(x)_H^{}^X^{}f(y)`$ via some $`\pi H^{}H_0`$ we can still conclude that $`xE_Hy`$. In particular, although the $`X`$ we will start with will be finite, this encoding will require $`X^{}`$ to be infinite. Moreover, we will be forced to restrict the $`x`$’s to some subset of $`X^{}`$, say $`YX^{}`$, so we will also need to make sure that $`E_HY`$ is universal.
We will now implement this idea. We fix some notation first:
For any $`X`$ and countable group $`G`$, we have the shift action of $`G`$ on $`X^G`$ given by
$$gx(h)=x(g^1h).$$
This induces for any subgroup $`HG`$ an action of $`H`$ on $`X^G`$ and we denote the corresponding equivalence relation by $`E(H,X^G)`$. If $`G`$ is infinite, fixing a one-to-one enumeration of $`G`$, we can view this as an action of $`H`$ on $`X^{}`$.
Now fix a one-to-one enumeration $`\{g_n:n\}`$ of the free group $`F_2`$ on two generators, with $`g_0=1`$ where $`1`$ is the identity element of $`F_2`$. Define $`\stackrel{~}{a}`$ and $`\pi _a`$ as above by the formulas $`g_{\stackrel{~}{a}(n)}=g_ng_a`$ and $`\pi _a(n,m)=\stackrel{~}{a}(n),m`$, and let
$$S_0=\{\pi _a:a\}.$$
If $`\{g_n:n\}`$ and $`,`$ are chosen appropriately, then $`S_0`$ is primitive recursive. Fix also any countable group $`S`$ such that $`S_{\mathrm{}}SS_0`$; we will show that $`_S^{}`$ is universal. Say $`S=\{\rho _i:i\}`$.
We call $`i`$ bad if
(i) $`nmn^{}(\rho _i(n,m)=n^{},m)`$; and
(ii) if $`\rho _i(0,m)=n_m,m`$ for all $`m`$, then $`n_m\mathrm{}`$ as $`m\mathrm{}`$.
We can now easily define $`n_j^{(i)},m_j^{(i)}`$ for $`i,j`$ such that:
(a) $`0<n_j^{(i)}<n_{j+1}^{(i)}`$ and $`0<m_j^{(i)}<m_{j+1}^{(i)}`$;
(b) $`(i,j)(i^{},j^{})m_j^{(i)}m_j^{}^{(i^{})}`$;
(c) if $`i`$ is bad, then $`n_{m_j^{(i)}}=n_j^{(i)}`$.
Also, for the free group $`F_k`$ with $`k`$ generators and $`gF_k`$, $`m`$, let $`B_k(g,m)`$ be the ball of radius $`m`$ around $`g`$ in the tree of $`F_k`$; i.e., $`B_k(g,m)`$ is the set of all products $`gh`$ where $`h`$ is a word in $`F_k`$ of length at most $`m`$.
Now consider the shift action of $`F_2`$ on $`9^{F_3}`$ (9 is a large enough number here) and the Borel set $`A9^{F_3}`$ defined by
$`yAij`$ $`[[(g_{n_j^{(i)}}y)B_3(1,m_j^{(i)})=(g_{n_{j+1}^{(i)}}y)B_3(1,m_j^{(i)})]`$
$`g_{n_j^{(i)}}y=g_{n_{j+1}^{(i)}}y],`$
where $`1`$ is the identity element of $`F_3`$.
###### Lemma 4.2.
$`E(F_2,9^{F_3})A_B(_S^{})`$.
###### Proof.
Fix an injection $`c`$ from the countable set $`_m9^{B_3(1,m)}`$ to $``$. Now define $`f^{}:A^{}`$ by $`f^{}(x)=x^{}`$, where $`x^{}(n,m)=c((g_nx)B_3(1,m))`$. Thus $`x^{}(n,m)`$ encodes the values of $`g_nx`$ at the ball of radius $`m`$ around $`1F_3`$. In particular, $`x^{}(n,m)`$ encodes (i.e., uniquely determines) $`m`$ as well. (If we were to take $`f(x)`$ as in the intuitive explanation in the beginning of this proof, then $`f(x)(n,m)`$ would be just $`g_nx(p_m)`$, where $`\{p_m:m\}`$ is a one-to-one enumeration of $`F_3`$.)
We claim that
$$xE(F_2,9^{F_3})yx^{}_S^{}y^{},$$
which completes the proof.
$``$: Clearly $`y=g_axy^{}=x^{}\pi _a`$.
$``$: Say now $`\pi S`$ is such that $`y^{}=x^{}\pi `$, i.e., $`y^{}(n,m)=x^{}(\pi (n,m))`$. Since $`x^{}(n,m)`$ encodes $`m`$, it follows that there is a function $`\pi ^{}:`$ such that $`\pi (n,m)=\pi ^{}(n,m),m`$ for all $`n`$ and $`m`$; that is, the second coordinate is left fixed by $`\pi `$. (Note that all $`\pi _a`$ have this property, of course. By our encoding we have forced any $`\pi `$ as above to have it as well.)
We now have two cases:
(I) $`\pi ^{}(0,m)`$ does not tend to $`\mathrm{}`$ as $`m\mathrm{}`$. So there must exist a number $`\mathrm{}`$ such that, for infinitely many $`m`$, $`\pi ^{}(0,m)=\mathrm{}`$. For any such $`m`$, we have $`y^{}(0,m)=x^{}(\mathrm{},m)`$, i.e., $`yB_3(1,m)=(g_{\mathrm{}}x)B_3(1,m)`$; since there are arbitrarily large such $`m`$, it follows that $`y=g_{\mathrm{}}x`$, so $`xE(F_2,9^{F_3})y`$.
(II) $`\pi ^{}(0,m)\mathrm{}`$ as $`m\mathrm{}`$. So if $`\pi =\rho _i`$, then $`i`$ is bad. For any $`j`$, we have $`y^{}(0,m_j^{(i)})=x^{}(n_j^{(i)},m_j^{(i)})`$, i.e., $`yB_3(1,m_j^{(i)})=(g_{n_j^{(i)}}x)B_3(1,m_j^{(i)})`$; but we also have $`yB_3(1,m_{j+1}^{(i)})=(g_{n_{j+1}^{(i)}}x)B_3(1,m_{j+1}^{(i)})`$, and $`m_j^{(i)}<m_{j+1}^{(i)}`$, so we get $`(g_{n_j^{(i)}}x)B_3(1,m_j^{(i)})=(g_{n_{j+1}^{(i)}}x)B_3(1,m_j^{(i)})`$. So, since $`xA`$, we have $`g_{n_j^{(i)}}x=g_{n_{j+1}^{(i)}}x`$ for all $`j`$, i.e., $`g_{n_0^{(i)}}x=g_{n_1^{(i)}}x=g_{n_2^{(i)}}x=\mathrm{}`$. It follows that $`yB_3(1,m_j^{(i)})=(g_{n_0^{(i)}}x)B_3(1,m_j^{(i)})`$ for all $`j`$; since $`m_j^{(i)}\mathrm{}`$ as $`j\mathrm{}`$, we have $`y=g_{n_0^{(i)}}x`$, so $`xE(F_2,9^{F_3})y`$ again. ∎
It remains to show that $`E(F_2,9^{F_3})A`$ is universal. For that we will show that
$$E(F_2,2^{F_2})_BE(F_2,9^{F_3})A,$$
which is enough, since $`E(F_2,2^{F_2})`$ is universal (see, e.g., Dougherty-Jackson-Kechris \[DJK\]).
###### Lemma 4.3.
There is a Borel injection $`f:2^{F_2}9^{F_3}`$ with $`f(2^{F_2})A`$ which preserves the group action of $`F_2`$ (i.e., for all $`gF_2`$ and $`x2^{F_2}`$, $`f(gx)=gf(x)`$). So in particular
$$E(F_2,2^{F_2})E(F_2,9^{F_3})A.$$
To prove this lemma, we will need the following technical sublemma.
###### Sublemma.
For each $`wF_2\{1\}`$, there is a Borel injection $`f_w:2^{F_2}6^{F_2}`$ which preserves the group action of $`F_2`$ and satisfies
$$f_w(x)(g)=f_w(x)(gw)g^1x=w^1g^1x$$
for all $`gF_2`$ and $`x2^{F_2}`$.
We will assume this and complete the proof.
###### Proof of Lemma 4.3.
Let $`\{\alpha _1,\alpha _2\}`$ be the generators of $`F_2`$ and $`\{\alpha _1,\alpha _2,\alpha _3\}`$ the generators of $`F_3`$. Define $`f(x)`$ for $`x2^{F_2}`$ as follows:
(i) If $`gF_2`$, then $`f(x)(g)=x(g)`$.
(ii) If $`g=h\alpha _3^pg^{}`$, with $`hF_2`$, $`p>0`$, and $`g^{}`$ not starting with $`\alpha _3^{\pm 1}`$, then $`f(x)(g)=2`$.
(iii) If $`g=h\alpha _3^pg^{}`$, with $`h,g^{}`$ as in (ii) and $`p>0`$, $`pm_j^{(i)}`$ for all $`i,j`$, then $`f(x)(g)=2`$.
(iv) If $`g=h\alpha _3^{m_j^{(i)}}g^{}`$, with $`h,g^{}`$ as in (ii), then $`f(x)(g)=f_{w_j^{(i)}}(x)(h)+3`$, where
$$w_j^{(i)}=g_{n_j^{(i)}}g_{n_{j+1}^{(i)}}^1.$$
It is easy to check that $`f`$ is one-to-one and preserves the action of $`F_2`$. So it remains to verify that $`f(x)A`$.
So fix $`i,j`$ with
$$(g_{n_j^{(i)}}f(x))B_3(1,m_j^{(i)})=(g_{n_{j+1}^{(i)}}f(x))B_3(1,m_j^{(i)}).$$
If $`d=\alpha _3^{m_j^{(i)}}`$, then $`dB_3(1,m_j^{(i)})`$, so
$$f(x)(g_{n_j^{(i)}}^1d)=f(x)(g_{n_{j+1}^{(i)}}^1d),$$
thus
$`f_{w_j^{(i)}}(x)(g_{n_j^{(i)}}^1)`$ $`=f_{w_j^{(i)}}(x)(g_{n_{j+1}^{(i)}}^1)`$
$`=f_{w_j^{(i)}}(x)(g_{n_j^{(i)}}^1w_j^{(i)}).`$
By the sublemma, $`g_{n_j^{(i)}}x=(w_j^{(i)})^1g_{n_j^{(i)}}x=g_{n_{j+1}^{(i)}}x`$, so
$`g_{n_j^{(i)}}f(x)`$ $`=f(g_{n_j^{(i)}}x)`$
$`=f(g_{n_{j+1}^{(i)}}x)`$
$`=g_{n_{j+1}^{(i)}}f(x);`$
since $`i,j`$ were arbitrary, $`f(x)A`$. ∎
It remains to prove the sublemma.
###### Proof of Sublemma.
View $`F_2`$ as a rooted tree in the usual way (1 is the root of this tree, and there is an edge between $`g`$ and $`g\alpha _i`$ for any group element $`g`$ and generator $`\alpha _i`$). Thus $`x2^{F_2}`$ is a labeling of this tree using labels 0,1. Similarly for $`6^{F_2}`$. Then $`g^1x`$ is the same labeling except that the root of the tree is at $`g`$ instead of $`1`$. So the condition
$$g[g^1xw^1g^1xf_w(x)(g)f_w(x)(gw)]$$
just means that if $`x`$, viewed from root $`g`$, is different from $`x`$ viewed from $`gw`$, then the label of $`f_w(x)`$ at $`g`$ is different from the label of $`f_w(x)`$ at $`gw`$. Moreover, to guarantee that $`f_w(g^{}x)=g^{}f_w(x)`$ for each $`g^{}F_2`$, we will make sure that the value of $`f_w(x)`$ at any $`g`$ depends only on the labeling $`x`$ viewed from root $`g`$ (and not on $`g`$ itself).
Given $`x2^{F_2}`$ and $`gF_2`$, we have two cases:
(I) $`g^1x=w^1g^1x`$, i.e., $`x`$ looks the same from root $`g`$ and root $`gw`$ (note that this only depends on how $`x`$ looks from root $`g`$).
Then put $`f_w(x)(g)=x(g),0`$, where $`,`$ is a bijection of $`2\times 3`$ with $`6`$.
(II) $`g^1xw^1g^1x`$. So $`x`$ looks different from roots $`g,gw`$. In particular there is a least $`n=n_g(x)`$ so that for some $`i,j`$ and $`hF_2`$ of length $`n`$ we have $`x(gw^ih)x(gw^jh)`$. Clearly $`n_{gw^i}(x)=n_g(x)`$ for any integer $`i`$ (note that $`(gw^i)^1xw^1(gw^i)^1x`$ as well).
The functions $`p_j:B_2(1,n_g(x))2`$ given by
$$p_j(h)=x(gw^jh)$$
are thus not all equal. So fix $`p2^{B_2(1,n_g(x))}`$ with $`Z=\{j:p_j=p)\mathrm{}`$ and $`p`$ least such (in some ordering of $`2^{B_2(1,n_g(x))}`$ fixed in advance). The value of $`p`$ would be the same if we started with $`gw^i`$ instead of $`g`$; the set $`\stackrel{~}{Z}`$ we would get from $`gw^i`$ is a translate of $`Z`$ ($`j\stackrel{~}{Z}`$ iff $`j+iZ`$).
Also $`\{j:p_jp\}\mathrm{}`$. If $`Z`$ has a largest element $`i_0`$, let $`f_w(x)(g)=x(g),0`$, if $`i_0`$ is even, and $`f_w(x)(g)=x(g),1`$, if $`i_0`$ is odd. If $`Z`$ has no largest element but has a least element $`i_0`$, define $`f_w(x)(g)`$ the same way. Proceed similarly if $`Z`$ has a least or largest element. So assume both $`Z`$ and $`Z`$ are unbounded in both directions. Put
$$Z^{}=\{jZ:j+1Z\}.$$
Let finally $`f_w(x)(g)=x(g),0`$ if $`0Z^{}`$, $`f_w(x)(g)=x(g),1`$ if $`0Z^{}`$, but the least positive element of $`Z^{}`$ if odd, and $`f_w(x)(g)=x(g),2`$ if this least positive element is even.
This completes the definition of $`f`$; it is straightforward to verify that it has the desired properties. ∎
This completes the proof of Theorem 4.1. ∎
We conclude with another application of these ideas.
For a countable group $`G`$ consider the shift action of $`G`$ on $`X^G`$. We call $`xX^G`$ a left-free point if for all distinct $`g,g^{}G`$ there exists $`hG`$ such that $`x(hg)x(hg^{})`$. We call $`xX^G`$ a right-free or just free point, if for all distinct $`g,g^{}G`$ there exists $`hG`$ such that $`x(gh)x(g^{}h)`$; equivalently, $`gxg^{}x`$ for $`gg^{}`$, or simply $`gxx`$ for all $`g1_G`$. Denote by $`LF`$ the set of left-free points and $`F`$ the set of free points. Note that $`LF`$ and $`F`$ are Borel $`G`$-invariant subsets of $`X^G`$. If $`G`$ is abelian, clearly $`LF=F`$. But $`LF`$ and $`F`$ are very different for free groups in the following sense.
###### Theorem 4.4.
The equivalence relation $`E(F_3,4^{F_3})LF`$ is universal for countable Borel equivalence relations but $`E(F_3,4^{F_3})F`$ is not.
###### Proof.
The equivalence relation $`E(F_3,4^{F_3})F`$ is not universal because it is treeable; see Kechris \[K2\]. For the first assertion we will show that $`E(F_2,2^{F_2})_BE(F_3,4^{F_3})LF`$.
Fix a left-free point $`z_0`$ in $`\{2,3\}^{F_2}`$. Define then $`f:2^{F_2}4^{F_3}`$ by:
(i) If $`hF_2,f(x)(h)=x(h)`$.
(ii) If $`hF_2`$, express the reduced word for $`h`$ in the form $`h=h_1\alpha _3^{\pm 1}h^{}`$ with $`h^{}F_2`$, and put $`f(x)(h)=z_0(h^{})`$.
It is easy to check that $`xE(F_2,2^{F_2})yf(x)E(F_3,4^{F_3})f(y)`$. It remains to verify that $`f(x)LF`$. Let $`g`$ and $`g^{}`$ be distinct elements of $`F_3`$; we must find $`hF_3`$ such that $`f(x)(hg)f(x)(hg^{})`$.
Consider two cases:
(1) $`g^1g^{}F_2`$. Then let $`pF_2`$ be such that $`z_0(p)z_0(pg^1g^{})`$, and let $`h`$ be such that $`hg=\alpha _3p`$. Then $`f(x)(hg)=f(x)(\alpha _3p)=z_0(p)z_0(pg^1g^{})=f(x)(hgg^1g^{})=f(x)(hg^{})`$.
(2) $`g^1g^{}F_2`$. Let $`h=g^1`$. Then
$$f(x)(hg)=f(x)(1)=x(1)\{0,1\}$$
but
$$f(x)(hg^{})=f(x)(g^1g^{})=z_0(h^{})\{2,3\}$$
for some $`h^{}F_2`$. ∎ |
warning/0001/astro-ph0001204.html | ar5iv | text | # A large high-energy gamma-ray flare from the blazar 3C 273
## 1 Introduction
Since its discovery as an extragalactic object in 1962 (Schmidt (1963)), the quasar 3C 273 is one of the best studied Active Galactic Nuclei (AGN). With a redshift of 0.158 ($``$800 Mpc for H<sub>0</sub> = 60 km/s/Mpc) it is relatively close, and by being bright in all wavelength regions from radio to $`\gamma `$-ray energies, it is an excellent candidate for multiwavelength studies.
After six years of operation, the EGRET experiment aboard CGRO has now detected $`\gamma `$-ray emission from more than 70 AGN at energies above 100 MeV (e.g. Hartman et al. (1999)). These observations have dramatically changed our picture of these sources. With the exception of Cen A (Sreekumar et al. (1999)), all of them are identified with blazars (e.g. Mattox et al. (1997)), the AGN subgroup consisting of either flat-spectrum radio quasars or BL Lacertae objects. Two remarkable $`\gamma `$-ray characteristics of these sources are that 1) they are highly variable down to time scales of a day or even shorter, and 2) that during flaring states the $`\gamma `$-ray luminosity can dominate their bolometric power.
The quasar 3C 273 is one of these blazar-type $`\gamma `$-ray loud AGN. It was first detected at $`\gamma `$-rays by the COS-B satellite at energies above 50 MeV (Swanenburg et al. (1978)), and – until the launch of CGRO in 1991 – remained the only identified extragalactic point source at these energies. 3C 273 was redetected at $`\gamma `$-ray energies by the EGRET experiment in 1991 (von Montigny et al. (1993)). Analysing the first four years of EGRET data, von Montigny et al. (1997) found a time-variable $`\gamma `$-ray flux, consisting of detections as well as non-detections in individual observational periods. Spectral variability was observed as well, showing the trend of spectral hardening with increasing flux. The third EGRET source catalogue (Hartman et al. (1999)) lists 3C 273 with an average flux value of (15.4$`\pm `$1.8) $`\times `$ 10<sup>-8</sup> ph cm<sup>-2</sup> s<sup>-1</sup> for energies above 100 MeV and the time period between April ’91 and October ’95.
3C 273 was first discovered to be an emitter of low-energy $`\gamma `$-rays by COMPTEL in 1991 (Hermsen et al. (1993)). The source is frequently detected in individual CGRO pointings (e.g. Collmar et al. (1999)), however, non-detections occur as well proving time variability of the MeV-flux on time scales of months (Williams et al. (1995)). In time-averaged analyses 3C 273 is detected very significantly and shows in the 0.75-30 MeV band a soft spectrum, i.e. photon index $`\alpha >`$ 2 (E) in combined data (Collmar et al. (1996)). However, combining contemporaneous high-energy data reveals that the MeV-band is a transition region for the spectrum of 3C 273 showing a turnover from a harder ($`\alpha `$1.7) spectrum at hard X-ray energies to a softer one ($`\alpha `$2.5) at high-energy ($`>`$100 MeV) $`\gamma `$-rays (e.g. Lichti et al. (1995), von Montigny et al. 1997).
At hard ($`>`$50 keV) X-rays 3C 273 is always significantly detected by the Oriented Scintillation Spectrometer Experiment (OSSE) showing flux variations in the 50-150 keV band up to a factor of 8 during 5 years (McNaron-Brown et al. (1997)). A power-law spectrum with a photon index of roughly 1.7 (Johnson et al. (1995)) up to $``$1 MeV is typically observed. Above $``$1 MeV a spectral softening is found, consistent with the results of the multiwavelength campaigns. Recently, during the highest flux state as observed by OSSE, evidence for a low-energy spectral break at about 0.3 MeV was found (McNaron-Brown et al. (1997)) suggesting an anticorrelation between flux and break energy.
In this paper we report on 7 weeks of continuous $`\gamma `$-ray observations by the EGRET and COMPTEL experiments aboard CGRO in December 1996 and January 1997. In Sect. 2 we describe the observations and the data analyses, in Sect. 3 we give the results, and discuss their implications in Sect. 4. Finally, the conclusions are presented in Sect. 5.
## 2 Observations and data analysis
During Cycle 6 (October 15, 1996 to November 11, 1997) of its mission, CGRO was pointed continuously to the Virgo sky region for 7 weeks, beginning on December 10, 1996 and ending on January 28, 1997. The main target was the blazar 3C 279 which is located at ($`\alpha ,\delta )_{2000}`$ = (12<sup>h</sup>56<sup>m</sup>11<sup>s</sup>, -5.847’21.5”), at a distance of $``$10.5 from 3C 273. The relevant observational parameters are given in Table 1.
The spark-chamber telescope EGRET covers the energy range from $``$30 MeV to $``$30 GeV. The instrument and its calibration is described in detail by Thompson et al. (1993) and Esposito et al. (1999). The analysis of the EGRET data followed the standard EGRET procedure i.e. using count and exposure maps as well as predictions of the diffuse $`\gamma `$-ray background (e.g. Hunter et al. (1997)). The maps containing events with energies above 100 MeV were used for source detection and determination of the source position, and the ones of the 10 standard energy intervals for determination of the source spectrum by assuming a power law. The analysis applied the standard maximum-likelihood method (Mattox et al. (1996)) and spatial selections ($`<`$25 off axis). Empirical flux correction were applied for energies below 70 MeV.
The imaging Compton telescope COMPTEL covers the energy band $``$0.75 to $``$30 MeV. For a detailed description of the COMPTEL instrument see Schönfelder et al. (1993). The COMPTEL data have been analysed following the COMPTEL standard maximum-likelihood analysis procedures, which for point-source analyses including background generation are described in sufficient detail by Bloemen et al. (1994) and which derive quantitative source parameters like detection significances, fluxes, and flux errors. For consistency checks maximum-entropy images (see e.g. Strong et al. (1992)) have been generated as well.
The fluxes of 3C 273 were determined by simultaneously fitting further sources or source candidates which are indicated by the maps (see Fig. 1). This approach leads iteratively to simultaneous flux determinations of several sources, including the generation of a background model which takes into account the possible presence of further sources. This analysis has been carried out for the four standard COMPTEL energy bands. The flux results given in Sect. 3 have been derived with point spread functions assuming an E<sup>-2.0</sup> power-law shape for the sources, which is approximately the correct shape for the MeV-spectrum of 3C 273.
## 3 Results
### 3.1 Detections
The EGRET data analysis of the combined 7 weeks of Virgo data revealed a strong source consistent with the location of 3C 273 (Fig. 1). The overall detection significance at the position of the quasar for energies above 100 MeV is 10.4$`\sigma `$ assuming $`\chi _1^2`$-statistics for a known source. The 7-week mean flux (E$`>`$100 MeV) is (43.4$`\pm `$5.8) $`\times `$ 10<sup>-8</sup> ph cm<sup>-2</sup> sec<sup>-1</sup>, which is roughly 3 times the average flux listed in the third EGRET catalogue (Hartman et al. (1999)). A flux level in excess of that – (48.3$`\pm `$11.8) $`\times `$ 10<sup>-8</sup> ph cm<sup>-2</sup> sec<sup>-1</sup> in VP 308.6 – had been reported previously only once (Hartman et al. (1999)). 3C 273 is identified with the $`\gamma `$-ray source on the basis of its sky location.
Simultaneously with these EGRET findings, COMPTEL observed significant emission from the sky position of 3C 273 in three (1-3 MeV, 3-10 MeV, 10-30 MeV) out of its 4 standard energy bands. The source is not detected at the lowest COMPTEL energies (0.75-1 MeV). The overall detection significance is 7.7$`\sigma `$. The average flux values in the individual COMPTEL bands (see Table 2) are among the largest ever measured in this energy window, showing that the source was active in $`\gamma `$-rays during this period.
### 3.2 Time variability
To check for time variability we have subdivided the total 7 week observation into slices of the 7 individual VPs covering typically one week each (Table 1). 3C 273 was not detected by EGRET during the first two VPs (detection-significance threshold 3.5$`\sigma `$). It then appeared just above the detection threshold and increased further to show the largest $`\gamma `$-ray flare observed during the CGRO era. This maximum flux level of (77$`\pm `$20) $`\times `$ 10<sup>-8</sup> ph cm<sup>-2</sup> sec<sup>-1</sup> is reached in VP 610. Thereafter the flux returned to an intermediate level, which is clearly detectable. Fits of the two EGRET light curves shown in Figs. 2 and 3 (upper panel) assuming a constant flux resulted in $`\chi _{min}^2`$-values of 19.9 for the light curve containing 7 flux points (individual VPs) and 12.5 for the one containing 3 flux points. According to $`\chi ^2`$-statistics these values correspond to probabilities of 2.9 $`\times `$ 10<sup>-3</sup> and 1.9 $`\times `$ 10<sup>-3</sup>, respectively, for a constant flux, which convert to 3.0$`\sigma `$ and 3.1$`\sigma `$ evidence for a time-variable $`>`$100 MeV flux. The largest change in flux occurred between VPs 610 and 610.5 when the flux dropped by a factor of $``$4.1 within 7 days. This is the shortest time variability at $`\gamma `$-ray energies ever reported for 3C 273. The significance that the two flux values are different is 2.7$`\sigma `$. During the two-week outburst (VPs 609 and 610) the source reached a flux level even slightly in excess to the COS-B flux reported by Swanenburg et al. (1978).
We checked for time variability in the COMPTEL bands by subdividing the data into the same time slices as chosen for EGRET. No obvious time variability is visible in either energy band, however, the statistics in the different COMPTEL bands became marginal, resulting in large error bars on the flux values. To improve the statistical significance we combined individual VPs. We defined three time intervals which were selected according to the EGRET light curve: a pre-flare period (VPs 606-608) covering 3 weeks which we shall call A in the following, a flare period (VPs 609 and 610) covering two weeks which we shall call B, and a post-flare period (VPs 610.5 and 611) which we shall call C (see also Table 1). The simultaneous EGRET and COMPTEL 3C 273 fluxes for different energy bands and various time periods are listed in Table 2 and are plotted in Fig. 3. In contrast to EGRET, COMPTEL observes no obvious time variability. The same $`\chi ^2`$-procedure as applied to the EGRET data showed that the different COMPTEL flux values are consistent with a constant level as is also obvious from Fig. 3. In particular the COMPTEL light curves do not show a hint of increased $`\gamma `$-ray emission during the EGRET flaring period. This result suggests that the observed flare is either solely a high-energy ($`>`$30 MeV) phenomenon, or a time offset of at least 2 weeks between the EGRET and COMPTEL $`\gamma `$-ray bands is required.
### 3.3 Energy spectra
The EGRET spectral analysis followed the standard EGRET procedure (see Sect. 2). The energy range between 30 MeV and 10 GeV was subdivided into 10 energy intervals and the likelihood analysis was used to estimate the number of source photons in each energy bin. The data were fit to a single power-law model of the following form
$$I(E)=I_0(E/E_0)^\alpha \mathrm{photons}\mathrm{cm}^2\mathrm{s}^1\mathrm{MeV}^1$$
(1)
with the parameters $`\alpha `$ (photon spectral index) and $`I_0`$ (intensity at the normalization energy $`E_0`$). $`E_0`$ was chosen such, that the two free parameters are minimally correlated. We derived 1$`\sigma `$-errors on the parameters by adding 1.0 to the minimum $`\chi ^2`$-value (Mattox et al. (1996)).
This approach was applied to the sum of all data as well as to selected subsets (see Fig. 3, Table 1) to check for a possible trend in time. In addition we summed the subsets A and C having roughly equal flux levels, which we shall call D, to check for a possible spectral trend with flux by using the improved event statistics. The results of the spectral fitting are given in Table 3. First of all the spectra are well fitted by simple power-law functions: the average spectral index in the EGRET band is $`\alpha `$ = 2.40$`\pm `$0.14, which is comparable to previous results. For instance, von Montigny et al. (1997) found spectral indices in the range between 2.2 and 3.2 with a trend of spectral hardening with increased source flux. The EGRET fits show this trend as well: the spectrum is hardest during the flaring period. However, this is not significant because the error on the spectral slope is quite large.
To derive the COMPTEL fluxes of 3C 273, we have applied the standard maximum-likelihood method as described in Sect. 2. Background-subtracted and deconvolved source fluxes in the 4 standard energy bands have been derived by taking into account the presence of further $`\gamma `$-ray sources (3C 279 at $`\alpha `$/$`\delta `$ = 194.1/-5.8, 4C -02.55 (1229-021) at $`\alpha `$/$`\delta `$ = 188.0/-2.4) and source candidates (at $`\alpha `$/$`\delta `$ = 193.5/0.5and $`\alpha `$/$`\delta `$ = 173.5/8.5), showing some evidence in the maps (Fig. 1). We note that inclusion of these further sources and source candidates has only a marginal effect on the flux of 3C 273. Their inclusion or exclusion changes the derived 3C 273 fluxes only within its error bars. Assuming a power-law shape, the quasar is fitted typically with a spectral index of $``$2 throughout the COMPTEL energy band. For example, the sum of all data (’ALL’) which has the best statistics, yields $`\alpha `$ = 1.92$`\pm `$0.13 between 0.75 and 30 MeV, which is significantly harder than $`\alpha `$ = 2.41$`\pm `$0.14 as found in the EGRET range for the same observational period (Fig. 4). This result indicates a spectral hardening towards lower energies with the turnover starting at a few MeV. This spectral behaviour of 3C 273 is well known and has been reported previously (e.g. Lichti et al. (1995), von Montigny et al. 1997).
To take advantage of the contemporaneous observations of both instruments in neighbouring energy regimes, we combined the deconvolved EGRET and COMPTEL spectra for the different time periods. Fitting the whole energy range (0.75 MeV to 10 GeV) with a power-law model, we derive harder spectral slopes and increased $`\chi _{red}^2`$-values compared to fitting solely the EGRET range. This effect weakens if we exclude the lowest-energy (0.75-1 MeV) spectral point, and disappears when we subsequently remove the next (1-3 MeV) COMPTEL flux point from the fitting procedure. This behaviour is easily explained by the spectral turnover at low energies, which seems to affect the fits only below 3 MeV. Fitting a broken power-law model
$$I(E)=\{\begin{array}{cc}I_0(E/E_0)^{\alpha _2}\hfill & \text{if }E>E_b\hfill \\ I_0(E_b/E_0)^{\alpha _2}(E/E_b)^{(\mathrm{\Delta }\alpha \alpha _2)}\hfill & \text{if }E<E_b\hfill \end{array}$$
(2)
where I<sub>0</sub> describes the differential source flux at the normalization energy E<sub>0</sub>, $`\alpha _2`$ the high-energy spectral photon index, $`\mathrm{\Delta }\alpha `$ the break in spectral photon index towards lower energies ($`\mathrm{\Delta }\alpha `$ = $`\alpha _2`$ \- $`\alpha _1`$), and E<sub>b</sub> the break energy, provides consistent results. The best-fit value for the break energy for the sum of all data is found to be at $``$5 MeV, which, however, is not well defined due to the small lever-arm towards lower energies. Considering these facts, we conclude, that the EGRET power-law spectrum extends into the COMPTEL band down to $``$3 MeV before it is substantially altered by the spectral turnover. This is illustrated in Fig. 4. By chance the quasar was observed simultaneously in the X-ray band by the BeppoSax satellite. Haardt et al. (1998) published the X-ray results, i.e. fluxes and spectra, of 3C 273 covering a monitoring period of 4 days (January 13, 15, 17, and 22, 1997). These X-ray observations are coincident in time with the transition from the $`\gamma `$-ray flaring period to the moderate post-flare $`\gamma `$-ray level (Fig. 2). Fig. 4 shows a broad-band high-energy spectrum of 3C 273, containing the COMPTEL and EGRET spectra for the sum of all 7 weeks and the best-fit power-law shape for the quasi-simultaneous X-ray measurements (covering only one day) provided by Haardt et al. (1998).
To investigate this high-energy component in more detail, we fitted the different observational subsets between 3 MeV and 10 GeV with single power-law models (Fig. 5). We take advantage of this enlarged (with respect to only EGRET) energy band, for which the spectral index can be determined more accurately. So, for the 3 MeV to 10 GeV band, the trend that during the flare the spectrum hardens as suggested by the EGRET analysis (see above), is observed more significantly. Especially, if we compare the periods B and D. The 1$`\sigma `$ statistical errors in spectral index during the flare and non-flare intervals do not overlap anymore. The fit results for the EGRET band only and this enlarged band are given in Table 3, and the latter ones are shown graphically in the Figs. 5 and 6. Along the 7-week observation, 3C 273 is observed to have a steep spectrum at the beginning, which hardens during the two-week flaring period, and turns back to roughly the same shape in the 2 week post flare period (Fig. 7). This result is consistent with the constant flux observed at COMPTEL energies. The flare occurs mainly at energies above 100 MeV, not affecting the COMPTEL points which results in a hardening of the overall $`\gamma `$-ray spectrum.
From this spectral analysis, we conclude too, that we either observed a phenomenon which occurs solely at high $`\gamma `$-ray energies, or that there are time delays between the different $`\gamma `$-ray bands. The first case would require an additional spectral component triggered by some mechanism which is only effective at EGRET energies. The second possibility would require that the COMPTEL energies are either delayed by two weeks or would be in advance by three weeks with respect to EGRET.
## 4 Discussion
At the end of 1996 and early 1997 the CGRO experiments EGRET and COMPTEL observed the well-known quasar 3C 273 at $`\gamma `$-ray energies continuously for 7 weeks. The blazar was $`\gamma `$-ray active and therefore it was significantly detected by both experiments. Assuming isotropic emission, H<sub>0</sub>=60 km/s/Mpc, and a q<sub>0</sub>=0.5 cosmology, we derive an average 7-week luminosity in the EGRET band (100 MeV - 10 GeV) of $``$1.7$`\times `$10<sup>46</sup>erg/s, for the flaring period (period B in Table 2) $``$2.7$`\times `$10<sup>46</sup>erg/s, and $``$0.9$`\times `$10<sup>46</sup>erg/s for the periods outside the two-week flare. The 7-week average luminosity in the COMPTEL band (1-30 MeV) is derived to be $``$12$`\times `$10<sup>46</sup>erg/s.
The blazar was simultaneously observed in the X-ray band by the BeppoSax satellite in a monitoring fashion covering a period of 1.5 weeks with four individual pointings. In X-rays the quasar was about 15% brighter in the first observation than in the last one (Haardt et al. (1998)), which is consistent with the $`\gamma `$-ray behaviour above 100 MeV. In addition, they note that during these monitoring observations the source was, on average, a factor of 2 brighter as observed half a year earlier. This suggests a correlation of X- and $`\gamma `$-ray behaviour of 3C 273.
Simultaneous X-ray and $`\gamma `$-ray measurements provide the possibility of estimating some physical source parameters, if one assumes that both photon populations are generated co-spatially. This is not an unreasonable assumption, given the indication of correlated variability as mentioned above. It was observed also in other blazars like 3C 279 for example (Wehrle et al. (1998)). Using the simultaneous X-ray spectra of the BeppoSax Medium Energy Concentrator Spectrometers (MECS) published by Haardt et al. (1998) we derive a flux of $``$16$`\mu `$Jy at 1 keV for 3C 273. Applying the expression for the lower limit on the Doppler factor, $`\delta `$, given by Dondi & Ghisellini (1995) and assuming H<sub>0</sub>=60 km/s/Mpc, we derive
$$\delta \left(684t_{var}^1E_\gamma ^\alpha \right)^{\frac{2}{(4+2\alpha )}},$$
(3)
where t<sub>var</sub> is the variability time scale in days, E<sub>γ</sub> the highest unabsorbed $`\gamma `$-ray energy in GeV, and $`\alpha `$ the spectral index in X-rays. With a variability scale of one week as observed by EGRET at a significance level of 2.7$`\sigma `$, an E<sub>γ</sub> of 1 GeV, and an $`\alpha `$ of 0.6 (energy index) as observed simultaneously by the BeppoSax MECS, we derive a lower limit on the Doppler factor of
$$\delta 2.4.$$
(4)
A $`\delta `$ 2.4 implies that the $`\gamma `$-ray luminosities mentioned above overestimates the intrinsic luminosities at least by a factor of about 33 since $`L_{obs}=\delta ^{3+\alpha }L_{intr}`$.
Although EGRET observed time variability at energies above 100 MeV, the COMPTEL experiment between 1 and 30 MeV simultaneously measures a constant flux of 3C 273. In particular, COMPTEL observes no hints of increased $`\gamma `$-ray emission in any of its energy bands during the two-week flaring period. This is a surprising result, and in some respects different to simultaneous COMPTEL/EGRET observations of other flaring blazars. For example, during the major outburst of 3C 279, the COMPTEL 10-30 MeV flux followed the flux trend as observed by EGRET at higher energies (Collmar et al. 1997a ). Also, by analysing the COMPTEL data of the first 3.5 years on the blazar PKS 0528+134, Collmar et al. (1997b) found the trend that the COMPTEL upper energy band follows the EGRET light curve, while the emission in the COMPTEL 1-3 MeV band was independent of the EGRET-observed behaviour. According to the measurements presented here, the largest $`\gamma `$-ray flux ($`>`$100 MeV) of 3C 273 is only a high-energy phenomenon, because it is (at least simultaneously) restricted to energies above $``$30 MeV. This is consistent with the hardening of the $`\gamma `$-ray spectrum during the flare. This observation suggests either an additional spectral component which becomes important at energies above 100 MeV or a time-offset between the high- and low-energy $`\gamma `$-rays is required.
A different behaviour in the MeV- and $`>`$100 MeV band had been observed in the so-called ’MeV-blazars’ GRO J0506-609 (Bloemen et al. (1995)) and PKS 0208-512 (Blom et al. (1995)), which - in contrast to the presented case - showed simultaneously strong MeV-emission compared to weak emission above 100 MeV. These sources indicated first that several emission components or mechanisms can be operating at $`\gamma `$-ray energies.
In the standard models, the $`\gamma `$-ray emission is generated within a jet, where blobs, filled with relativistic leptons, are moving at relativistic speeds along the jet axis. The $`\gamma `$-ray emission is generated by inverse-Compton interactions of these blob leptons with soft photons which either are provided by the environment (e.g. accretion disk) of the jet or are self-generated synchrotron photons. In such a picture, the observed behaviour of 3C 273 could be qualitatively explained by a change in the energy distribution of the blob leptons, by a change of the energy distribution of the soft target photons or by both. If the $`\gamma `$-ray flare is triggered by a change in the energy distribution of the blob leptons, the observation would require that only the high-energy end of the distribution, responsible for the $`\gamma `$-rays in the EGRET band, is increased in energy as well as in number density, while at lower energies the distribution has to remain constant to keep the MeV-emission unchanged. This case, if applicable, might provide hints on the particle acceleration mechanism. If a variation of the soft photon distribution is responsible for the $`\gamma `$-ray flare, a flare in a certain wavelength band, e.g. UV flare of accretion disk photons, could trigger the event.
An apparently natural explanation of this uncorrelated EGRET-COMPTEL behaviour would be if the MeV- and $`>`$100 MeV-emissions emanate from spatially different regions. However, we consider this explanation unlikely because – as mentioned above – simultaneous $`\gamma `$-ray variations have been observed by EGRET and COMPTEL in several blazars. Also during the observational period presented here, both experiments found 3C 273 in an active $`\gamma `$-ray state, which suggests a common region of photon generation.
In our opinition the most plausible scenario appears to be that the MeV and the GeV emissions are dominated by different radiation mechanisms with different flaring amplitudes, but emitted co-spatially by the same population of relativistic electrons which are also responsible at least for the high-frequency part of the synchrotron component. Such a two-component scenario for spectral variability during high-energy flares has first been suggested for PKS 0528+134 by Collmar et al. (1997b); see also Böttcher & Collmar (1998) and Mukherjee et al. (1999). Given typical parameter values for the relativistic electron distribution in AGN jets, and scaling $`\delta `$ in units of 3 (due to the results of Eq. 4), the synchrotron emission peaks at
$$ϵ_{Sy}710^8B_0\gamma _3^2\left(\frac{\delta }{3}\right),$$
(5)
where $`ϵ=h\nu /(m_ec^2)`$ is the dimensionless photon energy, $`B_0`$ is the co-moving magnetic field in Gauss, and $`\gamma _3`$ is the average Lorentz factor of the electrons in units of $`10^3`$. With $`B_0`$ and $`\gamma _3`$ being of order unity, the synchrotron peak is at $`\nu _{sy}10^{13}`$ Hz, as generally observed for 3C 273. We will discuss two possible radiation mechanisms for the observed high-energy flare: a) Comptonization of accretion disk radiation which is reprocessed by broad-line region clouds (ECC for External Comptonization of radiation from Clouds; Sikora et al. (1994)), and b) Comptonization of synchrotron radiation from the jet, reprocessed by broad-line region clouds (RSy for synchrotron reflection) as proposed by Ghisellini & Madau (1996). The ECC spectrum, which assumes the accretion disk (’blue bump’) photons to be the soft target photons, is expected to have a rather narrow spectral distribution, peaking around
$$ϵ_{ECC}310^3ϵ_D_4\mathrm{\Gamma }_{10}\gamma _3^2\left(\frac{\delta }{3}\right),$$
(6)
where $`ϵ_D_4`$ ($`ϵ_D`$/10<sup>-4</sup>) – being of order unity – is the dimensionless $`\nu F_\nu `$ peak energy of the accretion disk spectrum and $`\mathrm{\Gamma }_{10}`$ is the bulk Lorentz factor in units of 10, which is also assumed to be of order unity. Consequently, this peak is expected to be at $`1`$ GeV. In contrast, the RSy spectrum is expected to be much broader (similar to the synchrotron self-Compton spectrum) and peaks around
$$ϵ_{RSy}7B_0\gamma _3^4\mathrm{\Gamma }_{10}^2\left(\frac{\delta }{3}\right),$$
(7)
typically at MeV energies. The factors $`\mathrm{\Gamma }_{10}`$ and $`\mathrm{\Gamma }_{10}^2`$ in Eqs. (6) and (7) arise from the Lorentz boost of the external radiation field into the comoving rest frame of the blob and of the jet synchrotron radiation into the stationary frame of the BLR.
The estimates derived from Eqs. (6) and (7) indicate that a flare in the EGRET energy range without a significant variation of the MeV emission is more likely to be caused by the ECC mechanism than by the RSy scenario. We point out that these interpretations are based on time-variability of EGRET which is signifcant at 3.1$`\sigma `$.
Recently McHardy et al. (1999) reported simultaneous mm, infrared (IR), and X-ray (3-20 keV) observations of 3C 273 which cover in part the same time period as the $`\gamma `$-ray observations reported here. In particular they observed a simultaneous IR and X-ray flare lasting for about 10 days, which in fact is simultaneous to the high-energy $`\gamma `$-ray flare observed by EGRET. These simultaneous flares add important information in light of the two-component hypothesis for the $`\gamma `$-ray spectrum proposed here and provide additional constraints for the modelling of the emission processes in this quasar. A discussion and interpretation of the peculiar variability pattern of 3C 273 in IR, X-rays and $`\gamma `$-rays as observed in early 1997 will be given by Böttcher & Collmar (2000).
## 5 Conclusion
From December 10, 1996 to January 28, 1997 the CGRO instruments EGRET and COMPTEL observed the Virgo sky region continuously for 7 weeks, detecting 3C 273 in an active $`\gamma `$-ray state. EGRET ($`>`$100 MeV) observed a time-variable flux, peaking during a 2-week flaring period at its highest level observed during the CGRO-era. COMPTEL, however, does not observe any contemporaneous $`\gamma `$-ray flare at energies below $``$30 MeV, showing that this outburst is restricted to $`\gamma `$-ray energies above 30 MeV. This is consistent with the spectral hardening observed in the 3 MeV to 10 GeV energy band during the flaring period.
The peculiar variability properties of the flare may be explained in terms of a two-component spectral model with the emission in the EGRET energy range produced by Comptonization of reprocessed accretion disk emission. The different variability behaviour in $`\gamma `$-rays is inconsistent with the synchrotron-reflection model being the cause of the $`\gamma `$-ray flare.
This observation covers an opportune sequence of low pre-flare, high flare, and again low post-flare emission in $`\gamma `$-rays. In general, this $`\gamma `$-ray observation could turn out to be important for further modelling of blazar emission processes because the $`\gamma `$-ray flare is well located in time and therefore can possibly be correlated to flux measurements of monitoring observations in other wavelength regions.
###### Acknowledgements.
This research was supported by the German government through DLR grant 50 QV 9096 8, by NASA under contract NAS5-26645, and by the Netherlands Organisation for Scientific Research NWO. |
warning/0001/hep-ph0001033.html | ar5iv | text | # hep-ph/0001033 FTUV/00-7 IFIC/00-07 LC-TH-2000-005 Light Stop: MSSM versus R–parity violation
## 1 Introduction
The search for supersymmetry (SUSY) plays an important rôle in the experimental program at the colliders LEP2 and Tevatron. It will be even more important at future colliders, e.g. an upgraded Tevatron, LHC, an $`e^+e^{}`$ linear collider. Therefore many phenomenological studies have been carried out in recent years (see e.g. and references therein). Most of them have been carried out in the context of the minimal supersymmetric standard model (MSSM) . However, neither gauge invariance nor supersymmetry requires the conservation of R-parity. Indeed, there is considerable theoretical and phenomenological interest in studying possible implications of alternative scenarios in which R-parity is broken . The violation of R-parity could arise explicitly as a residual effect of some larger unified theory , or spontaneously, through nonzero vacuum expectation values (vev’s) for scalar neutrinos . In realistic spontaneous R-parity breaking models there is an $`SU(2)U(1)`$ singlet sneutrino vev characterizing the scale of R-parity violation which is expected to be the same as the effective supersymmetry breaking scale.
There are two generic cases of spontaneous R-parity breaking models to consider. In the absence of any additional gauge symmetry, these models lead to the existence of a physical massless Nambu-Goldstone boson, called majoron (J) which is the lightest SUSY particle, massless and therefore stable. If lepton number is part of the gauge symmetry and R-parity is spontaneously broken then there is an additional gauge boson which gets mass via the Higgs mechanism, and there is no physical Goldstone boson . As in the standard case in R-parity breaking models the lightest SUSY particle (LSP) is in general a neutralino. However, it now decays mostly into visible states, therefore diluting the missing momentum signal and bringing in increased multiplicity events which arise mainly from three-body decays such as
$$\stackrel{~}{\chi }_1^0f\overline{f}\nu ,$$
(1)
where $`f`$ denotes a charged fermion. The neutralino also has the invisible decay mode
$$\stackrel{~}{\chi }_1^03\nu .$$
(2)
as well as
$$\stackrel{~}{\chi }_1^0\nu J,$$
(3)
in the case the breaking of R-parity is spontaneous . This last decay conserves R-parity since the majoron has a large R-odd singlet sneutrino component.
Owing to the large top Yukawa coupling the stops have a quite different phenomenology compared to those of the first two generations of up–type squarks (see e.g. and references therein). The large Yukawa coupling implies a large mixing between $`\stackrel{~}{t}_L`$ and $`\stackrel{~}{t}_R`$ and large couplings to the higgsino components of neutralinos and charginos. The large top quark mass also implies the existence of scenarios where all two-body decay modes of $`\stackrel{~}{t}_1`$ (e.g. $`\stackrel{~}{t}_1t\stackrel{~}{\chi }_i^0,b\stackrel{~}{\chi }_j^+,t\stackrel{~}{g}`$) are kinematically forbidden. In these scenarios higher order decays of $`\stackrel{~}{t}_1`$ become relevant: : $`\stackrel{~}{t}_1c\stackrel{~}{\chi }_{1,2}^0`$, $`\stackrel{~}{t}_1W^+b\stackrel{~}{\chi }_1^0`$, $`\stackrel{~}{t}_1H^+b\stackrel{~}{\chi }_1^0`$, $`\stackrel{~}{t}_1b\stackrel{~}{l}_i^+\nu _l`$, $`\stackrel{~}{t}_1b\stackrel{~}{\nu }_ll^+`$, where $`l`$ denotes $`e,\mu ,\tau `$. In it has been shown that the three-body decay modes are in general much more important than the two body decay mode in the framework of the MSSM. Recently it has been demonstrated that not only LSP decays are sign of R-parity violation but that also the light stop is possible candidate for observing R-parity violation even if R-parity violation is small . It has been demonstrated that there exists a large parameter region where the R-parity violating decay
$$\stackrel{~}{t}_1b\tau $$
(4)
is much more important than
$$\stackrel{~}{t}_1c\stackrel{~}{\chi }_{1,2}^0$$
(5)
in scenarios where only those decay modes are possible. It is therefore natural to ask if there exist scenarios where the decay $`\stackrel{~}{t}_1b\tau `$ is as important as the three–body decays. Note that in the R-parity violating models under consideration the neutral (charged) Higgs–bosons mix with the neutral (charged) sleptons. These states are denoted by $`S_i^0`$, $`P_j^0`$, and $`S_k^\pm `$ for the neutral scalars, pseudoscalars and charged scalars, respectively. Therefore in the R-parity violating case one has the following three-body decay modes:
$`\stackrel{~}{t}_1`$ $``$ $`W^+b\stackrel{~}{\chi }_1^0`$ (6)
$`\stackrel{~}{t}_1`$ $``$ $`S_k^+b\stackrel{~}{\chi }_1^0`$ (7)
$`\stackrel{~}{t}_1`$ $``$ $`S_k^+b\nu _l`$ (8)
$`\stackrel{~}{t}_1`$ $``$ $`bS_i^0l^+,`$ (9)
$`\stackrel{~}{t}_1`$ $``$ $`bP_j^0l^+.`$ (10)
We will demonstrate that $`\stackrel{~}{t}_1b\tau ^+`$ can indeed be the most important decay mode. In particular we will consider a mass range of $`\stackrel{~}{t}_1`$, where it is difficult for the LHC to discover the light stop within the MSSM due to the large top background . The rest of paper is organized in the following way: in the next section we will introduce the model. In Sect. 3 numerical results for stop decays are presented and their implications for LC. In Sect. 4 we present our conclusions.
## 2 The model
The supersymmetric Lagrangian is specified by the superpotential $`W`$ given by
$$W=\epsilon _{ab}\left[h_U^{ij}\widehat{Q}_i^a\widehat{U}_j\widehat{H}_2^b+h_D^{ij}\widehat{Q}_i^b\widehat{D}_j\widehat{H}_1^a+h_E^{ij}\widehat{L}_i^b\widehat{R}_j\widehat{H}_1^a\mu \widehat{H}_1^a\widehat{H}_2^b\right]+\epsilon _{ab}ϵ_i\widehat{L}_i^a\widehat{H}_2^b,$$
(11)
where $`i,j=1,2,3`$ are generation indices, $`a,b=1,2`$ are $`SU(2)`$ indices, and $`\epsilon `$ is a completely antisymmetric $`2\times 2`$ matrix, with $`\epsilon _{12}=1`$. The symbol “hat” over each letter indicates a superfield, with $`\widehat{Q}_i`$, $`\widehat{L}_i`$, $`\widehat{H}_1`$, and $`\widehat{H}_2`$ being $`SU(2)`$ doublets with hypercharges $`1/3`$, $`1`$, $`1`$, and $`1`$ respectively, and $`\widehat{U}`$, $`\widehat{D}`$, and $`\widehat{R}`$ being $`SU(2)`$ singlets with hypercharges $`\frac{4}{3}`$, $`\frac{2}{3}`$, and $`2`$ respectively. The couplings $`h_U`$, $`h_D`$ and $`h_E`$ are $`3\times 3`$ Yukawa matrices, and $`\mu `$ and $`ϵ_i`$ are parameters with units of mass.
Supersymmetry breaking is parametrized by the standard set of soft supersymmetry breaking terms
$`V_{soft}`$ $`=`$ $`M_Q^{ij2}\stackrel{~}{Q}_i^a\stackrel{~}{Q}_j^a+M_U^{ij2}\stackrel{~}{U}_i^{}\stackrel{~}{U}_j+M_D^{ij2}\stackrel{~}{D}_i^{}\stackrel{~}{D}_j+M_L^{ij2}\stackrel{~}{L}_i^a\stackrel{~}{L}_j^a+M_R^{ij2}\stackrel{~}{R}_i^{}\stackrel{~}{R}_j`$ (12)
$`+m_{H_1}^2H_1^aH_1^a+m_{H_2}^2H_2^aH_2^a`$
$`[\frac{1}{2}M_3\lambda _3\lambda _3+\frac{1}{2}M\lambda _2\lambda _2+\frac{1}{2}M^{}\lambda _1\lambda _1+h.c.]`$
$`+\epsilon _{ab}[A_U^{ij}h_U^{ij}\stackrel{~}{Q}_i^a\stackrel{~}{U}_jH_2^b+A_D^{ij}h_D^{ij}\stackrel{~}{Q}_i^b\stackrel{~}{D}_jH_1^a+A_E^{ij}h_E^{ij}\stackrel{~}{L}_i^b\stackrel{~}{R}_jH_1^a`$
$`B\mu H_1^aH_2^b+B_iϵ_i\stackrel{~}{L}_i^aH_2^b],`$
Note that, in the presence of soft supersymmetry breaking terms the bilinear terms $`ϵ_i`$ can not be rotated away, since the rotation that eliminates it reintroduces an R–Parity violating trilinear term, as well as a sneutrino vacuum expectation value .
For our discussion it suffices to assume R-parity Violation (RPV) only in the third generation. However we do allow for R-parity-conserving Flavour Changing Neutral Currents (FCNC) effects, such as the process $`\stackrel{~}{t}_1c\stackrel{~}{\chi }_1^0`$ involving the three generations of quarks. In this case we will omit the labels $`i,j`$ in the soft breaking terms. In order to study the R–Parity violating decay mode $`\stackrel{~}{t}_1b\tau `$ it is sufficient to consider the superpotential
$$W=h_t\widehat{Q}_3\widehat{U}_3\widehat{H}_2+h_b\widehat{Q}_3\widehat{D}_3\widehat{H}_1+h_\tau \widehat{L}_3\widehat{R}_3\widehat{H}_1\mu \widehat{H}_1\widehat{H}_2+ϵ_3\widehat{L}_3\widehat{H}_2$$
(13)
This amounts to neglecting the effects of RPV on the two first families. A short discussion on $`\stackrel{~}{t}_1bl^+`$ in the three generation model will be given at the end of Sect. 3.
The bilinear term in Eq. (13) leads to a mixing between the charginos and the $`\tau `$–lepton which in turn leads to the decay $`\stackrel{~}{t}_1b\tau `$. The mass matrix is given by
$$𝐌_𝐂=\left[\begin{array}{ccc}M& \frac{1}{\sqrt{2}}gv_2& 0\\ \frac{1}{\sqrt{2}}gv_d& \mu & \frac{1}{\sqrt{2}}h_\tau v_3\\ \frac{1}{\sqrt{2}}gv_3& ϵ_3& \frac{1}{\sqrt{2}}h_\tau v_d\end{array}\right]$$
(14)
As in the MSSM, the chargino mass matrix is diagonalized by two rotation matrices $`𝐔`$ and $`𝐕`$
$$𝐔^{}𝐌_𝐂𝐕^1=\left[\begin{array}{ccc}m_{\stackrel{~}{\chi }_1^\pm }& 0& 0\\ 0& m_{\stackrel{~}{\chi }_2^\pm }& 0\\ 0& 0& m_\tau \end{array}\right].$$
(15)
The lightest eigenstate of this mass matrix must be the tau lepton ($`\tau ^\pm `$) and so the mass is constrained to be 1.7771 GeV. To obtain this the tau Yukawa coupling becomes a function of the parameters in the mass matrix, and the full expression is given in .
The stop mass matrix is given by
$$M_{\stackrel{~}{t}}^2=\left[\begin{array}{cc}M_Q^2+\frac{1}{2}v_u^2h_{t}^{}{}_{}{}^{2}+\mathrm{\Delta }_{UL}& \frac{h_U}{\sqrt{2}}\left(v_uA_t\mu v_d+ϵ_3v_3\right)\\ \frac{h_U}{\sqrt{2}}\left(v_uA_t\mu v_d+ϵ_3v_3\right)& M_U^2+\frac{1}{2}v_u^2h_{t}^{}{}_{}{}^{2}+\mathrm{\Delta }_{UR}\end{array}\right]$$
(16)
with $`\mathrm{\Delta }_{UL}=\frac{1}{8}\left(g^2\frac{1}{3}g_{}^{}{}_{}{}^{2}\right)\left(v_d^2v_u^2+v_3^2\right)`$ and $`\mathrm{\Delta }_{UR}=\frac{1}{6}g_{}^{}{}_{}{}^{2}(v_d^2v_u^2+v_3^2)`$. The sum of the $`v_i^2`$ is given by $`m_W^2=g^2(v_d^2+v_u^2+v_3^2)/2`$. The mass eigenstates are given by $`\stackrel{~}{t}_1=\stackrel{~}{t}_L\mathrm{cos}\theta _{\stackrel{~}{t}}+\stackrel{~}{t}_R\mathrm{sin}\theta _{\stackrel{~}{t}}`$ and $`\stackrel{~}{t}_2=\stackrel{~}{t}_R\mathrm{cos}\theta _{\stackrel{~}{t}}\stackrel{~}{t}_L\mathrm{sin}\theta _{\stackrel{~}{t}}.`$ The sfermion mixing angle is given by
$$\mathrm{cos}\theta _{\stackrel{~}{t}}=\frac{M_{\stackrel{~}{t}_{12}}^2}{\sqrt{(M_{\stackrel{~}{t}_{11}}^2m_{\stackrel{~}{t}_1}^2)^2+(M_{\stackrel{~}{t}_{12}}^2)^2}},\mathrm{sin}\theta _{\stackrel{~}{t}}=\frac{M_{\stackrel{~}{t}_{11}}^2m_{\stackrel{~}{t}_1}^2}{\sqrt{(M_{\stackrel{~}{t}_{11}}^2m_{\stackrel{~}{t}_1}^2)^2+(M_{\stackrel{~}{t}_{12}}^2)^2}}.$$
(17)
In addition the charged Higgs bosons mix with charged sleptons and the real (imaginary) parts of the sneutrino mix the scalar (pseudoscalar) Higgs bosons. The formulas can be found e.g. in . Their main decay modes for the mass range considered in this study are:
$`S_i^0`$ $``$ $`b\overline{b},\tau ^+\tau ^{},\stackrel{~}{\chi }_1^0\nu _\tau `$ (18)
$`P_j^0`$ $``$ $`b\overline{b},\tau ^+\tau ^{},\stackrel{~}{\chi }_1^0\nu _\tau `$ (19)
$`S_k^{}`$ $``$ $`s\overline{c},\tau ^{}\nu _\tau ,\stackrel{~}{\chi }_1^0\tau ^{}`$ (20)
## 3 Numerical results
In this section we present our numerical results for the branching ratios of the higher order decays of $`\stackrel{~}{t}_1`$. Here we consider scenarios where all two-body decays induced at tree-level are kinematically forbidden. Before going into detail it is useful to have some approximate formulas at hand :
$`\mathrm{\Gamma }(\stackrel{~}{t}_1b\tau )`$ $``$ $`{\displaystyle \frac{g^2|U_{32}|^2h_b^2\mathrm{cos}_{\theta _{\stackrel{~}{t}}}^2m_{\stackrel{~}{t}_1}}{16\pi }}`$ (21)
$`\mathrm{\Gamma }(\stackrel{~}{t}_1c\stackrel{~}{\chi }_1^0)`$ $``$ $`10^6h_b^4m_{\stackrel{~}{t}_1}`$ (22)
where $`|U_{32}||ϵ_3/\mu |`$ if $`|ϵ_3||\mu |`$ and $`v_3m_W`$. The complete formulas are given in . For the three–body decays the formulas given in can be used as a good approximation if the mixings induced by R-parity violation are small. The complete formulas for the three–body decays in the R–parity violating case will be given elsewhere .
We have fixed the parameters as in to avoid colour breaking minima: we have used $`m_{\stackrel{~}{t}_1}`$, $`\mathrm{cos}\theta _{\stackrel{~}{t}}`$, $`\mathrm{tan}\beta `$, and $`\mu `$ as input parameters in the top squark sector. For the sbottom (stau) sector we have fixed $`M_{\stackrel{~}{Q}},M_{\stackrel{~}{D}}`$ and $`A_b`$ ($`M_{\stackrel{~}{E}},M_{\stackrel{~}{L}}`$, and $`A_\tau `$) as input parameters. In addition we choose the R-parity violating parameters $`ϵ_3`$ and $`v_3`$ in such a way that the tau neutrino mass is fixed ( and references therein):
$$m_{\nu _\tau }\frac{(g^2M^{}+g^2M)\mu _{}^{}{}_{}{}^{2}}{4MM^{}\mu _{}^{}{}_{}{}^{2}2(g^2M^{}+g^2M)\mu ^{}v_uv_d^{}\mathrm{cos}\xi }v_{d}^{}{}_{}{}^{2}\mathrm{sin}^2\xi $$
(23)
with
$`\mathrm{sin}\xi `$ $`=`$ $`{\displaystyle \frac{ϵ_3v_d+\mu v_3}{\sqrt{\mu ^2+ϵ_3^2}\sqrt{v_d^2+v_3^2}}}`$ (24)
$`\mu ^{}`$ $`=`$ $`\sqrt{\mu ^2+ϵ_3^2},v_d^{}=\sqrt{v_d^2+v_3^2}.`$ (25)
For simplicity, we assume that the soft SUSY breaking parameters are equal for all generations.
In Fig. 1(a) and (b) we show the branching ratios of $`\stackrel{~}{t}_1`$ as a function of $`\mathrm{cos}\theta _{\stackrel{~}{t}}`$. The parameters and physical quantities are given in Tab. 1. In Fig. 1(a) we show $`\text{BR}(\stackrel{~}{t}_1b\tau ^+)`$, $`\text{BR}(\stackrel{~}{t}_1c\stackrel{~}{\chi }_1^0)`$, $`\text{BR}(\stackrel{~}{t}_1bW^+\stackrel{~}{\chi }_1^0)`$, $`\text{BR}(\stackrel{~}{t}_1be^+\stackrel{~}{\nu }_e`$) + $`\text{BR}(\stackrel{~}{t}_1b\nu _e\stackrel{~}{e}_L^+)`$. The branching ratios for decays into $`\stackrel{~}{\mu }_L`$ or $`\stackrel{~}{\nu }_\mu `$ are practically the same as those into $`\stackrel{~}{e}_L`$ or $`\stackrel{~}{\nu }_e`$. We have summed up those branching ratios for the decays into sleptons that give the same final state, for example:
$`\stackrel{~}{t}_1b\nu _e\stackrel{~}{e}_L^+be^+\nu _e\stackrel{~}{\chi }_1^0,\stackrel{~}{t}_1be^+\stackrel{~}{\nu }_ebe^+\nu _e\stackrel{~}{\chi }_1^0`$ (26)
Note that in Fig. 1 we have also summed the decay branching ratios $`\text{BR}(\stackrel{~}{t}_1bS_k^+\nu _\tau )`$ \+ $`\text{BR}(\stackrel{~}{t}_1b\tau ^+S_i^0)`$ \+ $`\text{BR}(\stackrel{~}{t}_1b\tau ^+P_j^0)`$.
In the above cases the assumption $`m_{\stackrel{~}{t}_1}m_b<m_{\stackrel{~}{\chi }_1^+}`$ implies $`m_{\stackrel{~}{\chi }_1^+}>m_{\stackrel{~}{l}}`$. Therefore, charginos can not arise as decay products of sleptons. The latter can only decay into the corresponding lepton plus $`\stackrel{~}{\chi }_1^0`$ except for a small parameter region where the decay into $`\stackrel{~}{\chi }_2^0`$ is possible. However, this decay is negligible due to kinematics in that region. In addition there exists the possibility of R-parity violating decays. However, these will be small because the neutrinos mix mainly with higgsinos implying that the partial decay widths are proportional to the squared product of an R-parity violating mixing parameter and small Yukawa coupling. For this set of parameters $`\text{BR}(\stackrel{~}{t}_1c\stackrel{~}{\chi }_1^0)`$ is $`O(10^4)`$ independent of $`\mathrm{cos}\theta _{\stackrel{~}{t}}`$ and therefore negligible. Near $`\mathrm{cos}\theta _{\stackrel{~}{t}}=0.3`$ one has $`\stackrel{~}{t}_1bW^+\stackrel{~}{\chi }_1^0`$ as dominant decay channel, since the $`\stackrel{~}{t}_1`$-$`\stackrel{~}{\chi }_1^+`$-$`b`$ coupling vanishes implying that the main contribution for the decays into the scalars vanishes. Moreover, the width for $`\stackrel{~}{t}_1b\tau ^+`$ is somewhat suppressed because of the $`\mathrm{cos}^2\theta _{\stackrel{~}{t}}`$ factor in Eq. (21). Note, from the figure that the branching ratios for the various decays into selectrons $`e`$-sneutrino is roughly a factor two smaller than the sum of the decays into $`S_k^\pm `$, $`S_i^0`$, $`P_j^0`$. The reason is that, for this choice of parameters the $`P_2^0`$ is mainly the pseudoscalar Higgs boson $`A^0`$ of the MSSM with mass 110 GeV. In this case the R-parity violating channel $`\stackrel{~}{t}_1b\tau ^+P_2^0`$ is comparable to the corresponding R-parity conserving decays. This state appears additional to the states which carry tau–lepton number in the MSSM limit giving rise to the observed difference. Note, that one has to expect additional jets from the states containing the scalars $`S_i^0`$, $`P_j^0`$, and $`S_k^\pm `$ because they have admixtures of the original Higgs boson. In case of negative $`\mathrm{cos}\theta _{\stackrel{~}{t}}`$ the decay into $`\stackrel{~}{t}_1b\tau ^+`$ is important and can be even the most important one. Therefore one has events with $`\tau ^+\tau ^{}b\overline{b}`$ in the final state which can be used for a full mass reconstruction of the light stop.
In Fig. 1(b) the $`\mathrm{tan}\beta `$ dependence of branching ratios is shown. For this specific choice of $`\mathrm{cos}\theta _{\stackrel{~}{t}}`$ the decay $`\stackrel{~}{t}_1b\tau ^+`$ is the most important one for $`\mathrm{tan}\beta <15`$. Above this value the final states which contain the scalars corresponding to the lighter MSSM stau are the most important ones. The growth of the decay branching ratios into these states with $`\mathrm{tan}\beta `$ is a feature independent of $`\mathrm{cos}\theta _{\stackrel{~}{t}}`$.
The assumption that no tree-level-induced two-body decays are kinematically allowed implies that $`m_{\stackrel{~}{\chi }_1^+}>m_{\stackrel{~}{t}_1}m_b`$. Therefore, one expects an increase of $`\text{BR}(\stackrel{~}{t}_1bW^+\stackrel{~}{\chi }_1^0)`$ if $`m_{\stackrel{~}{t}_1}`$ increases, because the decay into $`bW^+\stackrel{~}{\chi }_1^0`$ is dominated by the $`t`$ exchange, whereas for the decays into scalars $`\stackrel{~}{\chi }_j^+`$ exchange dominates. This trend is indeed observed in Fig. 2, where we show the branching ratios for $`m_{\stackrel{~}{t}_1}=350`$ GeV. Here we concentrate on the range of $`\mathrm{cos}\theta _{\stackrel{~}{t}}`$ where $`A_t1`$ TeV to avoid possible minima in the scalar potential which break either color or electric charge. Notice that for the heavy stop case the decay $`\stackrel{~}{t}_1bW^+\stackrel{~}{\chi }_1^0`$ is the most important one, independently of $`\mathrm{cos}\theta _{\stackrel{~}{t}}`$ and $`\mathrm{tan}\beta `$.
Note however that also in this case R-parity violation implies a distinct signature compared to what is expected in the MSSM due to the decays of $`\stackrel{~}{\chi }_1^0`$. One gets the following high-multiplicity final states:
$`\stackrel{~}{t}_1`$ $``$ $`bW^+f\overline{f}\nu `$ (27)
$``$ $`bW^+f\overline{f}^{}l^\pm `$ (28)
$``$ $`bW^+\nu J`$ (29)
where $`f`$ denotes a Standard Model fermion. Here the decays of the $`W`$–boson will give additional leptons and jets. Therefore, one has in general additional jets and leptons compared to the MSSM case.
In the event that $`ϵ_{1,2}`$ are of the same order of magnitude, as suggested by a solution to the present neutrino anomalies one has in addition the decays into $`be^+`$ and $`b\mu ^+`$. If one passes from the 1-generation model to the 3-generation model the situation changes as follows. From Eq. (21) it follows that the sum of the modes $`\mathrm{\Gamma }(\stackrel{~}{t}_1bl^+)`$ in the 3–generation model is nearly equal to $`\mathrm{\Gamma }(\stackrel{~}{t}_1b\tau ^+)`$ in the 1–generation model, if $`(ϵ_1^{})^2+(ϵ_2^{})^2+(ϵ_2^{})^2=ϵ_3^2`$ where the $`ϵ_i^{}`$ ($`ϵ_3`$) are the parameters of the 3–generation (1–generation) model. In Fig. 3 we show the ratios of branching ratios for different $`\stackrel{~}{t}_1bl^+`$ modes versus the ratios of different $`ϵ`$’s squared and for different values of $`\mathrm{cos}\theta _{\stackrel{~}{t}}`$. In both cases we have fixed $`ϵ_1^2+ϵ_2^2+ϵ_3^2=1`$ GeV<sup>2</sup> and in Fig. 3a $`ϵ_1=ϵ_3`$ whereas in Fig. 3b $`ϵ_2=ϵ_3`$. One can see that the dependence is nearly linear even for rather small $`\mathrm{cos}\theta _{\stackrel{~}{t}}`$. This result depends on $`(ϵ_1^2+ϵ_2^2+ϵ_3^2)/\mu ^2`$ and on the neutrino mass $`m_{\nu _\tau }`$, since both determine the mixings of the leptons with the charginos. The lines indicated in the figure come closer to the diagonal if $`(ϵ_1^2+ϵ_2^2+ϵ_3^2)/\mu ^2`$ increases and $`m_{\nu _\tau }`$ decreases.
## 4 Conclusions
We have studied the phenomenology of the lightest stop in scenarios where the R-parity violating decay $`\stackrel{~}{t}_1b\tau ^+`$ competes with three–body decays. We have found that for $`m_{\stackrel{~}{t}_1}<250`$ GeV there are regions of parameter where $`\stackrel{~}{t}_1b\tau ^+`$ is an important decay mode if not the most important one. This implies that there exists the possibility for full stop mass reconstruction from $`\tau ^+\tau ^{}b\overline{b}`$ final states. Moreover, in this mass range the discovery of the lightest stop might not be possible at the LHC (certainly this is the case in the MSSM). This implies that one has to take into account the importance of this new decay mode when designing the stop search strategies at a future $`e^+e^{}`$ Linear Collider. Spontaneously and bilinearly broken R-parity violation imply additional leptons and/or jets in stop cascade decays. Looking at the three generation model the decays into $`\stackrel{~}{t}_1bl^+`$ imply the possibility of measuring $`ϵ_e^2/ϵ_\tau ^2`$ and $`ϵ_\mu ^2/ϵ_\tau ^2`$ and thereby probing the parameters associated with the present solar and atmospheric neutrino anomalies.
## Acknowledgments
This work was supported by DGICYT under grants PB95-1077 and by the TMR network grant ERBFMRXCT960090 of the European Union. D.R. was supported by Colombian COLCIENCIAS fellowship. W.P. was supported by the Spanish ’Ministerio de Educacion y Cultura’ under the contract SB97-BU0475382. |
warning/0001/cond-mat0001388.html | ar5iv | text | # Universality Class of Thermally Diluted Ising Systems at Criticality
## Abstract
The universality class of thermally diluted Ising systems, in which the realization of the disposition of magnetic atoms and vacancies is taken from the local distribution of spins in the pure original Ising model at criticality, is investigated by finite size scaling techniques using the Monte Carlo method. We find that the critical temperature, the critical exponents and therefore the universality class of these thermally diluted Ising systems depart markedly from the ones of short range correlated disordered systems. Our results agree fairly well with theoretical predictions previously made by Weinrib and Halperin for systems with long range correlated disorder.
I. Introduction
During last decades the systems with quenched randomness have been intensively studied . The Harris criterion predicts that weak dilution does not change the character of the critical behavior near second order phase transitions for systems of dimension $`d`$ with specific heat exponent lower than zero in the pure case (the so called P systems), $`\alpha _{pure}<0\nu _{pure}>2/d`$, being $`\nu `$ the correlation length critical exponent. This criterion has been supported by renormalization group (RG) , and scaling analyses . The effect of strong dilution was studied by Chayes et al. . For $`\alpha _{pure}>0`$ (the so called R systems), the Ising 3D case for example, the system fixed point flows from a pure (undiluted) fixed point towards a new stable fixed point at which $`\alpha _{random}<0`$. Recently Ballesteros et al. have used the Monte Carlo approach to study the diluted Ising systems in two, three and four dimensions . The existence of a new universality class for the random diluted Ising system (RDIS), different from that of the pure Ising model and independent of the average density of occupied spin states $`(p)`$, is proved, using an infinite volume extrapolation technique based upon the leading correction to scaling. The critical exponents obtained this way could be compared with the experimental critical exponents for a random disposition of vacancies in diluted magnetic systems .
In all cases previously mentioned frozen disorder was always produced in a random way, that is, vacancies were distributed throughout the lattice randomly. Real systems, however, can be realized with other kinds of disorder, where the vacancy locations are correlated. In particular, long range correlation (LRC) has been found in X-ray and Neutron Critical Scattering experiments in systems undergoing magnetic and structural phase transitions. This effect has been modeled by assuming a spatial distribution of critical temperatures obeying a power law $`g(x)x^a`$ for large separations $`x`$ . In general these systems behave in a way very different from RDIS, since systems with randomly distributed impurities may be considered as the limit case of short range correlated (SRC) distributions of the vacancies. The basic approach to the critical phenomena of LRC systems was established by Weinrib and Halperin almost two decades ago. They found that the Harris criterion can be extended for these cases, showing that for $`a<d`$, the disorder is irrelevant if $`a\nu _{pure}2>0`$, and that in the case of relevant disorder a new universality class (and a new fixed point) with correlation length exponent $`\nu =2/a`$, and a specific heat exponent $`\alpha =2(ad)/a`$ appears. In contrast, if $`a>d`$, the usual Harris criterion for SRC systems is recovered. LRC disorder has been studied also by the Monte Carlo approach . In this case a correlation function $`g(x)=x^a`$ with $`a=2`$ (defects consisting in randomly oriented lines of magnetic vacancies inside a three dimensional Ising system) confirmed the theoretical predictions of Weinrib and Halperin.
In the present paper we will study Ising three dimensional systems where the long range correlated dilution has been introduced as a thermal order-disorder distribution of vacancies in equilibrium, governed by a characteristic ordering temperature ($`\theta `$), in a way similar to the thermal disordering in a binary alloy of magnetic (spins) and non-magnetic atoms (vacancies). \[See previous work and also the application of this kind of disorder in percolation problems \]. Making more explicit the analogy with order-disorder in alloys we may distinguish clearly between:
i) thermally diluted Ising system realizations (TDIS), in which the quenched randomness is produced by considering a ferromagnetic Ising system at $`(\theta =T_c^{3D})`$: after thermalization, the spins of the dominant type (concentration $`c0.5`$) are taken as the locations of the magnetic atoms, and the rest are taken as the magnetic vacancies. The structure of the realization is fixed thereafter for all temperatures at which the magnetic interactions are subsequently investigated.
ii) random diluted Ising system realizations (RDIS) (equivalent to $`\theta >>T_c^{3D}`$), also fixed thereafter for all temperatures at which the spin interactions are investigated. Of course, we chose for comparison a vacancy probability $`p=0.5`$, resulting in $`c0.5.`$
iii) as a third possibility, not considered here, one could investigate what might be called an antiferromagnetically diluted Ising system (with $`\theta <<T_N^{3D}`$, which would lead to a disposition of non-magnetic atoms (vacancies) strictly alternating with spins, with $`c=0.5`$.
So, if this ordering temperature $`\theta `$ determining the particular realization is high enough, the equilibrium thermal disorder will be very similar to the random (short range correlated) disorder of the usual previous investigations. On the other hand, if $`\theta `$ happens to coincide with the characteristic magnetic critical temperature ($`T_c^{3D}=4.511617`$) of the undiluted system, we will have vacancies in randomly located points, but with a long range correlated distribution. \[Note that the situation differs markedly from that of previously studied LRC systems, in which lines or planes of vacancies were considered\]. The correlation of our TDIS is given by a value $`a=2\eta _{pure},`$ where $`\eta _{pure}`$ is the correlation function exponent for the pure system. Since $`d=3`$ and $`\eta _{pure}=0.03`$ for the three dimensional Ising system, we have a long range correlated disorder with $`a=1.97<3=d`$. So we are in the case where LRC disorder is relevant and we should detect a change of universality class with respect to the SRC case (following Weinrib and Halperin we expect for the thermally diluted Ising system $`\nu 1`$ and $`\alpha 1)`$. Details about the construction of these thermally diluted Ising systems (TDIS) can be found in Ref.. In the present work we study the critical behavior and the university class of three dimensional TDIS at criticality using the Monte Carlo approach. We will compare our results with the critical behavior of the RDIS.
The structure of the paper is as follows: In Section II we study the dependence of the critical temperature (and of the self-averaging at criticality) with the size of the system for both TDIS and RDIS. Once the critical point is determined, we investigate whether or not TDIS and RDIS belong to the same universality class. In order to proceed, we study the critical behavior of both kinds of systems by applying finite size scaling techniques (Section III) and by using the effective-exponent approach (Section IV). A summary of the main results and concluding remarks are given in Section V.
II. Transition temperature and self-averaging of thermally diluted systems
For a hypercubic sample of linear dimension $`L`$ and number of sites $`N=L^d`$, any observable singular property $`X`$ has different values for the different random realizations of the disorder, corresponding to the same dilution probability $`p`$ (grand-canonical constraint). This means that X behaves as a stochastic variable with average $`\overline{X}`$ (in the following, the bar indicates average over subsequent realizations of the dilution and the brackets indicate MC average). The variance would then be $`(\mathrm{\Delta }X)^2`$, and the normalized square width, correspondingly:
$$R_X=(\mathrm{\Delta }X)^2/\overline{X}^2$$
(1)
A system is said to exhibit self-averaging (SA) if $`R_X0`$ as $`L\mathrm{}`$. If the system is away from criticality, $`L>>\xi `$ (being $`\xi `$ the correlation length). The central limit theorem indicates that strong SA must be expected in this case. However, the behavior of a ferromagnet at criticality (with $`\xi >>L)`$ is not so obvious. This point has been studied recently for short range correlated quenched disorder. Aharony and Harris (AH), using a renormalization group analysis in $`d=4\epsilon `$ dimensions, proved the expectation of a rigorous absence of self-averaging in critically random ferromagnets . More recently, Monte Carlo simulations were used to investigate the self-averaging in critically disordered magnetic systems . The absence of self-averaging was confirmed. The normalized square width $`R_X`$ is an universal quantity affected just by correction to scaling terms. LRC diluted systems are expected to have different critical exponents and different normalized square widths with respect to those of the usual randomly disordered systems studied previously.
We perform Monte Carlo calculations of the magnetization and the susceptibility $`(\chi =(M^2M^2)/T)`$ per spin at different temperatures for different realizations of TDIS, and for randomly diluted systems with $`p=0.5`$ (restricting to $`c>0.5`$) using in both cases the Wolff single cluster algorithm with periodic boundary conditions, on lattices of different sizes $`L=10,20,40,60,80,100`$. Results for susceptibility vs. temperature are shown in Fig. 1.
Note how due to the existence of randomness the susceptibility points do not collapse into a single curve, since each realization has a different value of the critical temperature and of the concentration. This is even more clear for small values of L and for the critically thermal case (white points). Fig. 1 indicates that the critical temperature of the TDIS clearly differs from that of the RDIS, and also than the effect of the dilution on the lack of self-averaging seems to be stronger in the thermally diluted Ising case. The critical temperature may be obtained at the point where the normalized square width for the susceptibility, $`R_\chi `$, reaches its maximum. A different value of the critical temperature does not imply, of course, a different universality class. However from the dependence of the critical temperature with the length of the system we expect to detect a change in universality between TDIS and RDIS following the scaling relation:
$$T_c(L)=T_c(\mathrm{})+AL^{1/\nu }$$
(2)
being $`\nu `$ the critical exponent associated with the specific system’s correlation length. This critical exponent has been determined by means of Monte Carlo data for the random case by Ballesteros et al. . They found a value $`\nu _{random}=0.683`$. On the other hand the result by Weinrib and Halperin indicates that the critical exponent expected for the thermal case should be $`\nu _{thermal}=2/a=1.015`$. Fig. 2 represents the dependence of the critical temperature with respect to the length of the systems for random and thermal dilutions.
In both cases a fit to Eq. 2 has been performed for both values $`\nu =\nu _{random}=0.683`$ (continuous line) and $`\nu =\nu _{thermal}=1.015`$ (dashed line). Note how the thermal data fit better Eq. 2 for $`T_c(L)`$ using $`\nu _{thermal}`$, indicating a possible change in universality class with respect to the random case. \[If we fit the data leaving all parameters free, we find $`\nu _{random}\mathrm{.0.7}`$ and $`\nu _{thermal}1.2`$, which are very near the expected results\]. The extrapolated values of critical temperatures for infinite systems obtained this way are $`T_c^{random}(\mathrm{})=1.845\pm 0.003`$ (close to the values previously obtained by Ballesteros et al.) and $`T_c^{thermal}(\mathrm{})=3.269\pm 0.002`$ (clearly different from $`T_c(\mathrm{})`$ for the SRC case). Incidentally $`\nu _{thermal}`$ can be compared with $`\nu `$ for the observed sharp component in neutron scattering line shapes, which is around 1.3 for Tb . This point deserves more careful analysis and will be taken up elsewhere.
Once the critical temperatures are known we can perform simulations for the magnetization and the susceptibility at criticality for several realizations of thermal and random diluted systems in an effort to determine the value of $`R_\chi `$, the normalized square width for the susceptibility. We went up to 500 realizations for $`L=10,20,40`$ and up to 200 realizations for $`L=60,80`$. Results are shown in Fig.3.
The arrow represents the (concentration independent) $`R_\chi `$ value obtained by Ballesteros et al . The straight continuous line represents the average value obtained for random dilution data, and the straight dashed line gives the average value obtained for thermal dilution data. Note that the TDIS presents a lack of self-averaging around one order of magnitude larger than the RDIS. We have already presented a similar analysis for both kinds of dilution , but only at the critical temperature characteristic of the random system. Our results are not precise enough to specify accurately the evolution of the normalized square width as a function of L governed by corrections to scaling terms. However the average we obtain for $`R_\chi ^{random}=0.155`$ is close to the value previously reported by means of infinite volume extrapolations. For the thermal case we obtain $`R_\chi ^{thermal}=1.19`$, about one order of magnitude larger than for random dilution. In this case an evolution of $`R_\chi ^{thermal}`$ vs. $`L`$ given by correction to scaling terms may be also expected, but according to Weinrib and Halperin the analysis would be even more complicated, due to the fact that the long range correlated disordered systems present complex oscillating corrections to scaling.
III. Critical behavior and exponents of thermally diluted systems
The dispersion in concentration and magnetization at criticality between the different realizations is shown at a glance in scattered plots as in Ref. . Each point in Fig. 4 represents a single realization with magnetization at criticality ($`M`$) and concentration ($`c`$). Note that in both cases (TDIS and RDIS), the dispersion on the magnetization and on the concentration decreases with L, but this is more clearly shown in the thermal case. Fig. 4 shows clearly the difference in behavior between the random and the thermal cases, at least up to the values of $`L`$ considered.
From Fig. 4 we can extract averaged values for the magnetization and the inverse susceptibility, $`\overline{M}`$ and $`\overline{\chi ^1}`$ (with $`\chi =M^2`$ at the critical point) for the TDIS. Both averaged values are expected to fit the following scaling laws at criticality:
$$\overline{M(L)}L^{\beta /\nu }$$
(3)
$$\overline{\chi ^1(L)}L^{\gamma /\nu }$$
(4)
Considering $`(1/\nu )_{thermal}=a/2=0.985`$, we can obtain from our data the values of $`\beta _{thermal}`$ and $`\gamma _{thermal}`$ (see f.i. the fitting for the inverse of the susceptibility in Fig. 5). We get $`\beta _{thermal}=0.56\pm 0.05`$ and $`\gamma _{thermal}=1.91\pm 0.06`$, very close to the predicted values by Weinrib and Halperin . Using the scaling relation:
$$\alpha =22\beta \gamma $$
(5)
we obtain the following specific heat critical exponent: $`\alpha _{thermal}=1\pm 0.1`$. Weinrib and Halperin give for LRC systems $`\alpha =1`$, in good agreement with our result. An analogous analysis has been performed but using the dispersion in magnetization and inverse susceptibility instead of the averaged values. The results are similar. The same study has been made for the random case (also shown in Fig.5).
The critical exponents obtained are in agreement with those of Ballesteros et al. within our error bars.
The average concentration for the thermal system is also expected to show a scaling law behavior given by:
$$(\overline{c(L)}0.5)L^{(\beta /\nu )_{3D}}$$
(6)
where $`(\beta /\nu )_{3D}0.52`$ gives the values corresponding to the pure three dimensional Ising case, because in critically thermally diluted Ising systems, vacancies are distributed with the same long range correlation spin distribution function as in the pure case. The fit to the average value of the concentrations shown in Fig. 4 give, for the thermal case, a value $`(\beta /\nu )_{3D}0.55\pm 0.08`$. This implies also a clear difference between RDIS and TDIS, since Ising systems with vacancies randomly distributed are not expected to follow an scaling behavior with $`(\beta /\nu )_{3D}`$. \[Fitting to an scaling law the results for RDIS gives an exponent around $`1.4`$, which implies a much faster convergence to $`c=0.5`$\].
IV. Effective exponents and possible crossover for finite systems
The difference between the universality class of RDIS and TDIS can be detected also by means of the effective critical exponents. In the case of the magnetization the effective critical exponent is defined by:
$$\beta _{eff}\mathrm{log}(M)/\mathrm{log}(t)(t=T_c^iT)$$
(7)
with $`T_c^i`$ the critical temperature of the particular realization (i) characterized by a maximum of the susceptibility ($`\chi =M^2M^2`$). For $`L\mathrm{}`$ and $`t0`$, $`\beta _{eff}=\beta `$. Finite size effects force the effective critical exponent to drop to zero before the critical value is attained. However, since $`\beta _{random}=0.3546`$ (Ref. ), and $`\beta _{thermal}`$ is expected to be around $`0.5`$ (Ref. ) , effective critical exponents (for the thermal case) may rise to values greater than 0.3546 and lower than 0.5, before finite size effects appear, indicating the difference of universality class between both kinds of systems. In the RDIS the effective critical exponent is expected to be always lower than 0.3546. Monte Carlo simulations of magnetization vs. temperature have been performed for randomly $`(p=0.5)`$ and thermally diluted systems, with $`L=80`$. In Fig. 6 we show the results for $`\beta _{eff}`$ vs. $`\mathrm{log}(t)`$ for two samples of the TDIS and the RDIS type respectively. The random effective critical exponent is always under 0.35 and it seems to tend towards this value for large enough $`L`$, as expected, but for thermal systems the behavior is completely different: In the figure the value of the critical thermal effective exponent is between 0.35 and 0.5. An analogous investigation has been done for different values of L and different realizations. The same effect has been found for lower L values. However, we may note that in TDIS the effective critical exponent arrives at the maximum in a very different way depending on the realization. The reason is twofold: (1) the different disposition of the vacancies in each particular realization, and (2) the large differences in concentration for the thermal case (the rise of $`\beta _{eff}`$ towards the expected diluted universality value should be faster for c closer to 0.5). It may be noted that as the size $`(L)`$ of the sample increases the possibility of local inhomogeneities in the TDIS realizations increases, giving rise to such phenomena as pseudo double peaks in the susceptibility, reduced values for the overall critical exponents, etc. However, it is important to remark than in all realizations investigated the $`\beta _{eff}`$ values of the TDIS, (before finite size effects take over), have been clearly superior to the $`\beta `$ for the RDIS.
In Fig. 6 the effective critical exponent for the thermal case seems to produce a crossover towards $`\beta =0.35`$ as the temperature approaches the critical point, before finite size effects finally appear. In principle this might be influenced by the indetermination in the measurement of the particular critical temperature of the realization. However, this crossover may point out that different length scales might be important at different distance from the critical point. In principle characteristic lengths such as the size of the fixed vacancies and spin clusters, the size of the system itself and the thermal spin fluctuations might all play a role. The possibility that for $`L\mathrm{}`$ and $`c=0.5`$ the apparent crossover could disappear all together should not be excluded, entailing no crossover from critical thermal to random universality class.
V. Concluding remarks
To summarize, a new way to produce diluted Ising systems with a long range correlated distribution of vacancies has been analyzed by means of Monte Carlo calculations and finite size scaling techniques. We find an universality class different from that reported for diluted Ising systems with short range correlated disorder. Our systems may be included in the universality class predicted by Weinrib and Halperin for $`a2`$. This kind of thermal disorder had been already applied in percolation problems, but as far as we know it had never been applied to magnetic systems, in which long range correlated disorder has been previously produced mostly by random distribution of lines or planes of vacancies. The present dilution procedure may be applicable to systems where the long range correlated disorder is not due to dislocations, preferential lines or planes of vacancies, but to systems where the vacancies (points) are critically distributed in clusters as in the case of order-disorder systems.
We thank H.E. Stanley for encourament and we are grateful to H.G. Ballesteros, L.A. Fernández, V. Martín-Mayor, and A. Muñoz Sudupe for helpful correspondence. We thank P.A. Serena for helpful comments and for generous access to his computing facilities. We acknowledge financial support from DGCyT through grant PB96- 0037 and from the Basque Regional Government (J.I.). |
warning/0001/gr-qc0001022.html | ar5iv | text | # Atom made from charged elementary black hole
## I Introduction
Fluctuations in the density of the very early universe led to parts with extremely high density to collapse and form black holes . The density of such objects would be greatly reduced in models of the early universe that include an inflation, but they do not have to disappear entirely. These particles would help explain the mass deficit of the universe.
A black hole is essentially an elementary particle in the sense that it is completely described by a set of quantum numbers and can have no detectable internal structure (the famous “black holes have no hair” theory, see e.g. ). Recently it has been demonstrated that the black hole mass is quantised in units of the Planck mass $`M_p=(\mathrm{}c/G)^{1/2}=1.2\times 10^{19}`$ GeV $`=2\times 10^5`$ g (see also ). One intriguing reason for quantisation of black holes comes from a classical formula. The horizon area of Kerr-Newman black holes is
$$A=4\pi \left[\left(M+\sqrt{M^2a^2Q^2}\right)^2+a^2\right]$$
(1)
where $`a`$ is the angular momentum divided by the mass $`M`$, and $`Q`$ is the charge of the black hole. This implies that for a given $`a`$ and $`Q`$ there is a minimum mass in order to avoid the singularity of the radical. In ordinary units (the previous equation was written in gravitational units where $`\mathrm{}=c=G=1`$) one obtains
$$M_{min}=M_p\left(Z^2\alpha /2+\sqrt{Z^4\alpha ^2/4+J(J+1)}\right)^{1/2},$$
(2)
where $`\alpha =e^2/\mathrm{}c`$; $`Ze`$ and $`J`$ are the black hole charge and spin. A black hole with spin can not have mass smaller than 0.93 $`M_p`$. For spinless black holes this equation gives the minimal mass $`M_{min}=M_p\sqrt{\alpha }=0.085M_p`$. This classical consideration does not take into account the mass renormalization problem. One can not completely exclude that after the “renormalization” the minimal mass of the charged black hole may be as small as the electron mass.
Such elementary black holes may appear in the very early Universe or as a result of “evaporation” of heavier black holes. We know that black holes radiate with a discrete spectrum in a black body envelope (Hawking radiation), but we cannot say whether a final elementary black hole vanishes completely or what the lifetime of such a process would be. For example, it may not be “easy” for a black hole with $`J=1/2`$ (or any half-integer spin) to decay since any such decay involves violation of lepton or barion number (e.g. black hole $`e^{}+\gamma `$; however, if one views this decay as a “disconnection” of the black hole “inner” universe from our Universe no separate conservation laws for each “universe” should be assumed until they are totally disconnected). We will assume that the elementary black holes do not undergo a final radiation.
The first thing we would like to do is discuss a method to search for elementary black holes. An elementary charged black hole moving through matter at less than a few thousand km s<sup>-1</sup> can capture electrons in the same way as an ordinary atomic nucleus, creating a “black hole atom” . One can search for spectral lines in these systems. Normal nuclei are unstable for very large $`Z`$, but a black hole can have any charge at all. Furthermore they can have a negative charge, giving rise to whole new types of systems. In fact just about any electrically charged particle can be bound in a black hole atom.
We would like to know whether or not we can do experiments on black holes in a laboratory. Obviously a neutral black hole can simply fall between atoms in the floor of our laboratory and to the centre of the earth, since it is small and the only force acting on it is gravity. However if we have a charged black hole then there is electromagnetic repulsion between atoms and the black hole which may be large enough to keep it in the laboratory. The contact Coulomb force between a neutral atom and the neutral black hole atom (with a charged black hole nucleus) is $`e^2/a_0^210^7`$ N where $`a_0`$ is the Bohr radius. This equation is just the Coulomb force on the radius of an external electron orbit. Suprisingly, the force due to gravity $`M_pg`$ is also $`10^7`$ N for an elementary black hole with the Plank mass. This means that a black hole has large probability to “tunnel” between the atoms and fall through. A lighter spinless black hole may be an exception since the classical equation (2) gives in this case the minimal mass $`M_p/12`$ . The electromagnetic force in this case is an order of magnitude larger than the gravitational force , and such black hole may stay for some time on the surface of the Earth. In this situation it may be reasonable to do both laboratory and astronomical observations. It is easy to calculate the “isotopic” shifts of the black hole lines relative to the usual atomic lines. We discuss the “atomic” spectrum of an elementary black hole in Sec. II with a view to observation of black hole atoms, and verification of their existence. Due to the large mass the black hole spectral lines do not have Doppler broadening and hyperfine structure. This also may help in identification of such lines.
Note that the upper limit on concentration of the black holes follows from the estimate of their mass density. The Plank mass $`M_p`$ is $`10^{18}`$ times larger than the proton mass. If we assume that the minimal mass of elementary black holes is $`M_p/12`$, and that the dark matter is 100 times heavier than the hydrogen matter and consists of the elementary black holes only, we obtain that the abundance of the elementary black holes in cosmic space does not exceed $`10^{15}`$ of the hydrogen abundance (one can compare this with the abundance of uranium $`3.10^{13}`$). The limit on the concentration of the elementary black holes can be much stronger if one considers other effects - see, for example, explaination of the deficit of solar neutrinos based on the catalysis of nuclear reactions in the sun by charged black holes . This may practically exclude the possibilty of the observation of the black holes in the spectra of very distant objects. We must rely on the observation of the close objects (like sun spectra) or laboratory data. We should also recall another motivation to search for atoms with superheavy nuclei and shifted atomic lines which is related to so called “strange matter” that have nuclei made from the “normal” (up, down) and strange quarks.
There is a finite probability that an orbiting electron or proton could fall into the central black hole and neutralise it’s charge . This gives a lifetime which is estimated in Sec. III. It is found that the lifetime of low $`Z`$ electronic atoms is many orders of magnitude larger than the age of the universe, but that it decreases exponentially with $`Z`$. In fact, we conclude from sections III and IV that primordial black holes would now not have charge greater than about $`Z=70`$.
We also discuss a few questions which may be of the theoretical interest. In Sec. IV we discuss the K-shell states of black hole atoms for $`Z>137`$, where there is a well known singularity in the equations. While the single particle solutions of the Dirac equation for high $`Z`$ has been known for some time, it is worth revisiting these because now we have a physical system where the critical field is realised. We show in this section that the ground state reaches the lower continuum ($`E=mc^2`$) before $`Z=138`$ which means that there is no room for negative energy bound states (recall that for $`Z=1/\alpha 137`$ the relativistic energy $`E0`$) and hence spontaneous positron emission will occur for all $`Z>137`$. We also consider a similar problem for a charged scalar particle where the singularity appears for $`Z>68`$.
If a black hole has a negative charge then it can become a “protonic” atom. The strong force between protons leads to a peculiar ground state characterised by a nucleus orbiting our elementary black hole (Sec. V). It leads to a new stability curve for the captured nucleus, shifted towards more protons (some of the protons still undergo $`\beta `$-decay to neutrons).
Furthermore we show in Sec. VI that the protonic black hole atom leads to a mechanism by which the black hole can gain a colour charge, by capturing a single quark. This naturally leads to an extension of the “no hair” theorem (which states that any black hole can be characterised by it’s mass, electric charge and spin) to include a colour. Discussion of the solutions of Einstein-Yang-Mills equations for coloured black holes can be found, e.g. in Ref. . We estimate the lifetime of these systems and find that a black hole with colour charge (that may be called a “superheavy quark”) will persist for $`10^7`$ years.
We should note that the single quark capture may be forbidden if the energy of the coloured black hole is higher than that of electrically charged black hole. In this case the lifetime of the protonic black hole can be very long (it contains extra factor $`(r_g/fm)^610^{150}`$ where the Planck length $`r_g=\mathrm{}/M_pc=1.6\times 10^{33}`$ cm, $`fm=10^{13}`$ cm) since the probabilty of all three quarks to be near black hole horizon is very small. The same argument may be valid for the electron black hole atom if electron is not an elementary particle, i.e. consists of “pre-quarks”. Therefore, one may view lifetime calculations in the present work as the estimates of the minimal lifetimes.
Finally, in Sec. VII, we briefly consider atoms formed when much heavier (with masses $`10^{12}`$ kg) black holes capture electrons. The electric charge is shown to neutralise very quickly, but there still remains the possibility of short-lived gravitationally bound systems.
## II Black Hole Spectrum
The spectrum of a single electron in a pure central Coulomb potential is given by
$$E=\frac{mZ^2e^4}{2\mathrm{}^2n^2}$$
(3)
where $`n`$ is the principal quantum number and $`Z`$ is the charge at the centre. If we have many electrons, then the outer electron energy is usually described by a Rydberg formula
$$E=\frac{mZ_a^2e^4}{2\mathrm{}^2\nu ^2}$$
(4)
where $`\nu `$ is an effective principle quantum number, $`Z_a1`$ is the ion charge ($`Z_a`$ is the charge that the outer electron “sees”; for neutral atom $`Z_a=1`$). To search for black holes we should calculate the level shifts from normal atoms with the same charge. The normal atom spectra thus provides us with a calibration point.
Normal atom spectra are shifted from the “ideal” spectra because of the finite mass and volume of their nuclei. In an elementary black hole atom, however, the mass of the nucleus is practically infinite (over $`10^{18}`$ proton masses) and it’s volume is zero (around $`10^{60}`$ of nuclear volume). Thus we wish to find the energy levels of normal atoms in terms of the ideal (black hole) spectra plus energy shifts due to the finite mass and volume of their nucleus. To improve the accuracy of the calculations one can use experimental data for isotopic shifts in normal atoms.
The theory of isotopic shifts is presented in numerous books (see, e.g. ). However, it would be instructive to present some results here with a particular application to the black hole line shifts. The mass shift is given by
$$\mathrm{\Delta }E_m=E_AE_{\mathrm{}}=E_{\mathrm{}}.\left(\frac{\mu }{m}1\right)(1+S)$$
(5)
where $`m`$ is the electron mass, $`\mu =mM/(m+M)`$ is the reduced mass, $`M`$ is the mass of the nucleus, $`E_A`$ is the energy of atomic level, $`E_{\mathrm{}}`$ is the energy of the black hole level (corresponding to infinite $`M`$) , $`S`$ is the correction due to the specific mass shift which exists in many-electron atoms. The contribution of the specific shift to the transition frequencies can be calculated or extracted from the experimental data on isotopic shifts (see below).
The volume shift is given by
$$\mathrm{\Delta }E_V=e\left(\varphi Ze/r\right)\psi ^2\left(r\right)𝑑V$$
(6)
where $`\varphi `$ is the Coulomb nuclear potential. While this integral is formally extended to all space, it is practically zero outside of the nucleus. For relativistic s-wave electrons
$$\mathrm{\Delta }E_VE_{\mathrm{}}\frac{6}{\nu }\frac{\gamma +1}{\gamma \left(2\gamma +1\right)\left(2\gamma +3\right)\left[\mathrm{\Gamma }\left(2\gamma +1\right)\right]^2}\left(2Z\frac{r_0}{a_0}\right)^{2\gamma }$$
(7)
where $`\gamma =\sqrt{1\left(Z\alpha \right)^2}`$, $`a_0=\mathrm{}^2/me^2`$ is the Bohr radius, $`r_0=1.2A^{1/3}`$ is the nuclear radius and $`A`$ is the mass number of nucleus. To calculate the parameter $`\nu `$, we simply use eq. 4, since $`E`$ is the known ionisation energy . The estimates of the normal mass shifts and volume shifts of electron levels in atoms relative to “ideal” black hole atoms are presented in fig. 1 and fig. 2.
Consider two examples. In hydrogen the volume shift is very small. We will specifically look at the $`2p_{3/2}1s_{1/2}`$ transition. The spectral data is given in , which gives for hydrogen and deuterium
$`\omega _H`$ $`=`$ $`82259.279\mathrm{cm}^1`$
$`\omega _D`$ $`=`$ $`82281.662\mathrm{cm}^1`$
Using the nuclear masses given in we obtain
$$\omega _{\mathrm{}}\omega _H=44.801\mathrm{cm}^1$$
(8)
It turns out that this shift is the same for the $`2p_{1/2}1s_{1/2}`$ transition.
Now consider Mg II. For magnesium we must include the volume shift. The two isotopes used are <sup>24</sup>Mg and <sup>26</sup>Mg. From eq. (5) we then obtain
$`\omega _{26}\omega _{24}`$ $`=`$ $`\omega _{\mathrm{}}\left({\displaystyle \frac{\mu _{26}}{m}}{\displaystyle \frac{\mu _{24}}{m}}\right)(1+S)\mathrm{\Delta }E_{V26}+\mathrm{\Delta }E_{V24}`$ (9)
$`\omega _{\mathrm{}}\omega _{24}`$ $`=`$ $`\omega _{\mathrm{}}\left(1{\displaystyle \frac{\mu _{24}}{m}}\right)(1+S)+\mathrm{\Delta }E_{V24}`$ (10)
Data for the MgII 2796 line (corresponding to the $`3p_{3/2}3s_{1/2}`$ transition) have been obtained which give
$`\omega \left({}_{}{}^{24}\mathrm{MgII}2796\right)`$ $`=`$ $`35760.834\pm 0.004\mathrm{cm}^1`$
$`\mathrm{\Delta }\nu \left({}_{}{}^{26}\mathrm{Mg}^{24}\mathrm{Mg}\right)`$ $`=`$ $`3.050\pm 0.1\mathrm{GHz}`$
$`\omega _{26}\omega _{24}`$ $`=`$ $`0.1017\pm 0.003\mathrm{cm}^1`$
Now using the formula for volume shift eq. (7) we obtain
$`\mathrm{\Delta }E_{V24}`$ $`=`$ $`0.0315\mathrm{cm}^1`$
$`\mathrm{\Delta }E_{V26}`$ $`=`$ $`0.0333\mathrm{cm}^1`$
And so eqs. (9,10) give
$$\omega _{\mathrm{}}\omega _{24}=1.378\mathrm{cm}^1$$
(12)
## III Lifetime
There is a finite probability that the electrons orbiting a black hole atom enter the black hole. This “direct capture” mechanism leads to a lifetime for the black hole atom. The simplest estimate of the capture probability is given by the product of black hole horizon area $`4\pi r_g^2`$ and flux for a particle moving with speed c near this horizon, $`j=c\left|\psi _s\left(0\right)\right|^2`$.
$$\frac{1}{\tau }=w\left|\psi _s\left(0\right)\right|^24\pi r_g^2c$$
(13)
Now the K-shell electron is the most likely to be captured. Using relativistic electron wave function from we obtain
$$w\frac{4r_g^2cZ^3}{a_0^3}\left(\frac{a_0}{2Zr_g}\right)^{2(1\gamma )}$$
(14)
A complete picture is given in fig. 3. In all numerical estimates we assume that the mass of the black hole is equal to $`M_p`$. For $`Z=1`$ this gives lifetime about $`10^{22}`$ years, for $`Z=70`$ about the age of the Universe. Thus, elementary black holes formed at the Big Bang with $`Z70`$, would have captured electrons in the K-shell, reducing the charge of the black hole. This process would have continued until the present day. Since the lifetime exponentially decreases with $`Z`$ we would not expect black holes with $`Z`$ larger than about 70 to exist today.
Note that after the electron has been captured by the black hole, the quantisation rules of the black hole (see Introduction) are no longer obeyed. Thus the black hole must undergo some process to correct itself, such as radiate.
## IV Critically Charged Black Holes
When a black hole forms it inherits the (conserved) quantum numbers of the particles used to form it. This means that most black holes are formed with some non-zero electric charge $`Z`$. The evaporation process should reduce this charge, however a final charge can still be large. There are some results which are not possible in normal atoms but which may be observed in black hole atoms. For example there are two well known singularities in the single particle solution of the Dirac equation for a Coulomb field, corresponding to particular $`Z`$ (see e.g. ). These single particle solutions are very good approximations for the K-shell (ground state) energy because these electrons are closest to the nucleus and only very lightly screened by the other electrons.
The first singularity occurs when $`Z\alpha `$ becomes greater than one (Sec. IV A). This singularity is removable when the finite size of the nucleus is taken into account. The second occurs when the ground state energy reaches the lower Dirac sea (Sec. IV B). We call the charge corresponding to this the supercritical charge.
### A First critical charge
The solution of the Dirac equation for a Coulomb field $`V=Z\alpha /r^2`$ has a singularity at $`Z\alpha =1`$ (or $`Z=137.04`$). At this point the parameter $`\gamma =\sqrt{1\left(Z\alpha \right)^2}`$ and the ground state energy $`\epsilon =m\gamma `$ become imaginary. In fact nothing much should change when this critical point is passed in a system with a finite nucleus. The energy, $`\epsilon `$, becomes negative, but this is entirely allowable.
The physical reason for this singularity is that when $`Z\alpha >1`$ the particle can “fall to the centre”. The effective potential $`U\left(r\right)`$, which arises when the Dirac equation is squared and put into a Schrödinger equation - like form, behaves, for a Coulomb field, in a singular manner $`U\left(r\right)\left(j\left(j+1\right)Z^2\alpha ^2\right)r^2`$ as $`r0`$. This leads, as it does non-relativistically, to all bound state wavefunctions having an infinite number of nodes when $`Z\alpha >j+1/2`$.
To determine the level energies it is necessary to specify the potential $`V\left(r\right)`$ and a boundary condition at zero. In any physical system this is obtained by an alteration of the potential at small $`r`$, due to finite size of the nucleus, or in our case, the black hole. The form of the potential at small $`r`$ is not important when $`Z\alpha <1`$ (in fact in our case the black hole should not significantly affect the energy levels even at $`Z=137`$, since the scale of the cutoff, $`r_0=10^{35}`$ m, is so small). Yet it does become very important when $`Z\alpha >1`$.
A lot of interesting physics comes into play when the energy of the K-shell electron goes below zero (but not into the lower continuum). The electrons will have an energy $`\epsilon =\mathrm{const}.m`$. Therefore it is energetically favourable for the electron to undergo a “ beta decay” into heavier and heavier particles, like a muon, tauon or perhaps some new grand unification particle. This process would be very fast, and in fact it may be possible for this to occur before the particle is captured by the black hole.
Unfortunately in the elementary black hole system none of these effects are realised because the ground state energy falls into the Dirac sea before $`Z=138`$ (Sec. IV B, see figure 4), and hence there are no negative energy bound states. But the critical value of $`Z`$ is dependent on the radius of the black hole, and hence it may be possible for bound states with $`m<\epsilon <0`$ to exist in heavier black holes.
### B Supercritical charge
If a black hole has a charge larger than some critical value $`Z_c`$ then the K-shell electron energy will reach the lower continuum (the lower Dirac sea) corresponding to energy $`\epsilon =mc^2`$. This corresponds to a binding energy of $`2mc^2`$. This field can spontaneously polarise the vacuum to create an electron - positron pair.
In Appendix A we follow the method outlined by Popov to find that the supercritical charge $`Z_c=137.29`$. Thus when $`Z=138`$ we are already in the lower continuum, and so we can conclude that there are no K-shell electron bound states with energy $`1<\epsilon <0`$.
When $`Z>137`$, then, the field can spontaneously polarise the vacuum to create an electron - prositron pair. The electron goes into the K-shell and the positron goes to infinity (in an ordinary nucleus there would be spontaneous emission of two positrons, after which the effective charge of the nucleus decreases by two units, corresponding to filling of the K-shell). In the black hole atom the electron would immediately fall into the black hole, decreasing it’s charge by 1.
This process has a finite lifetime, since the positron must tunnel to infinity to overcome the positive potential it would feel from the black hole. The spontaneous emission of positrons would continue until the charge of the black hole fell below $`Z=138`$. The ground states with $`\epsilon >0`$ also have a finite lifetime, so the charge of the black hole then continues to fall until the charge is neutralised with the lifetimes described in Sec. III.
It is also interesting to consider the binding between black hole and charged scalar particle: Higgs boson, $`\pi `$-meson, etc. For a scalar particle, the Klein-Gordon equation in the Coulomb field has a singularity at $`Z\alpha =1/2`$ (or $`Z=68.5`$). After this we must take into account the finite size of the nucleus. We wish to know at which point the energy reaches the lower continuum, to see whether there are any bound states with negative energy (in much the same fashion as was done for the electron case).
The calculation to find the critical charge in the case of the scalar particle is done in Appendix B, again following the method of Popov . It is found that the bound state energy of the lowest shell equals $`mc^2`$ at $`Z_c=69.001`$. Due to uncertainty in the size and boundary condition of the black hole we cannot, from these results, tell if there is a $`Z=69`$ bound state or not. Anyway, it would be a very short-lived state.
In a case of a finite-size particle like $`\pi `$-meson the low-$`Z`$ Klein-Gordon equation may be treated in a similar fashion to the “finite-size nucleus” problem with the radius of $`\pi `$-meson instead of nuclear radius. However, accurate high-$`Z`$ results may be obtained only by solution of the three-body problem and depends on the strong interactions between the quarks. For example, one of the quarks can be rapidly captured by a black hole and they can form a “superheavy quark” interacting with a remaining quark - see next section devoted to a protonic black hole.
## V Protonic Black Hole Atoms
If a black hole has a negative charge, then it can capture positively charged particles including protons and alpha particles. At first glance, one might think that this would form a “protonic black hole atom”, something akin to usual atoms, but with orbiting protons instead of electrons. But in fact the physics of this system is quite different.
### A Ground state
Consider a protonic black hole. Because of the strong nuclear force, the ground state of this system would consist of a “nucleus” of protons (some of which may decay into neutrons) orbiting a black hole. For a singly charged black hole, there would simply be a single bound proton and nothing very exciting occurs (until it decays - see Sec. VI).
It is known that for two nucleons there are no bound isotopic triplet states. This means that in a system of two protons orbiting a black hole of charge $`Z=2`$, one of the protons should decay via the weak interaction as
$$p^+n+e^++\nu $$
(15)
and thus acheive a stable deuterium nucleus configuration. However this is not the end of the story. The positron would be emitted with a large kinetic energy and so it would only be weakly bound if at all. Thus the system will still have an effective charge of $`1`$, which can attract another proton. This would join with the deuterium nucleus and form a <sup>3</sup>He nucleus.
Also it is concievable that the doubly charged black hole could pick up an extra neutron, if a sufficiently slow one chanced by the orbiting nucleus. Then the tritium nucleus would be formed. This would certainly be unstable, as it is usually, and decay to <sup>3</sup>He.
We now have the idea that the negatively charged black hole can have a nucleus orbiting it, having a certain number of neutrons and protons which follow the nuclear stability curve. But we have not considered the effect of the Coulomb field on this nucleus. The binding energy of the nucleus will decrease when protons decay to neutrons because the neutrons are not charged. This may provide enough energy in some cases to form a bound state with an unusually high number of protons. So the stability curve is effectively pushed towards the proton side.
### B Lifetime
The lifetime found from eq. 14 is proportional to $`m^{2\gamma 1}`$, where $`m`$ is the mass of the orbiting particle. This means that a single proton orbiting around an elementary particle of charge $`Z=1`$ has a decay probability $`\left(m_p/m_e\right)^3=1836^3`$ times larger, leading to a lifetime of $`\tau =3\times 10^{12}`$ years.
But the mass of our particle will be that of the orbiting nucleus, which must in turn follow the new stability curve for a nucleus orbiting a black hole (Sec. V A). Thus the lifetimes will decrease even faster than they do in the electron case because the entire nucleus mass is of importance, not the mass of just one proton. However, if we consider black hole with only one proton the lifetime is equal to the age of the Universe for $`Z=5`$.
There is an additional complication to the lifetime, considered in Section VI.
## VI Colour Charge in Black Holes
We have established in Section V that negatively charged black holes can have orbiting protons. We even estimated the probability that the bound protons fall into the black hole. But we did not consider that protons are not elementary particles, but are made up of three quarks, each with a different colour charge. This means that due to the extremely small size of the elementary black hole, it is not the entire proton which is captured, but a single quark!
The black hole then obtains the colour of the quark it captured, and becomes a “superheavy quark”. This would form a strongly bound state with the remaining two quarks of the original proton. While the other two quarks will eventually fall into the black hole, this will happen with a finite lifetime $`\tau 10^7`$ years. Here we assumed that the quark wave function in the proton has the radius $`R0.5`$ fm and $`\psi (0)^2=1/(4\pi R^3/3`$).
## VII Heavier Black Holes
Let us discuss briefly non-elementary black holes. Because the density fluctuations in the early universe would have occurred on all scales , there may conceivably have been black holes formed with much larger masses. The smaller ones of these would have evaporated into elementary black holes. The minimal mass of a black hole that would not have evaporated entirely by now is $`M5\times 10^{11}`$ kg . In this chapter we will consider black holes of mass $`M10^{12}`$ kg. Such an object has a radius of 1.5 fm.
These objects are interesting to study because at this mass they have gravitational fields comparable to the electric field, as well as sizes comparable to that of ordinary nuclei. Unfortunately there is another complication in determining orbits - the constant flux of Hawking radiation being emitted will interfere with any particle orbits.
Consider a singly charged black hole. Neglecting the gravitational potential and using equation 14 for capture probability with a new Schwarzschild radius of 1.5 fm, we obtain the lifetime
$`\tau _{\mathrm{Coulomb}}=5\times 10^{11}\mathrm{sec}`$
Obviously increasing the charge makes this smaller still, as does including the gravity term. Thus we can conclude from this, and the fact that if it was charged, evaporation would favour neutralisation of charge, that any charge of the black hole would have been neutralised by now.
Having concluded that there would be no electric charge on black hole atoms, we turn our attention to gravitational atoms. These would consist of particles (we will confine ourselves to electrons, but any particle would do) orbiting heavy black holes of the mass discussed. The potential is
$$U_{\mathrm{grav}}=\frac{GMm}{r}=\frac{g}{r}$$
(16)
In relativistic units, $`g=1/520`$ for electrons. The radius of the ground state for such a system is approximately $`a=1/gm=2\times 10^{11}`$ m $`=0.2`$ nm. This does not take into account the continual radiation of the black hole which can interact with the particles and may potentially destroy the bound states.
## VIII Conclusion
We have seen that elementary black holes with charge can capture electrons and form bound states similar to those of an ordinary atom. The electronic spectra of these black hole atoms can be searched.
The electrons orbiting a black hole can fall into it. This “direct capture” leads to a finite lifetime for black hole atoms. This was found to be many times the age of the universe for black holes with small charges, but decreasing exponentially with increasing Z. From these calculations we concluded that primordial black holes would today have charges no larger than $`Z70`$.
The black hole atoms may give rise to new physical phenomena. When $`Z>137`$ ($`Z>68`$ for scalar particles) we have a physical realisation of the much theorised supercritical Coulomb fields, where the single particle K-shell energy becomes negative. We found that these energies immediately drop into the lower continuum $`\epsilon <mc^2`$. Although the ground state falls into the Dirac sea at $`Z\alpha =1.002`$, the upper states ($`p_{3/2}`$, $`d`$, etc.) do not until after $`Z\alpha =\left|\kappa \right|`$ (see ). Therefore it remains as an interesting question whether there are any states with negative energy (including $`mc^2`$) in the upper bound states. If there are, then there is still the possibility of beta decay of electrons to muons, tauons, and so on.
We also discussed protonic black hole atoms which are formed with negatively charged elementary black holes. The ground state of such systems is a nucleus orbiting a black hole, with the nuclear stability curve pushed towards the proton side. Because protons are not elementary particles, the black holes would capture just one of the quarks of the proton (or of any constituent nucleon in the orbiting nucleus). This leads to a black hole with a colour charge. The lifetime of this “superheavy quark” was found to be about $`10^7`$ years.
Finally we briefly considered much heavier black holes (with masses of $`10^{12}`$ kg) and found that their charge is neutralised very rapidly. There remains the possibility of gravitational atoms. However, because of Hawking radiation and the capture of the electron, such atoms may be very short-lived.
We are grateful to M.Yu.Kuchiev for valuable discussion. This work was supported by the Australian Research Concil.
## A Supercritical charge - electron case
Here we proceed to find the supercritical $`Z`$, following the method outlined by Popov . To find $`Z_c`$ we need to solve the Dirac equation for $`\epsilon =1`$ (we use relativistic units $`\mathrm{}=c=m=1`$). We write the potential as
$$V\left(r\right)=\{\begin{array}{c}\xi /r,r>r_g\hfill \\ \xi /r_g,0<r<r_g\hfill \end{array}$$
(A1)
where $`\xi =Z\alpha `$. This assumes that the wavefunction can extend inside the black hole, and that the charge of the black hole is concentrated entirely on the surface of it. The Dirac equation can be expressed
$`{\displaystyle \frac{dF}{dr}}`$ $`=`$ $`{\displaystyle \frac{\kappa }{r}}F+\left(1+\epsilon V\right)G`$ (A2)
$`{\displaystyle \frac{dG}{dr}}`$ $`=`$ $`\left(1\epsilon +V\right)F+{\displaystyle \frac{\kappa }{r}}G`$ (A3)
where $`F\left(r\right)=rf\left(r\right)`$ and $`G\left(r\right)=rg\left(r\right)`$, with $`f`$ and $`g`$ the radial wavefunctions, $`\kappa `$ is the eigenvalue of the operator $`K=\beta \left(L\sigma +1\right)`$, conserved in a spherically symmetric field. For K-shell electrons, $`\kappa =1`$.
We eliminate $`G`$ from eq. A3 to obtain the second order differential equation
$$F^{\prime \prime }+\frac{V^{}}{1+\epsilon V}\left(F^{}+\frac{\kappa }{r}F\right)+\left[\left(\epsilon V\right)^21\frac{\kappa \left(\kappa +1\right)}{r^2}\right]F=0$$
(A4)
The wavefunction outside the black hole can be obtained from this, setting $`\epsilon =1`$ and $`V=\xi /r`$ where $`\xi `$ now becomes it’s critical value when the bound state reaches the lower continuum, $`\xi _c`$,
$$r^2F^{\prime \prime }+rF^{}+\left(\xi ^2\kappa ^22\xi r\right)F=0$$
(A5)
Transforming this equation by $`x=\sqrt{8\xi r}`$ we obtain the Bessel equation
$$x^2F^{\prime \prime }+xF^{}\left(x^2\nu ^2\right)F=0$$
(A6)
with $`\nu =2\sqrt{\xi ^2\kappa ^2}`$. This leads to the solution (for $`\kappa =1`$ and $`\xi >1`$)
$$F=\mathrm{constant}K_{i\nu }\left(x\right)=\mathrm{constant}K_{i\nu }\left(\sqrt{8\xi r}\right)$$
(A7)
where $`K_{i\nu }\left(x\right)`$ is the MacDonald function of imaginary order (a Bessel function of the second kind). The second independent solution $`I_{i\nu }\left(\sqrt{8\xi r}\right)`$ is unacceptable because of it’s growth at infinity.
Now we must choose a boundary condition at $`r=r_g`$, hence specifying completely the wavefunction $`F`$ and allowing us to find $`\xi _c`$. For the potential considered in eq. A1, we equate the logarithmic derivatives of the wavefunctions inside and outside the black hole (for details see ) and obtain the transcendental equation
$$xK_{i\nu }^{}\left(x\right)=2\xi \mathrm{cot}\xi K_{i\nu }\left(x\right)$$
(A8)
Solving this for $`r_g=1.6\times 10^{35}\mathrm{m}=4\times 10^{23}`$ in the relativistic units used we obtain $`\xi =Z_c\alpha =1.00187`$ corresponding to a critical charge
$$Z_c=137.29$$
(A9)
Thus we conclude that there are no (K-shell) bound states with energy $`1<\epsilon <0`$ because when $`Z=138`$ we are already in the lower Dirac sea.
One may think that the boundary condition chosen is somewhat artificial, but in fact the actual boundary condition chosen is not important compared to the scale of $`r_g`$. Consider, for example a very general (and incomplete) boundary condition $`K_{i\nu }\left(\sqrt{8\xi r_g}\right)=\mathrm{constant}`$. The constant should necessarily be positive because the ground state should have no nodes. The maximum value of $`\xi `$ will be realised when a node does exist at the boundary. So let’s try $`K_{i\nu }\left(\sqrt{8\xi r_g}\right)=0`$. This gives us a value of $`\xi _c=1.00199`$ corresponding to $`Z_c=137.31`$. So our conclusion that no negative energy bound states exist is valid.
## B Supercritical charge - scalar case
We can discuss the case of a point-like scalar particle in a Coulomb field in an analagous manner to the electron case . Firstly, we must solve the Klein-Gordon equation, outside the nucleus
$$\phi _l^{\prime \prime }+\left(\epsilon ^21+\frac{2\epsilon \xi }{r}+\frac{\xi ^2l\left(l+1\right)}{r^2}\right)\phi _l=0$$
(B1)
We can solve this in it’s present form, but it is easier if we set $`\epsilon =1`$ immediately, since we are interested in finding $`\xi _c`$, where the K-shell energy meets the lower continuum. The wavefunction in this case has the form
$`\phi _l\left(r\right)`$ $`=`$ $`\sqrt{r}K_{i\mu }\left(\sqrt{8\xi R}\right)`$ (B2)
$`\mu `$ $`=`$ $`2\sqrt{\xi ^2\left(l+{\displaystyle \frac{1}{2}}\right)^2}`$ (B3)
If we again use the cut-off potential for $`V\left(r\right)`$ then we obtain for $`\xi _c`$ the transcendental equation (for the lowest level $`n=1`$, $`l=0`$)
$`xK_{i\nu }^{}\left(x\right)`$ $`=`$ $`\left(2\beta \mathrm{cot}\beta 1\right)K_{i\nu }\left(x\right)`$ (B4)
$`\nu `$ $`=`$ $`\sqrt{4\xi ^21}`$ (B5)
$`\beta `$ $`=`$ $`\sqrt{\xi \left(\xi 2R\right)}`$ (B6)
$`x`$ $`=`$ $`\sqrt{8\xi R}`$ (B7)
Solving this numerically we obtain that $`\xi _c=0.50353`$ which means that the energy of the scalar particle reaches the Dirac sea at $`Z_c=69.001`$. Due to uncertainties in the size and boundary condition of the black hole, we cannot tell if there is a bound state with negative energy at $`Z=69`$. |
warning/0001/astro-ph0001518.html | ar5iv | text | # Molecules as Tracers of PN Structure
## 1. Introduction
The matter that evolves into a planetary nebula (PN) is initially ejected during the AGB and proto-PN phases, largely in the form of molecular gas. The structure in this gas and its response to ionizing radiation and fast winds or jets from the central star are important features of PN formation. In this paper we briefly describe and comment on some recent results which focus on this structural development.
## 2. Precursor Envelopes and Multiple Shells
We turn first to the structure of AGB envelopes which is basic for understanding PNe since it determines the environment in which the nebulae form. It is also crucial for understanding the mass-loss process and has been widely discussed (e.g., Olofsson 1999). A recent, valuable contribution on this issue is the atlas of circumstellar envelopes by Neri et al. (1998), which contains maps of 46 envelopes in the CO (1–0) and (2–1) lines, made using the IRAM telescopes. With respect to the structure of the envelopes, the authors conclude that within the sensitivity of their observations, AGB and post-AGB envelopes are for the most part spherical, with evidence for an inner shell and/or significant asymmetry in 30% of the sample. In fact, most of these cases are in the former category, so the truly asymmetrical envelopes are not very common. This is in marked contrast to the widespread asymmetries in the ionized gas of PNe, discussed throughout this volume.
One interesting question that arises in this context is: just how regular and symmetric is the mass loss of a typical AGB star? A good example with which one can address this question is the archetype of carbon-rich AGB stars, IRC+10216. The molecular envelope of IRC+10216 has been intensively observed for more than a decade using millimeter interferometry, in order to study the chemistry. Recent results are reviewed by Lucas & Guélin (1999). A key point relevant to the structure of the envelope, is that some molecular species are observed to form more than one layer around the star, and the peaks of several species coincide, which suggests the presence of an underlying physical shell structure. About three distinct shells can be identified within $`25^{\prime \prime }`$ of the star (see Fig. 4 of Lucas & Guélin 1999).
This multiple shell picture of the envelope of IRC+10216 has recently been extended at optical wavelengths (Mauron & Huggins 1999). Deep images obtained with the CFHT reveal the envelope of IRC+10216 in dust-reflected, Galactic light. The left hand panel in Fig. 1 shows a $`V`$-band image, with a field of 223<sup>′′</sup>, and the the right hand panel shows a close-up of the center, enhanced to show details.
The figures show that the envelope is not smooth, but consists of a series of nested, limb-brightened shells. The inner shells show a rough correspondence to those in the molecular gas. The shells are not complete in azimuth, and are separated by somewhat irregular intervals that correspond to time scales of $`200`$–800 yr. This structure appears to be a basic feature of the mass-loss process but is not yet understood: the time scale is much longer than the stellar pulsation period, but is much shorter than the interval between the thermal pulses.
The presence of these shells in the archetype AGB envelope is especially significant in the context of recent HST observations of multiple arcs (with similar spacings) in a half-dozen proto-PNe and young PNe (see Terzian, this volume). The number of cases found already suggests that they may be common, and may well be the norm for certain classes of objects. These arcs can almost certainly be identified with the shells seen in IRC+10216 – at a more advanced stage of evolution – so their origin can be traced back to the formation of the molecular/dust shells in the precursor AGB envelopes.
## 3. Molecular Gas in PNe
In spite of major changes in the mass loss and the onset of photo-ionization through the post-AGB phase, a significant component of molecular gas is found in many *bone fide* PNe, even in highly evolved cases. The most widely detected molecular signatures are the 1.3 and 2.6 mm lines of CO and the 2 $`\mu `$m vib-rotational lines of H<sub>2</sub>. Recent survey work has been reported by Huggins et al. (1996) and Kastner et al. (1996) for CO and H<sub>2</sub>, respectively, with more than 40 PNe detected in each species. Additional examples are continually being added to the lists (e.g., Josselin et al. 1999; Hora & Latter 1999). Other molecular species are also detected in the neutral gas, even in some evolved cases (Bachiller et al. 1997), and various aspects of the chemistry have recently been discussed by Natta & Hollenbach (1998), and Howe & Williams (1998).
The molecular gas is found predominantly in PNe at low Galactic latitude, and there is a strong correlation with morphology (e.g., see Fig. 2, right panel), which suggests that we detect the nebulae with higher mass progenitors. There is no doubt about the location of the gas. Except in the youngest PNe where the envelopes may still completely enshroud the nebula, the molecular gas is found around the waist of the ionized gas, in shapes variously described as rings, cylinders, or toroids.
One example which illustrates the relation of the molecular gas to the nebula is M2-9. This young ($`\tau _{\mathrm{exp}}1,500`$ yr), bipolar PN is well known from imaging with HST. It is seen nearly sideways-on, and is dominated by elongated, twin outflows or jets. The H<sub>2</sub> emission in M2-9 lies along the edges of the cylindrical structure formed by the jets (Hora & Latter 1994). The CO emission, which traces the densest gas, is found only in a slowly expanding ($`V_{exp}`$ = 7 km s<sup>-1</sup>) torus which fits tightly around the waist of the nebula (Zweigle et al. 1997).
A second example is the Helix nebula (NGC 7293) which is a much older system ($`\tau _{\mathrm{exp}}10,000`$ yr), seen nearly end-on. The 2 $`\mu `$m H<sub>2</sub> emission has been imaged by Kastner et al. (1996) and the whole nebula has recently been mapped in the CO(2–1) line by Young et al. (1999). The left panel of Fig. 3 shows an $`R`$-band image of the Helix (mainly H$`\alpha `$ and \[N ii\]), and the next panel shows the CO map. These two views are quite similar: it is evident that the ionized nebula abuts the envelope and has formed through photo-ionization of the neutral gas. Thus, in both the Helix and M2-9, and in many other PNe as well, the molecular gas is a main structural feature of the nebula and an important key to its morphology.
The mass of molecular gas in the PN envelopes is not easy to determine, even though the major constituent (H<sub>2</sub>) is directly observed. This is because the 2 $`\mu `$m emission arises from high lying levels ($`>6000`$ K) and typically samples only a small fraction of the gas. (For recent work on whether shocks or ultraviolet radiation excite the lines, see Hora & Latter 1999). More useful mass estimates come from the low lying lines of CO, which are thermalized under a wide range of conditions. The CO fluxes can be used to estimate the total number of CO molecules, and with reasonable assumptions on the CO abundance lead to an estimate of the mass. Masses estimated in this way range up to a few $`M_{}`$.
To illustrate the evolution of the envelopes, the mass ratio of molecular to ionized gas in a large number of PNe is shown in Fig. 2 (left panel) vs. PN size – which is a rough measure of the age. The upper envelope of this plot defines a striking evolutionary trajectory for PNe with a significant molecular component (statistically the bipolar PNe): the young PNe are dominated by the envelope, which remains a significant component of the circumstellar gas until they reach a size of $`0.1`$ pc. In other PNe, the molecular gas is photo-dissociated more rapidly.
In addition to the molecular gas, there is, not surprisingly, ample evidence for neutral atomic gas in PNe, e.g., from observations of the 21 cm line and infrared fine structure lines. This gas is in interface regions between the molecular and ionized gas, and in envelopes that are essentially completely atomic. The masses in these components are substantial (e.g., Taylor et al. 1990; Dinerstein et al. 1995; Young et al. 1997), but they have not yet been studied in large numbers of PNe or at high angular resolution, so they do not yet provide a systematic or detailed picture.
## 4. Large and Small Scale Structure in the Molecular Gas
The detailed structure of the molecular gas in PNe is of considerable interest since it contains information on the physical processes that produce the nebulae. For space considerations, we focus here on one example, the Helix nebula, which is among the nearest PNe and can be examined with high spatial resolution.
One aspect of the structure in the gas is the high degree of fragmentation. The cometary globules located within the inner, ionized gas are well known from optical imaging with HST, and have been discussed by O’Dell & Burkert (1997). Millimeter CO observations have shown that the globules have dense cores of molecular gas (Huggins et al. 1992); in fact without this structure it would be hard to account for their presence in large numbers, because they would be rapidly photo-ionized and disperse on short times scales, compared to the age of the nebula. Recent work by Meaburn et al. (1998) shows that the globules have kinematic similarities to the inner CO ring seen in Fig. 3, and it seems likely that they share a common origin with the more extended molecular envelope.
This envelope itself is also highly fragmented. The clumpy structure in the CO map in Fig. 3 (resolution 30<sup>′′</sup>) is seen to consist of many smaller clumps at higher resolution (see Huggins 1999). The masses of the clumps cannot be more than an order of magnitude more than the cometary globules in the ionized gas, and they share remarkably narrow line widths of $`1`$ km s<sup>-1</sup>. These fragments appear to be close cousins of the cometary globules, and probably represent a slightly earlier stage in the development of these structures.
This picture of a fragmented torus is confirmed by observations with ISO (Cox et al. 1998). Spectroscopy at 5–17 $`\mu `$m with ISOCAM reveals a strong, pure ($`v=00`$) rotational spectrum of H<sub>2</sub> in the Helix, from the S(2) to the S(7) lines. As one moves from the ionized cavity to the limb, the spectrum changes from nearly featureless, to being dominated by the H<sub>2</sub> emission. The line intensities indicate an excitation temperature of $`900`$ K, thus the H<sub>2</sub> is warmer than the CO, and probably forms a skin on the cooler gas. The S(5) line emission dominates the ISO LW2 filter so it has also been possible to image the large scale distribution of H<sub>2</sub> over the whole nebula. The image, at 6<sup>′′</sup> resolution, is shown in Fig. 3, right panel. It shows finer details than the CO map, including flocculent radial rays and clumps of globules around the inner periphery, and generally underscores the completely fragmented picture of the molecular gas.
A second aspect of the distribution of the molecular gas is its global structure. Fig. 3 shows that it is not simply a regular torus, and the apparent “double ring” seen in optical images has long been a target for speculation. The large scale geometry has recently been discussed by Young et al. (1999) using their CO observations. The CO map in Fig. 3 actually consists of 3425 spectra which record the velocity of the gas along the line of sight as well as the intensity. Since the system is expanding, these can be used to study the 3-dimensional structure.
When examined in this way, the data indicate that the inner ring in Fig 3, is a true ring, tilted $`37^{}`$ to the line of sight. The outer arcs to the east and west peel away from the ring (from the north and south, respectively), with a remarkable degree of point symmetry in their geometry (see Fig. 7 in Young et al. 1999). For each point in the structure, there is a similar one on the opposite side of the central star. There seems little doubt that this structure, which dominates the nebula, was formed at an early stage by the interaction of the envelope with collimated, bipolar outflows or jets.
## 5. Shaping the Envelopes
The actual shaping of envelopes by collimated outflows or jets from the central star system is well documented in a number of proto-PNe and young PNe by high resolution observations of the molecular gas. These observations are especially interesting in light of the growing number of PNe that are seen to exhibit point symmetries (especially young PNe observed with HST, see Sahai this volume). These symmetries indicate that collimated outflows – and their interactions – are common in young PNe, and can provide a general shaping mechanism.
The presence of a molecular (or atomic) envelope is an interesting part of this scenario: it not only affects the dynamics of the interaction, but the high densities and low temperatures of the neutral gas preserves the results of the interactions for much longer than the ionized gas. Thus at later stages, this structure can dominate the appearance of the PNe, as is the case in the Helix nebula described above.
Examples of outflow-envelope interactions in the proto-PN phase are described elsewhere in this volume (e.g., by Lucas, and Alcolea), and we focus here on young PNe. The most spectacular outflows are those in KjPn 8 . These consist of pairs of bipolar jets, the most recent of which have expansion velocities $`200`$ km s<sup>-1</sup> and extend over $`4^{}`$, even though the ionized nebula is only $`4^{\prime \prime }`$ in size (see López, this volume, for details). Using the IRAM interferometer, Forveille et al. (1998) have recently mapped the molecular emission from the KjPn 8 system (Fig. 4, left panel). The molecular gas forms a disk that surrounds the ionized nebula (which corresponds to the hole in the center of the figure), and exceeds it in mass by a factor of $`60`$, so it dominates the circumstellar environment.
The disk axis and the jets are aligned in KjPn 8, and their expansion time scales are similar (a few thousand yr). A likely scenario is that common (or related) mechanisms ejected the molecular gas and formed the jets close to the central star. The jets and the gas clearly interact. In the right panel of Fig. 4, the CO channel map at the systemic velocity shows a wind swept disk, with gas extending along the edges of the jet cavities; other channels show that this gas has the highest velocities, and is entrained in the flows. In this case, a primary result of the interaction is the shaping of the disk or torus.
A second example of outflow-envelope interactions occurs in He3-1475. In this young PN, high resolution optical imaging with HST (Borkowski et al. 1997) shows high velocity, bipolar outflows with large opening angles that are apparently focussed into narrow jets by interactions with their surroundings. Recent CO observations with the IRAM interferometer (Forveille, private communication) reveal a torus of molecular gas around the outflows. The torus extends a few arc seconds along the flows, roughly corresponding the region in which the flows are collimated; much of the molecular gas is also at high velocities and is entrained. The interesting perspective here is that the outflow-envelope interactions are complex: the molecular torus is probably the focusing agent of the jets, while at the same time the flows are acting to shape and disrupt the envelope. Further studies of these types of shaping interactions are currently in progress.
## 6. Concluding Remarks
The examples described here illustrate and underscore the important role of molecular gas in the structure of PNe.
* In the equatorial tori that dominate the evolution of bipolar and related PNe for much of their lives.
* As the dense cores of globules and other microstructures,
* As the origin of multiple arcs.
High resolution observations of the molecular gas are also beginning to reveal details of the shaping of PNe by envelope interactions with collimated outflows or jets at an early phase.
### Acknowledgments.
It is a pleasure to acknowledge collaborative programs with R. Bachiller, P. Cox, T. Forveille, and N. Mauron, which form part of this paper. This work was supported in part by NSF grant AST-9617941.
## References
Bachiller, R., Forveille, T., Huggins, P. J., & Cox, P. 1997, A&A, 324, 1123
Borkowski, K. J., Blondin, J., & Harrington, J. P. 1997, ApJ, 482, L97
Cox, P., Boulanger, F., Huggins, P. J., et al. 1998, ApJ, 495, L23
Dinerstein, H. L., Sneden, C., & Uglum J. 1995, ApJ, 447, 262
Forveille, T., Huggins, P. J., Bachiller, R., & Cox, P. 1998, ApJ, 495, L111
Hora, J. L., & Latter, W. B. 1994, ApJ, 437, 281
Hora, J. L., & Latter, W. B. 1999, ApJ, in press
Howe, D. A., & Williams, D. A. 1998, in The Molecular Astrophysics of Stars and Galaxies, ed. T. W. Hartquist & D. A. Williams (Oxford: Oxf. Univ. Press), 347
Huggins, P. J. 1999, in IAU Symp. 191, Asymptotic Giant Branch Stars, ed. T. Le Bertre, A. Lèbre, C. Waelkens (San Francisco: PASP), 425
Huggins, P. J., Bachiller, R., Cox, P., & Forveille, T. 1992, ApJ, 401, L43
Huggins, P. J., Bachiller, R., Cox, P., & Forveille, T. 1996, A&A, 315, 284
Josselin, E., Bachiller, R., Manchado, A., & Guerrero, M. 1999, A&A, in press
Kastner, J. H., Weintraub, D. A., Gatley, I., Merrill, K. M., & Probst, R. G. 1996, ApJ, 462, 777
Lucas R., & Guélin M., 1999, in IAU Symp. 191, Asymptotic Giant Branch Stars, ed. T. Le Bertre, A. Lèbre, C. Waelkens (San Francisco: PASP), 305
Mauron, N., & Huggins, P. J. 1999, A&A, 349, 203
Meaburn, J., Clayton, C. A., Bryce, M., Walsh, J. R., Holloway, A. J., & Steffen, W. 1998, MNRAS, 294, 201
Natta, A., & Hollenbach, D. 1998, A&A, 337, 517
Neri, R., Kahane, C., Lucas, R., Bujarrabal, V., & Loup, C. 1998, A&AS, 130, 1
O’Dell, C. R., & Burkert, A. 1997, in IAU Symp. 180, Planetary Nebulae, ed. H. J. Habing & H. J. G. L. M. Lamers (Dordrecht: Kluwer), 332
Olofsson, H. 1999, in IAU Symp. 191, Asymptotic Giant Branch Stars, ed. T. Le Bertre, A. Lèbre, C. Waelkens (San Francisco: PASP), 3
Taylor, A. R., Gussie, G. T., & Pottasch, S. R. 1990, ApJ, 351, 515
Young, K., Cox, P., Huggins, P. J., Forveille, T., & Bachiller, R. 1997, ApJ, 482, L101
Young, K., Cox, P., Huggins, P. J., Forveille, T., & Bachiller, R. 1999, ApJ, 522, 387
Zweigle, J., Neri, R., Bachiller, R., Bujarrabal, V., & Grewing, M. 1997, A&A, 324, 624 |
warning/0001/cond-mat0001351.html | ar5iv | text | # The sign problem in Monte Carlo simulations of frustrated quantum spin systems
## I INTRODUCTION
Recently, there have been several significant developments of more efficient Monte Carlo methods for interacting quantum many-body systems. The Trotter decomposition formula has traditionally been used as a starting point for finite-temperature simulation algorithms, such as the worldline and fermion determinant methods. It introduces a systematic error that can be removed only by carrying out simulations for several different imaginary time discretizations $`\mathrm{\Delta }\tau `$ and subsequently extrapolating to $`\mathrm{\Delta }\tau =0`$. Such extrapolations are not necessary with the stochastic series expansion method, which is based on sampling the power series expansion of exp$`(\beta H)`$ to all orders and is related to a less general method proposed much earlier by Handscomb. Results that are exact to within statistical errors can also be directly obtained with recent worldline and fermion determinant algorithms formulated in continuous imaginary time. Even more significant than the elimination of the Trotter decomposition error are generalizations to the quantum case of cluster algorithms developed for classical Monte Carlo. These “loop algorithms” (so called because the clusters are loops on a space-time lattice) can reduce the autocorrelation times by orders of magnitude and enable highly accurate studies of systems in parameter regimes where previous algorithms encountered difficulties due to long autocorrelation and equilibration times.
In spite of these developments, the class of models which can be studied using quantum Monte Carlo methods is still severely restricted due to the “sign problem”, i.e., the non-positive-definiteness of the weight function that can arise in transforming a quantum problem into a form resembling a classical statistical mechanics problem. There are two classes of systems for which this issue is particularly pressing — interacting fermions in more than one dimension and quantum spin systems with frustrated interactions (in any number of dimensions). For fermions in one dimension, and hopping between nearest-neighbor sites only, the sign problem can be avoided because the fermion anticommutation relations do not come into play (other than introducing a hard-core constraint) in the one-dimensional real-space path integral. In two or more dimensions (or even in one dimension if hopping further than between nearest neighbors is included), permutations of fermions during the propagation in imaginary time lead to a mixed-sign path integral which typically cannot be efficiently evaluated using Monte Carlo methods. The sign problem can be avoided with the fermion determinant algorithm in special cases, such as the half-filled Hubbard model (because of particle-hole symmetry), but in other cases simulations are restricted to high temperatures and/or small system sizes. For frustrated spin systems the source of the sign problem is different. A minus sign appears for every event in the path integral in which two antiferromagnetically interacting spins are flipped. This causes an over-all minus sign if the total number of spin flips is odd, which can be the case, e.g., for a triangular lattice or a square lattice with both nearest- and next-nearest-neighbor interactions. Simulations of quantum spin systems are therefore restricted to models with no frustration (in the off-diagonal part of the Hamiltonian), such as ferromagnets, or antiferromagnets on bipartite lattices.
A promising approach to solving the sign problem was recently suggested by Chandrasekharan and Wiese. They considered a system of spinless fermions on a two-dimensional square lattice within the context of the worldline loop algorithm. They showed that, for this particular model and for a certain range of nearest-neighbor repulsion strengths, the properties of the loops can be used to eliminate the sign problem. Flipping a loop can change the number of fermion permutations from odd to even, or vice versa, thereby also changing the over-all sign of the configuration. Such sign-changing loops are called “merons”. The magnitude of the weight is not affected by flipping a meron and therefore all configurations with one or more merons cancel in the partition function. The subspace of zero merons is positive definite and can be sampled without a sign problem. Typical operator expectation values of interest also contain contributions from configurations with two merons which therefore also have to be included in the simulation and introduce a “mild” sign problem. The relative weights of the zero- and two-meron subspaces to be sampled can further be chosen in an optimum way using a reweighting technique, which further reduces the sign problem.
In this paper we explore an analogous method for solving the sign problem for frustrated quantum spin models. We consider the operator-loop formulation of the stochastic series expansion method, in which sequences of two-spin operators are sampled by forming clusters (loops) of operators that can be simultaneously updated without changes in the weight function. The updated clusters contain operators acting on the same spins, but diagonal operators may be changed to off-diagonal ones and vice versa. For a model with frustrated interactions an operator-loop update can lead to a sign change. In analogy with Chandrasekharan and Wiese we will refer to such sign-changing loops as “merons”. The sign problem can be solved if the operator-loops for a given configuration can be uniquely defined and the weight function is positive definite in the configuration subspace containing no merons. Unfortunately, we find that these criteria are in general difficult to satisfy. Operator-loop algorithms with uniquely determined loops are typically non-ergodic for frustrated systems, and with supplemental local updates for ergodicity there are mixed signs in the zero-meron subspace. In fact, in such cases merons typically do not even exist, i.e., none of the operator-loops can change the sign when flipped. We have found only one spin system for which the sign problem can be eliminated using merons — the Heisenberg model with ferromagnetic couplings $`J_z(r)<0`$ along the $`z`$-axis and frustrated antiferromagnetic couplings $`J_{xy}(r)=J_z(r)`$ in the plane perpendicular to this axis, i.e., the Hamiltonian
$$H=\underset{i,j}{}J_{ij}[S_i^zS_j^z\frac{1}{2}(S_i^+S_j^{}+S_i^{}S_j^+)],$$
(1)
where $`J_{ij}>0`$ and the range of the couplings is arbitrary. We have implemented a meron algorithm for this model on a square lattice with nearest- and next-nearest-neighbor couplings $`J(1)`$ and $`J(\sqrt{2})`$. Standard algorithms for this model have a severe sign problem when using the $`z`$-axis as the quantization axis, however, it can be avoided by a simple rotation to the $`x`$-representation. Using the SSE algorithm and the meron concept, the sign problem can be eliminated also in the $`z`$ representation. With both representations accessible in simulations, correlation functions both parallel and perpendicular to the $`z`$ direction can be easily evaluated.
The model, Eq. (1), can be mapped onto a hard-core boson model with attractive interactions and frustrated hopping. Frustration in the potential energy has been investigated in this context as a possible mechanism to render a disordered bosonic ground state. Frustration in the hopping \[the $`xy`$-term in Eq. (1)\] should decrease the tendency to forming off-diagonal long-range order and could then lead to a normal fluid (non-superfluid). However, the highly symmetric case considered here has a trivial, ordered ground state; the fully polarized ferromagnetic state (corresponding to a completely filled lattice of hard-core bosons; a trivial case of normal solid). Effects of frustration only come into play at finite temperature, where the model is different from the corresponding isotropic ferromagnet (on non frustrated, bipartite lattices the two models are equivalent, since the sign of the $`xy`$-term can be switched by a spin rotation on one of the sublattices).
Although we have not been able to solve the sign problem for other cases, such as the Heisenberg model with completely antiferromagnetic interactions \[$`J_z(1)=J_{xy}(1)>0`$ and $`J_z(\sqrt{2})=J_{xy}(\sqrt{2})>0`$\], our work nevertheless gives some further insights into the meron concept and what is required in order to solve the sign problem for arbitrary couplings.
The outline of the rest of the paper is the following: In Sec. II we review the basics of the stochastic series expansion method and discuss operator-loop updating schemes for both ferromagnetic and antiferromagnetic couplings. In Sec. III we present the solution of the sign problem for the $`J_z(r)=J_{xy}(r)`$ model. The reweighting technique is analyzed in some detail in Sec. IV. In Sec. V we discuss some simulation results for the semifrustrated model and make comparisons with the isotropic Heisenberg ferromagnet. We summarize our work in Sec. VI.
## II OPERATOR-LOOP ALGORITHM
In this section we first briefly review the expansion underlying the SSE method and then discuss the operator-loop updates used to efficiently sample the expansion. We here assume a non-frustrated case and postpone the discussion of the sign problem for frustrated models to Sec. III. For definiteness we consider the $`S=1/2`$ Heisenberg model
$$H=\pm J\underset{i,j}{}S_iS_j,$$
(2)
where $`ij`$ denotes a pair of nearest-neighbor spins on a cubic lattice (in an arbitrary number of dimensions), and $`J>0`$. Depending on the sign, the model is an antiferromagnet ($``$) or a ferromagnet $`(+)`$. To construct the SSE configuration space the Hamiltonian is rewritten as a sum of diagonal and off-diagonal operators
$$H=\frac{J}{2}\underset{b=1}{\overset{M}{}}\left(H_{1,b}H_{2,b}\right)+C,$$
(3)
where the index $`b`$ denotes an interacting spin pair (bond) $`i(b),j(b)`$ and $`C`$ is an irrelevant constant equal to $`MJ/4`$, where $`M`$ is the total number of pairs of interacting spins. The bond-indexed operators are given by
$`H_{1,b}`$ $`=`$ $`2\left(\frac{1}{4}S_{i(b)}^zS_{j(b)}^z\right)`$ (4)
$`H_{2,b}`$ $`=`$ $`S_{i(b)}^+S_{j(b)}^{}+S_{i(b)}^{}S_{j(b)}^+.`$ (5)
Note that the eigenvalues of both the diagonal ($`H_{1,b}`$) and the off-diagonal ($`H_{2,b}`$) operators are $`0`$ and $`1`$, both for the antiferromagnet and the ferromagnet. The partition function $`Z=\text{Tr}\{\mathrm{exp}(\beta H)\}`$ is expanded as
$$Z=\underset{\alpha }{}\underset{n=0}{\overset{\mathrm{}}{}}\frac{(\beta )^n}{n!}\alpha |H^n|\alpha ,$$
(6)
in the basis $`\{|\alpha \}=\{|S_1^z,S_2^z,\mathrm{},S_z^N\}`$, where $`N`$ is the number of spins. Terms of order greater than $`nN\beta `$ give an exponentially vanishing contribution and for the purpose of a stochastic sampling the expansion can therefore be truncated at some $`n=L`$ of this order without loss of accuracy (see, e.g., Ref. for details on how to choose a sufficiently high truncation power). Additional unit operators $`H_{0,0}1`$ are introduced to rewrite Eq. (6) as
$$Z=\underset{\alpha }{}\underset{S_L}{}\frac{(1)^{n_2}(J\beta )^n(Ln)!}{2^nL!}\alpha \left|\underset{i=1}{\overset{L}{}}H_{a_i,b_i}\right|\alpha ,$$
(7)
where $`S_L`$ denotes a sequence of operator-indices:
$$S_L=(a_1,b_1)_1,(a_2,b_2)_2,\mathrm{},(a_L,b_L)_L,$$
(8)
with $`a_i\{1,2\}`$ and $`b_i\{1,\mathrm{},M\}`$, or $`(a_i,b_i)=(0,0)`$. The number of non-(0,0) elements in $`S_L`$ is denoted $`n`$, while $`n_2`$ denotes the number of off-diagonal operators in the sequence. Note that since the expectation value in Eq. (7) is always equal to zero or one, the sign of a term is negative only if $`n_2`$ is odd. This sign problem occurs (only) when frustration is present and is the main topic of this paper. However, for the discussion of the sampling procedures, in this section we assume a positive definite expansion. We introduce the notation $`|\alpha (p)`$ for a propagated state
$$|\alpha (p)=\underset{i=1}{\overset{p}{}}H_{a_i,b_i}|\alpha ,$$
(9)
where for an allowed configuration $`|\alpha (0)=|\alpha (L)=|\alpha `$ and the weight function corresponding to (7) is given by
$$W(\alpha ,S_L)=\frac{(J\beta )^n(Ln)!}{2^nL!}.$$
(10)
Having established the framework we will proceed to describe the procedures for importance sampling of the terms $`(\alpha ,S_L)`$ according to the weight (10). The initial state can be a sequence of the form $`(0,0)_1,(0,0)_2,\mathrm{},(0,0)_L`$ (subscripts on the index pairs will sometimes be used to denote the position in $`S_L`$) and a random state $`|\alpha `$. An ergodic procedure for sampling the terms is achieved using two types of basic updates; a simple substitution of single diagonal operators and the operator-loop update which involves simultaneous updates of a number (in principle varying between $`1`$ and $`n`$) of diagonal and off-diagonal operators.
The diagonal update is carried out by traversing the operator-index sequence $`S_L`$ from beginning ($`p=1`$) to end $`(p=L)`$. Operator substitutions of the form $`(0,0)_p(1,b)_p`$ are attempted where possible, while the stored state $`|\alpha `$ is updated every time an off-diagonal operator is encountered so that the state $`|\alpha (p)`$ is available when needed. With the weight function (10) detailed balance can be seen to be satisfied if the acceptance probabilities are taken to be
$`P\left[(0,0)_p(1,b)_p\right]`$ $`=`$ $`{\displaystyle \frac{M\beta \alpha _b(p)|H_{1,b}|\alpha _b(p)}{Ln}},`$ (11)
$`P\left[(1,b)_p(0,0)_p\right]`$ $`=`$ $`{\displaystyle \frac{Ln+1}{M\beta \alpha _b(p)|H_{1,b}|\alpha _b(p)}}.`$ (12)
Note that the diagonal update changes the expansion power $`n`$ by $`\pm 1`$. Off-diagonal operators cannot be introduced one-by-one because of the periodicity condition $`|\alpha (L)=|\alpha (0)`$. Local updates involving two operators can be used for this purpose, but are more complicated and far less efficient than the operator-loop update, which is discussed next.
We use a largely pictorial description of the operator loops. First we consider the antiferromagnet. Note that in this case the only non-zero matrix elements of the bond operators are
$`|H_1|`$ $`=`$ $`|H_1|=1,`$ (13)
$`|H_2|`$ $`=`$ $`|H_2|=1,`$ (14)
i.e., they can act only on antiparallel spins. An example of a term in the expansion for a four-site antiferromagnet is depicted in Fig. (1). This representation makes evident the close relationship between the SSE expansion and the Euclidean path integral. An imaginary time separation $`\tau `$ corresponds to a distribution of propagations $`\mathrm{\Delta }p`$ between states, centered around $`\mathrm{\Delta }p=(\tau /\beta )n`$. We will for convenience here refer to the propagation as the time dimension.
The general idea behind the loop update is to flip a cluster of spins in the configuration in such a way that the weight is not changed. With the SSE method there will also have to be changes made to the operators acting on the spins, since otherwise operators $`H_1`$ or $`H_2`$ may act on parallel spins, resulting in zero-valued matrix elements. Since one of the states $`|\alpha (p)`$ and the operator sequence $`S_L`$ uniquely define the whole spin configuration, the SSE loops are in practice treated as loops of operators, the exact meaning of which will be made clear below.
Consider one of the operators $`H_{1,1}`$ in Fig. (1). It can be depicted as a “vertex” with four legs associated with spin states $``$ or $``$. If we flip the upper left spin, a vanishing matrix element results. But if we flip both upper (or lower) spins and also change the operator type to off-diagonal $`H_{2,1}`$, an allowed matrix element is generated; see Fig. (2). Using this idea we can form a cluster of spins by choosing a random spin $`S_i^z(p)`$ in the configuration and traversing up or down until we encounter an operator (bond) acting on that spin, then switch to the second spin of the bond and change the direction of traversing the list. Eventually we will necessarily arrive back at the initial starting point, whereupon a closed loop has formed. All the spins on this loop can be flipped if the operators encountered are also switched \[$`(1,b)(2,b)]`$. Note that the same operator can be encountered twice, which results in no net change of operator type (but the spins at all four vertex legs are flipped).
The whole configuration can be uniquely divided up into a set of loops, so that each spin belongs to one and only one cluster; see Fig. (3), where our example configuration has been divided up into three clusters. All loops can be flipped independently with probability $`1/2`$ — in Fig. (3) we depict the result of flipping the largest loop of the example configuration. A full operator-loop update amounts to dividing up the configuration into all of its loops which are flipped with probability $`1/2`$. The (random) decision of whether or not to flip a loop can be made before the loop is constructed, so that each loop has to be traversed only once. Spin “lines” $`S_i^z(p)`$, $`p=0,\mathrm{},n`$, which are not acted upon by any of the operators in $`S_L`$ will not be included in any of the operator-loops. They correspond to “free” spins which can be flipped with probability $`1/2`$. Such a line can also be considered a loop, and then it will always be true that every spin $`S_i^z(p)`$ belongs to one loop. Free spins appear frequently at high temperatures, when the total number of operators $`nN`$, but are rare at low temperatures.
Note that the spin states at the four legs of the operator-vertices completely determine the full spin configuration, except for free spins that happen not to be acted upon by any of the operators in $`S_L`$. Hence, the operator-loop update can be carried out using only a linked list of of operators, i.e., an array of vertices with four spin states and associated pointers to the “previous” and “next” vertex associated with the same spin. The storage requirements and the number of operations needed for carrying out a full operator-loop update then scale as $`N\beta `$ instead of $`N^2\beta `$ if the full spin configuration were to be used.
In a simulation we first make a full cycle of diagonal updates in the sequence $`S_L`$ and then create the linked list of vertices in which the operator-loop updates are carried out. The vertex list is then mapped back onto the sequence $`S_L`$ and the state $`|\alpha (0)`$. Alternatively, the linked list can be updated simultaneously with each diagonal operator substitution, so that it does not have to be re-created for each Monte Carlo step — depending on the model studied there may be significant differences in execution time between the two approaches.
For the ferromagnet we can construct the loops in a similar manner, but the non-zero matrix elements are now
$`|H_1|`$ $`=`$ $`|H_1|=1,`$ (15)
$`|H_2|`$ $`=`$ $`|H_2|=1,`$ (16)
i.e., the off-diagonal operators act only on antiparallel spins, as before, whereas the diagonal ones can act only on parallel configurations. This implies qualitative changes in the structure of the loops, as depicted in Fig. (4). If we again consider an operator $`H_1`$ and flip the upper left spin, we note that we need to flip the lower right spin and change the operator to $`H_2`$. Hence, instead of changing the direction of traversing the configuration every time an operator is encountered we now continue in the same direction after switching to the second spin of the bond. Any configuration can still be uniquely divided up into loops that can be flipped with probability $`1/2`$. An example of a ferromagnetic configuration with its loops is shown in Fig. (5).
Note that since the loops for the ferromagnet never change direction as they go through the lattice, every single loop has to traverse each state $`|\alpha (p)`$ at least once. It follows that the number of sites $`N`$ is an upper bound of the number of loops. The antiferromagnetic loops, on the other hand, traverse the lattice in both directions and the number of loops is therefore instead limited by the total number of operators $`nN\beta `$. As a consequence of the change of direction, for the antiferromagnet the linked list of vertices must to be bi-directional, but for the ferromagnet it is sufficient to keep a singly-directional list.
The diagonal and operator-loop updates satisfy detailed balance and the combination of them leads to ergodic sampling for a ferromagnet on any lattice, and for antiferromagnets on bipartite lattices — for frustrated antiferromagnets there are complications, in addition to the sign problem, which will be discussed further in the next section. The operator-loop sampling is highly efficient, with integrated auto-correlation times typically less than one Monte Carlo step.
The loop construction described here relies on the rotational invariance of the models, i.e., the fact that both the diagonal and off-diagonal matrix elements in Eqs. (14) and (16) are equal to one. For a general anisotropic case the loops will lead to weight changes when flipped and must then be assigned “a-posteriori” acceptance probabilities which typically become small for large lattices at low temperatures. Other types of loops avoiding this problem have been proposed, but will not be discussed here.
## III THE SIGN PROBLEM
The notorious sign problem arises in stochastic sampling when the function used to weight the different configurations in not positive definite. A typical quantity that can be calculated by Monte Carlo (importance sampling) is of the form
$$A=\frac{_iA(x_i)W(x_i)}{_iW(x_i)}=A(x)_W,$$
(17)
where $`W`$ is the weight function and $`A`$ the measured quantity, which both depend on the general coordinate $`x`$ of the configuration space sampled. When the coordinates are sampled according to relative weight, the desired quantity is simply the arithmetic average of the measured function $`A(x)`$, as indicated by the notation $`A(x)_W`$ above. If the weight function is not positive definite, the sampling can be done using the absolute value of the weight, and the expectation value can be calculated according to
$$A=\frac{A(x)\text{s}(x)_{|W|}}{s(x)_{|W|}},$$
(18)
where $`s(x)`$ equals $`\pm 1`$, depending on whether the sign of the weight function is positive or negative. For most models, where a sign problem is present, the average sign $`s(x)_{|W|}`$ decreases exponentially to zero as a function of inverse temperature and system size, and the relative statistical errors of calculated quantities increase exponentially.
A meron-cluster solution to the sign problem using loop updates of fermion world-line configurations was recently proposed by Chandrasekharan and Wiese. This approach is based on the idea that if it is possible to map every configuration with negative weight uniquely to a corresponding configuration with equal weight magnitude but opposite sign, then the partition function can be sampled without a sign problem simply by not including any configuration which is a member of such a canceling pair. In the meron-cluster algorithm, flipping a loop of spins can lead to a sign change without change in the magnitude of the weight, and such a “meron” hence identifies a canceling pair of configurations. Here we will present a similar approach within the SSE operator-loop method for frustrated quantum spins.
Let us consider the Heisenberg antiferromagnet discussed in the previous section. From Eq. (7) we see that a configuration has negative weight if the total number $`n_2`$ of off-diagonal operators is odd. This can only be the case on frustrated lattices. As described in the previous section, any SSE configuration can be divided up uniquely into a number of loops. Flipping a loop interchanges the diagonal and off-diagonal operators, but leaves the total number of operators unchanged. It follows that the sign will change if and only if a loop passing through an odd number of operators is flipped (two passes through the same operator is counted as two operators). Since the total weight remains unchanged we have thus found the desired mapping between positive and negative configurations (assuming that there exist loops which change the sign when flipped, which in fact is not always the case). In analogy with previous work we call such a sign-flipping loop a “meron”. Let us now see how we can use this concept to measure observables.
As in Ref. we consider improved estimators that average over all loop configurations. Denote the number of loops in the system $`N_L`$. Since each loop can be in one of two states there is a total of $`2^{N_L}`$ configurations that can be reached by flipping all combinations of the loops present. The improved estimate therefore takes the form
$$A=\frac{A(x)s(x)_{|W|}}{s(x)_{|W|}},$$
(19)
where the double brackets denote an average over all the loop states for each generated SSE configuration, e.g.,
$$s(x)=\frac{1}{2^{N_L}}\underset{l=1}{\overset{2^{N_L}}{}}s(x_l).$$
(20)
The general coordinate $`x`$ here refers to the SSE configuration space $`(\alpha ,S_L)`$ and $`x_l`$ refers to one out of the $`2^{N_L}`$ possible outcomes of “flipping” a number of loops.
Let us consider this average. Denote the state of a loop with $`\delta `$, with two possible “orientations” $`\delta \{,\}`$. Since flipping one loop does not affect any other loops (in terms of their paths taken), the sign of a configuration factors according to
$$s(\delta _1,\delta _2,\mathrm{},\delta _{N_L})=\underset{i=1}{\overset{N_L}{}}s(\delta _i),$$
(21)
where $`s(\delta )=\pm s(\overline{\delta })`$, where $`\overline{\delta }`$ denotes a flipped loop, and the sign is negative for merons and positive otherwise. Since flipping any meron leads to two terms that cancel, it follows that
$$\frac{1}{2^{N_L}}\underset{l=1}{\overset{2^{N_L}}{}}s(x_l)=\pm \delta _{n_M,0},$$
(22)
where $`n_M`$ denotes the total number of merons. The sign in front of the delta function is the “inherent” sign of the configuration, independent of the loop orientation when there are no merons present. This sign has to be positive for the meron solution to be applicable in practice, and then the partition function can be sampled in the positive definite subspace of configurations with no merons.
Having found an expression for the denominator in Eq. (19) we need to consider the numerator for cases of interest. SSE estimators for a number of important operators have been discussed, e.g., in Ref. . The internal energy is given by
$$E=\frac{1}{\beta }n_W,$$
(23)
where $`n`$ denotes the total number of operators in the sequence $`S_L`$. This number is not affected by the loop updates, and hence it follows that
$$E=\frac{1}{\beta }\frac{n(x)s(x)_{|W|}}{s(x)_{|W|}}=\frac{1}{\beta }\frac{n(x)\delta _{n_M,0}_{|W|}}{\delta _{n_M,0}_{|W|}}.$$
(24)
Assuming that this sector has positive definite weight we have therefore completely eliminated the sign problem by restricting the simulation to the zero-meron sector. The energy is then simply given by
$$E=\frac{n(x^0)_W}{\beta },$$
(25)
where the superscript 0 indicates the restriction of the simulation to the zero-meron sector.
Next we will consider the magnetic susceptibility,
$$\chi =\beta \left(\underset{i}{}S_i^z\right)^2/N=\beta M^2/N.$$
(26)
Since $`M`$ is conserved by the Hamiltonian its value is the same in all propagated states; $`M(p)=M(0)M`$ $`(p=1,\mathrm{},n)`$. In a configuration uniquely divided up into loops, every spin $`S_i^z(p)`$ belongs to one and only one loop, if we count as loops also all “lines” of free spins, i.e., the spins $`S_i^z(p)`$, $`p=1,\mathrm{},n`$ for all sites $`i`$ which are not associated with any operator in the sequence (and therefore can be flipped). It follows that the change in $`M(p)`$ when flipping a loop must be the same for all $`p`$, and hence only loops that go through all states $`|\alpha (p)`$ (at least once) can change $`M`$ when flipped. We can therefore introduce a loop magnetization $`m_L`$, which is simply equal to the sum of the spins traversed by the loop for an arbitrary $`|\alpha (p)`$. In the estimator (19) corresponding to the susceptibility the numerator can hence be written as
$$\frac{1}{2^{N_L}}\underset{l=1}{\overset{2^{N_L}}{}}(m_1(x_l)+\mathrm{}+m_{N_L}(x_l))^2s(x_l),$$
(27)
The magnetization on a loop always changes sign when a loop is flipped; the over-all sign $`s(x_l)`$ only changes sign when a meron is flipped. Therefore, in summing over all loops in (27), a non-zero value results only if the configuration has zero or two merons. The full susceptibility estimator therefore takes the form
$$\chi =\frac{\underset{l=1}{\overset{2^{N_L}}{}}|m_l|^2\delta _{n_M,0}+2|m_{M_1}||m_{M_2}|\delta _{n_M,2}}{\delta _{n_M,0}},$$
(28)
where $`M_1`$ and $`M_2`$ are the indices of the loops corresponding to merons in a two-meron configuration. Hence, unlike in the case of the total energy, the sign problem has here not been completely eliminated, since the zero- and two-meron configuration contribute $`1`$ and $`0`$, respectively, to the average sign. When the SSE configuration volume grows the relative weight of the zero-meron sector should diminish, leading to a decreasing average sign. Chandrasekharan and Wiese stated that the statistical fluctuation in the improved susceptibility estimator increases quadratically with $`N\beta `$, i.e., much slower than the conventional exponential increase. They also pointed out that this remaining sign problem can be solved by reweighting the zero- and two-meron sectors with external weight factors $`w(0)`$ and $`w(2)`$. This changes the above formula to
$$\chi =\frac{_{l=1}^{2^{N_L}}|m_l|^2\delta _{n_M,0}w(2)+2|m_{M_1}||m_{M_2}|\delta _{n_m,2}w(0)}{\delta _{n_M,0}w(2)}$$
(29)
In the next section we will say more about reweighting.
As we have shown above, the meron concept within the SSE method formally leads to exactly the same equations as in the world-line simulations of fermion systems considered in Ref. . The difference is only in the structure of the meron itself; the fermionic meron changes the number of particle permutations from even to odd, or vice versa, whereas the SSE meron in the case of a frustrated spin system instead changes the number of antiferromagnetic spin flips from even to odd or vice versa. Applying the SSE operator-loop algorithm to a fermion system would lead to merons of exactly the same kind as those existing within the world-line framework, and, conversely, applying a world-line loop algorithm to a frustrated spin system should lead to merons similar to those discussed here (there are no diagonal operators in the world-line configurations, but spin flip events correspond to the SSE off-diagonal operators and can change from even to odd, or vice versa, in loop updates). These similarities are not surprising, considering the close relationship between the SSE expansion and the Euclidean path integral.
Now consider the application of the above ideas to the Heisenberg model on a square lattice with nearest- and next-nearest-neighbor couplings $`J(1)>0`$ and $`J(\sqrt{2})>0`$ (antiferromagnetic). This model has a sign problem since the total number $`n_2`$ of spin-flipping operators in a configuration can be odd, e.g., three operators on a triangle of spins, as shown in Fig. (6). Already this simple example illustrates that the operator-loop algorithm discussed above does not sample the full configuration space and that the meron concept therefore cannot be applied to solving the sign problem. Since all three bonds are antiferromagnetic, a loop will change direction every time an operator is encountered. In order for the loop to close, it therefore has to pass through an even number of operators and hence flipping the loop cannot change the number $`n_2`$ of off-diagonal operators from odd to even, or vice versa. This is illustrated in Fig. (6), where the only effect of flipping the single loop in the system is to flip all the spins; the operators remain unchanged. Hence, merons do not even exist within the operator-loop algorithm for this model, and the sampling is non-ergodic. A local update can in principle be used in combination with the operator-loops in order to make the sampling ergodic. However, a configuration can then have a negative sign (of which Fig. (6) is an example) even though there are no merons present. The principal requirement of positive-definiteness in the zero-meron subspace (which in this case is the full space) is hence not fulfilled. A similar problem seems to affect all models with frustration in all of the spin components. It is possible that some other way of constructing the loops could remedy this, e.g., by switching to some other, non-trivial basis in which the SSE expansion could be carried out. Other ways proposed for constructing loops in the standard basis considered here do not uniquely define the set of loops and can therefore not easily be used with the meron concept.
We have found one class of spin models for which the meron ideas can be successfully applied to solve the sign problem: Heisenberg models which are antiferromagnetic in the $`xy`$-plane but ferromagnetic along the $`z`$-axis., i.e., the “semi-frustrated” model (1), which can be written as
$$H=\underset{b=1}{\overset{M}{}}\frac{J_b}{2}\left(H_{1,b}H_{2,b}\right)+C,$$
(30)
where the bond-indexed operators are given by
$`H_{1,b}`$ $`=`$ $`2\left(\frac{1}{4}+S_{i(b)}^zS_{j(b)}^z\right),`$ (31)
$`H_{2,b}`$ $`=`$ $`S_{i(b)}^+S_{j(b)}^{}+S_{i(b)}^{}S_{j(b)}^+.`$ (32)
On a non-frustrated lattice this model is equivalent to an isotropic Heisenberg ferromagnet, since $`n_2`$ is always even and the sign in front of the operators $`H_{2,b}`$ in Eq. (30) is irrelevant as $`(1)^{n_2}=1`$ in Eq. (7) — the sign can also be transformed away by a spin rotation on one of the two sublattices. On a frustrated lattice, on the other hand, $`n_2`$ can be odd (the lattice is no longer bipartite so that the transformation mentioned above does not remove all signs), and the system is no longer equivalent to the isotropic ferromagnet. The model has a classical two-fold degenerate ferromagnetic ground state, but at finite temperatures the transverse spin components are frustrated, and the behavior will be different from the isotropic ferromagnet. When simulated in the $`z`$-basis using standard algorithms the semifrustrated model has a severe sign problem, but the zero-meron sector is positive-definite and the SSE meron solution can be applied. The structure of the operator-loops is the same as for the ferromagnet and the loop algorithm is therefore ergodic.
The meron solution can be implemented in several different ways and we briefly describe how it was done in this work: During the sequential diagonal updates the linked list, representing the loop structure, is updated simultaneously with each accepted diagonal update. The loops are numbered and information is stored on whether each loop is a meron or not. During an attempted diagonal update only the loops directly affected by the operator substitution are updated. This permits easy and fast checking of whether the number of merons in the system has changed or not. If the new number of merons is different form zero or two the update is rejected, whereas if the number of merons changes from $`i`$ to $`j`$ it is accepted with probability $`w(j)/w(i)`$, where $`i,j\{0,2\}`$, and $`w(i)`$ is the re-weighting factor assigned to meron sector $`i`$.
Note that the sign problem for the semifrustrated model can also be very easily transformed away by rotating the ferromagnetic component to the $`y`$-direction. The Hamiltonian then takes the form
$$H=\frac{J}{2}\underset{b=1}{\overset{M}{}}\left(H_{1,b}^{}+H_{2,b}^{}\right)+C,$$
(33)
where the bond-indexed operators are given by
$`H_{1,b}^{}`$ $`=`$ $`2\left(\frac{1}{4}S_{i(b)}^zS_{j(b)}^z\right),`$ (34)
$`H_{2,b}^{}`$ $`=`$ $`S_{i(b)}^+S_{j(b)}^++S_{i(b)}^{}S_{j(b)}^{},`$ (35)
and the fundamental spin-flips and operator exchange during a loop update is shown in Fig. (7). Being able to work in both bases we can easily measure all components of the susceptibility. Our main motivation for studying this model is to illustrate how the sign problem can be removed in the $`z`$-basis. Nevertheless, we will also show some results calculated in the $`x`$-basis.
## IV REWEIGHTING
An important technical aspect of the meron-solution is the reweighting of the zero- and two-meron sectors, which was briefly mentioned in the previous section. Eq. (29) gives the correct estimator for the susceptibility after reweighting, but it gives no information on how to do the reweighting in practice. How to determine the optimal reweighting and whether reweighting changes the scaling of the relative error are important questions to be considered in this section.
As a first, simple example of how reweighting affects the statistics of a simulation we will discuss a simple random process. Consider a random variable $`n`$, which can take two different values, $`0`$ or $`2`$. Let $`W_0`$ and $`W_2`$ designate the probability for these two outcomes. The expected fractions and standard deviations of these outcomes from $`N`$ random selections are given by
$`\delta _{n,0}`$ $`=`$ $`W_0\pm {\displaystyle \frac{\sqrt{W_0W_2}}{\sqrt{N}}}`$ (36)
$`\delta _{n,2}`$ $`=`$ $`W_2\pm {\displaystyle \frac{\sqrt{W_0W_2}}{\sqrt{N}}}.`$ (37)
We consider an expectation value of a form similar to the QMC susceptibility, Eq. (28);
$$f=\frac{\delta _{n,0}+\delta _{n,2}}{\delta _{n,0}}=\frac{1}{W_0}\pm \frac{1}{W_0}\sqrt{\frac{W_2}{W_0}}\frac{1}{\sqrt{N}},$$
(38)
with a relative standard deviation
$$\frac{\sigma _f}{f}=\sqrt{\frac{W_2}{W_0}}\frac{1}{\sqrt{N}}.$$
(39)
This formula becomes valid for large $`N`$, when the standard deviation is small. As $`W_0`$ decreases the standard deviation increases, but one can reweight the two outcomes by assigning an additional weight $`W`$ to the $`n=0`$ outcome such that a transition from $`n=0`$ to $`n=2`$ is accepted with probability $`1/W`$, while a transition from $`n=2`$ to $`n=0`$ is always accepted. After such a reweighting the probabilities of obtaining $`n=0`$ and $`n=2`$ are given by
$`W_0^{}`$ $`=`$ $`{\displaystyle \frac{W_0W}{W_0W+W_2}}`$ (40)
$`W_2^{}`$ $`=`$ $`{\displaystyle \frac{W_2}{W_0W+W_2}},`$ (41)
and $`f`$ is given by
$$f=\frac{\delta _{n,0}+W\delta _{n,2}}{\delta _{n,0}}.$$
(42)
When calculating the standard deviation for this case we have to be careful since the reweighting introduces correlations into the system. This is clearly visualized in Fig. 8, where in the upper graph a series of independent outcomes with equal probability ($`W_0=W_2=0.5`$) are shown , while in the lower graph a case with $`W_0=0.01`$ is shown with a reweighting factor of $`W=99`$ (leading to $`W_0^{}=W_2^{}=0.5`$).
Let us now calculate the standard deviations for this case. We are interested in minimizing the standard deviation of a variable evaluated in a simulation and also need to calculate its statistical error. In a standard MC simulation one usually wants to calculate the average and the standard deviation of the average for some quantity $`x`$. This is typically achieved by dividing the run into a number of bins, $`N`$, and saving the average of $`x`$ for each bin. If the bins are statistically independent the final average and standard deviation can be calculated according to
$$\overline{x}=\frac{1}{N}\underset{i=1}{\overset{N}{}}x_i$$
(43)
and
$$\sigma _{\overline{x}}=\sqrt{\frac{\overline{x^2}\overline{x}^2}{N}}.$$
(44)
When studying the behavior of the standard deviation itself, we also want to obtain an estimate of its accuracy. This can be done by dividing the $`N`$ bins into $`M`$ sets containing $`N/M`$ bins each. For each set a standard deviation $`\sigma _x`$ can be calculated according to
$$\sigma _x=\sqrt{\overline{x^2}\overline{x}^2},$$
(45)
where the bar denotes an average of the $`N/M`$ bins within the set. The final standard deviation and its statistical fluctuation are then given by
$$\overline{\sigma _x}=\frac{1}{M}\underset{i=1}{\overset{M}{}}\sigma _{x}^{}{}_{i}{}^{}$$
(46)
and
$$\sigma _{\overline{\sigma _x}}=\sqrt{\frac{\overline{\sigma _x^2}\overline{\sigma _x}^2}{M}}.$$
(47)
Eq. (46) represents the standard deviation of the distribution of the binned values $`x`$, and not the standard deviation of an average of these. It does not decrease as the number of measurements $`N`$ is increased, but it is dependent on the number of MC steps in each bin, $`N_{\text{bin}}`$, and will decrease as $`1/\sqrt{N_{\text{bin}}}`$. Hence it is important to state the number of MC steps in the bins for which the deviation is calculated. The statistical error of this standard deviation, $`\sigma _{\overline{\sigma _x}}`$, will, on the other hand, decrease as $`1/\sqrt{N}`$.
In this manner we can calculate the standard deviation for $`f`$. In order to show the standard deviation as a function of the reweighted probability $`W_0^{}`$, Eq. (41), can be inverted to express the necessary weight factor that causes the average to change from $`W_0`$ to $`W_0^{}`$;
$$W=\frac{W_0^{}(1W_0)}{W_0(1W_0^{})}.$$
(48)
Simulation results for the standard deviation of $`f`$ as a function of $`W_0^{}`$ is shown in Fig. 9. We see that the reweighting actually increases the standard deviation. This is due to the rapidly increasing auto correlation times. The autocorrelation function
$$C_\delta (t)=\frac{\delta _{n,0}(i+t)\delta _{n,0}(i)\delta _{n,0}(i)^2}{\delta _{n,0}(i)^2\delta _{n,0}(i)^2}$$
(49)
is shown in Fig. 10, and one can see that the autocorrelation times (inversely proportional to the slopes in Fig. 10) are proportional to $`W_0^{}`$. Notice that the longest autocorrelation times are significantly shorter than the individual bins (consisting of 2000 MC steps) used above, a criteria for the analysis to be valid.
This simple example seems to indicate that reweighting does not decrease the statistical errors. However, in a standard Monte Carlo simulation the measured quantities are not independent even with no reweighting, and formula Eq. (29) contains measured quantities different from Eq. (38) considered above. Therefore the reweighting will affect autocorrelation times differently than in the above example, and reweighting can actually decrease the standard error.
Using the above technique we can study how the relative error in the susceptibility of the semi-frustrated model is affected by reweighting. An initial run without reweighting has to be done first to determine the average sign $`\delta _{n,0}=W_0`$. Thereafter Eq. (48) can be used to determine the desired weight factors. In Fig. 11 the average sign in a simulation of the semifrustrated model with $`J(1)=J(\sqrt{2})=J`$ is shown as a function of lattice volume $`V=L\times L`$ at a temperature $`T/J=1.0`$. For comparison we first performed a standard simulation by sampling all meron sectors, which leads to a severe sign problem with a sign that decreases exponentially in system volume. Next we sampled only the zero- and two-meron sectors without reweighting, which dramatically increases the average sign. The scaling changes from exponential to quadratic in the volume, as can clearly be seen from the graph.
Having determined the average sign without reweighting we now use Eq. (48) to determine the desired weight factors. In Fig. 12 the standard deviation (46) of the susceptibility, calculated using bins containing $`1000`$ MC steps. Results are shown for systems of linear size $`N=8,12,16`$ and 20 at temperature $`T/J=1.0`$. Reweighting clearly helps to reduce the standard deviation, and there is a definite minimum in all these curves indicating an optimal reweighting. The optimally reweighted sign always appears to be less than 0.5, and decreases with decreasing sign (and increasing volume).
Having determined that there is an optimal reweighting we will next consider whether reweighting changes the scaling with system size of the relative statistical error. Let us first consider how the standard deviation scales with no reweighting. Since the sign decreases quadratically in the volume $`V`$ we can derive the scaling of the relative error in the sign, under the ideal (and typically false) assumption that individual measurements are completely independent. Using that $`s=s^2`$ we arrive at
$$\frac{\sigma _s}{\overline{s}}=\frac{\sqrt{\overline{s^2}\overline{s}^2}}{\overline{s}\sqrt{N}}\frac{1}{\sqrt{\overline{s}N}}=\frac{V}{\sqrt{N}}$$
(50)
and this indicates that in this case the statistics needed actually increases linearly in system volume, and not quadratically as stated in Ref. .
In order to study the scaling, the standard deviation for bins containing $`N_{\text{bin}}=10^4`$ MC measurements of the susceptibility is shown in Fig. 13. Four susceptibilites are shown: the $`z`$-component for the semifrustrated model without and with optimal reweighting, the $`x`$-component for the semifrustrated model and the rotationally invariant susceptibility of the isotropic ferromagnetic model. The two latter quantities can be obtained in simulations without sign problems, as discussed above.
Let us first consider the $`z`$-component of the susceptibility for the semifrustrated model without and with optimal reweighting. For both cases the graph suggests a linear increase in relative error. We have to keep in mind that Eq. (50) does not have to be valid, since there are autocorrelations in the simulation, and the results in Fig. 13 do not exclude that the scaling changes when approaching the thermodynamic limit (due to increasing autocorrelation times), but both results do support an approximately linear increase. It appears that the reweighting in this case changes only the prefactor of the volume scaling, but not the exponent. Whether this means that reweighting has completely eliminated the remaining sign problem is not clear (the error remains larger than that for the ferromagnet susceptibility), but it is clear that the reweighting reduces the standard deviations by a significant factor. In any case, an algorithm that changes the functional dependence of the size scaling of the statistics from exponential to polynomial can be considered a solution to the problem.
The $`x`$-susceptibility of the semifrustrated model, which is evaluated with an algorithm without sign problems, shows a constant standard deviation, which may be related to the fact that the susceptibility itself has converged to its thermodynamic limit for these system sizes (see next section). For the isotropic ferromagnetic model, the susceptibility still shows a linearly increasing error, but as already noted the slope is much smaller than for the semifrustrated case.
This concludes our discussion of the reweighting technique. In future work it would be interesting to explore how the optimal reweighting can be determined directly from quantities measured during one single test run, rather than by explicitly measuring the standard deviations as we have done here.
## V RESULTS
In this section we will present results for the semifrustrated and isotropic ferromagnetic models. We will demonstrate that it is feasible to obtain accurate results for large systems in the $`z`$-basis by using the meron-solution. The main motivation for this study is to analyze the meron-solution, and we will only briefly comment on the physics of the semifrustrated model and how it differs from the isotropic ferromagnet. We will primarily consider the semifrustrated model with $`J(1)=J(\sqrt{2})=J`$. To verify the correctness of our codes we have compared simulation results with exact diagonalization data for systems with $`4\times 4`$ spins.
In Fig. 14 the energy is shown for the semifrustrated and ferromagnetic model. The statistical errorbars are smaller than the symbol size and the results have converged to the thermodynamic limit. The largest system size used had $`128\times 128`$ spins. It can be seen that thermal fluctuations and finite-temperature quantum fluctuations more effectively destroy the ferromagnetic correlations for the isotropic ferromagnet than for the semifrustrated model.
The $`z`$-component of susceptibility is shown for both models in Fig. 15. The low-temperature susceptibility for finite-size systems will approach $`\beta N/4`$ for the semifrustrated model, and $`\beta N/12`$ for the ferromagnetic model (due to rotational averaging in the latter case). This can be clearly seen from Fig. 15, where the susceptibility is multiplied by temperature, so that a straight line at low temperatures indicates a Curie divergence. In the thermodynamic limit the uniform susceptibility should diverge exponentially for these models, but this can not be seen in Fig. 15 due to the strong finite-size effects.
In Fig. 16 the $`x`$-component of the susceptibility for the semifrustrated model is shown. This quantity has converged to its thermodynamic limit, and it exhibits a Curie divergence at low temperatures. The ground state value is dependent on the next-nearest neighbor coupling $`J(\sqrt{2})`$, as is clearly demonstrated by plotting two different ratios $`J\sqrt{2}/J(1)=1`$, and $`J\sqrt{2}/J(1)=0`$ (only nearest neighbor interactions) in Fig. 16.
## VI CONCLUSION
We have studied a recently introduced meron-solution to the sign problem within the SSE method. We investigated the sign problem arising in frustrated spin systems and showed that the meron-solution can be a applied to a particular model. The problems arising when applying the meron solution to general models of frustrated spins were discussed. We found that loop algorithm typically are not ergodic and merons do not exist. The sign problem then persists.
For models where the meron solution works we showed that the sign problem can be completely eliminated for certain variables and largely eliminated for other variables. This total and partial elimination comes form a mapping of positive- to negative-weight contributions and involves no approximation. For the variables where the sign problem is almost eliminated the statistical errors can be reduced using a reweighting technique. We study how the relative statistical error behaves as a function of lattice size with and without reweighting, and showed that the reweighting does not change the scaling behavior but typically significantly reduces the over-all magnitude of the fluctuations.
It is evident that the meron solution suffers from the same problem in both frustrated spin and fermionic systems — it is confined to a few special cases. It would be of great importance to be able to extend it to more general cases, e.g., by working in basis where the required loop structure appears. The possibility of finding such bases should be explored.
## VII ACKNOWLEDGMENTS
The research was supported by NSF Grants No. DMR-9629987, and DMR-9712765. P.H. acknowledges support by Finska Vetenskaps-Societeten and Suomalainen Tiedeakatemia. |
warning/0001/hep-ph0001171.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The study of strong interactions at finite baryon density has a long history. Phenomenological knowledge of nuclear forces allows one to obtain a good understanding of equilibrium nuclear matter which is relatively dilute. However, these results cannot be simply extended to higher densities relevant to neutron stars, supernova explosions and relativistic heavy ion collisions, where the microscopic degrees of freedom of QCD, quarks and gluons, become important.
Understanding QCD at finite density has been a tremendous challenge . Unlike finite temperature QCD, where significant progress has been achieved using lattice Monte Carlo simulations, QCD at finite baryon density is not amenable to such a numerical approach. The primary reason is that the determinant of the Euclidean QCD Dirac operator is not real at finite baryon chemical potential, $`\mu `$.
Recent progress has been achieved analytically by studying QCD at infinite density and by studying models of the Nambu-Jona-Lasinio type as well as the instanton liquid model . It was found that the QCD vacuum at sufficiently high baryon density could become a color superconductor . In other words, due to the attraction in the color anti-triplet isosinglet channel the quark Fermi surface becomes unstable towards the formation of a condensate of diquark pairs.
No first principle lattice calculation methods exist at this moment to study the phenomenon of color superconductivity. However, the mechanisms which lead to the formation of diquark condensates (gluon exchange or instanton induced attraction) can be investigated for QCD with two colors. This presents a tremendous advantage. QCD with two colors can be studied numerically on the lattice — the determinant of the Dirac operator is real. Another class of QCD-like theories with diquark condensation that can be studied on the lattice at nonzero chemical potential is QCD with $`any`$ number of colors and of quarks in the adjoint color representation. Clearly, the physics of both types of theories is different from that of the three-color QCD, but the differences are easy to understand and classify. We hope that understanding the behavior of such theories at finite baryon number density will provide us with an additional insight into the phenomenon of diquark condensation in three-color QCD. In addition, numerical simulations of both QCD with two colors and QCD with adjoint quarks are now being pursued by several groups .
The third class of theories which can be studied in the same way is QCD with the phase of the quark determinant quenched. In other words, for each quark such theory contains a conjugate quark, with opposite baryon charge. This leads to the appearance of colorless diquark states — baryonic pions. The zero-flavor limit $`N_f0`$ of such a theory is the quenched approximation of QCD, as has been demonstrated analytically in using a random matrix theory at nonzero $`\mu `$.
The unifying property of all such theories is the pseudo-reality of the quark representations, which manifests itself in the fact that the determinant of the Dirac operator is real. Another unifying property (which is ultimately related to this pseudo-reality) is the fact that diquarks can make color singlets — these are the baryons of such theories. The Bose-Einstein condensation of diquarks, with nonzero baryon charge, may be viewed as baryon charge superconductivity (rather than color superconductivity as in three-color QCD).
In this paper we construct the low-energy chiral Lagrangian describing mesons and baryons (diquarks) at finite baryon density for two-color QCD and QCD with adjoint quarks. As was pointed out in the $`\mu `$-dependence of this Lagrangian describing mesons and baryons can be fixed by global and certain local flavor symmetries. Within the domain of validity of this Lagrangian we study both its vacuum properties and the mass spectrum of the Goldstone modes. Our goal is twofold: (i) to understand and describe quantitatively the physics associated with diquark condensation; (ii) to provide lattice theorists with the qualitative and quantitative predictions aiding their data analysis.
This paper is organized as follows. In section 2 we give an overview of the different patterns of chiral symmetry breaking at $`\mu =0`$ for three-color QCD, two-color QCD and for QCD with adjoint quarks. We relate these patterns to the antiunitary symmetries of the Dirac operator, which provide us with a convenient classification in terms of the Dyson index $`\beta `$. Section 3 reviews the global symmetries of the theories, with the emphasis on the enlarged $`SU(2N_f)`$ symmetry. The spontaneous breaking of this symmetry by the chiral condensate is the subject of section 4, where we identify the Goldstone excitations. The effective Lagrangian at nonzero bare quark mass $`m`$ is introduced in section 5. In section 6 we review the introduction of the chemical potential into the effective Lagrangian by means of a local gauge principle . The vacuum alignment dictated by the effective Lagrangian is analyzed in sections 7 and 8. In section 9 we expand the static part of Lagrangian around the minimum and in section 10 we determine the masses of the excitations as functions of $`\mu `$ and $`m`$ from the pole of their respective propagators. In section 11 we introduce the diquark source $`j`$ and find its effect on the vacuum and the mass spectrum. In section 12 we present the dependence of the vacuum condensates and the baryon number density on $`\mu `$, $`m`$ and $`j`$. Finally, in section 13, we rederive the equation of state of the dilute Bose gas of diquarks with repulsion and show that it exactly matches the equation of state obtained using our mean field analysis of the previous sections. Concluding remarks and discussion are presented in section 14.
## 2 Overview and classification
At zero chemical potential the spontaneous breaking of the chiral symmetry is believed to be an essential low-energy property of three color QCD with fundamental fermions. In the limit of massless quarks, the QCD Lagrangian is invariant under $`U_L(N_f)\times U_R(N_f)`$ transformations, but the ground state is not. The analysis of the hadron spectrum and numerical simulations on the lattice strongly support this assertion . The order parameter of the spontaneous breaking of chiral symmetry is the chiral quark-antiquark condensate. In the case of two-color QCD with fundamental fermions and in the case of any-color QCD with adjoint fermions the symmetry of the Lagrangian is enlarged from $`U_L(N_f)\times U_R(N_f)`$ to $`U(2N_f)`$. In the case of $`N_c=2`$, this symmetry is sometimes referred to as the Pauli-Gürsey symmetry . Also in these cases there is strong evidence both from lattice simulations and from arguments based on supersymmetry that chiral symmetry is broken by a nonzero vacuum expectation value of the chiral condensate.
The pattern of chiral symmetry breaking is determined by two ideas. The Vafa-Witten theorem which tells us that vector symmetries cannot be spontaneously broken, and the idea of maximum breaking of the axial symmetry . For each of the three classes of theories chiral symmetry is broken spontaneously according to different patterns . The axial U(1)<sub>A</sub> subgroup of the global symmetry is broken explicitly by the axial anomaly in all three cases. The vector-like $`U(1)_B`$ symmetry, corresponding to baryon charge conservation, is intact. The remaining symmetries in QCD with three or more colors with fundamental fermions are broken according to $`SU_R(N_f)\times SU_L(N_f)SU_V(N_f)`$. For two-color QCD with fundamental fermions the symmetry is broken according to $`SU(2N_f)Sp(2N_f)`$, whereas for any-color QCD with adjoint fermions the pattern of symmetry breaking is given by $`SU(2N_f)O(2N_f)`$.
We can classify these above three cases by the Dyson index, $`\beta `$, of the Dirac operator with a value of $`\beta =2`$, $`\beta =1`$ and $`\beta =4`$, respectively. The value of $`\beta `$ is given by the number of independent degrees of freedom per matrix element and is determined by the antiunitary symmetries of the Dirac operator. It is a concept that originated in Random Matrix Theory , and is important for the Cartan classification of symmetric spaces .
In the case of two-color QCD the pseudo-real nature of $`SU(2)_{\mathrm{color}}`$ can be expressed as the antiunitary symmetry of the Dirac operator $`𝒟=\gamma _\nu D_\nu +m`$,
$$[𝒟,\tau _2C\gamma _5K]=0\text{ or }𝒟\tau _2C\gamma _5=\tau _2C\gamma _5𝒟^{},(\beta =1)$$
(1)
where $`\tau _2`$ is the color symmetry generator, $`C`$ is the Dirac charge conjugation matrix and $`K`$ is the complex conjugation operator. Since $`(\tau _2CK)^2=1`$ it is always possible to find a basis in which the Dirac operator becomes real which gives $`\beta =1`$.
The symmetry (1) persists even at $`\mu 0`$ . The reality of the Dirac determinant and the feasibility of lattice Monte Carlo simulations is the consequence of (1). This property also allows us to use QCD inequalities at finite $`\mu `$ to show that condensation can only occur in the scalar diquark channel .
In the case of QCD with adjoint quarks the antiunitary symmetry of the Dirac operator is
$$[𝒟,C\gamma _5K]=0\text{ or }𝒟C\gamma _5=C\gamma _5𝒟^{}.(\beta =4)$$
(2)
Since $`(CK)^2=1`$, it is always possible to find a basis in which the Dirac operator can be organized into selfdual (pseudoreal) quaternions. The value of the Dyson index is thus $`\beta =4`$.
In both above cases, $`\beta =1`$ and $`\beta =4`$, the antiunitary symmetry leads to enlargement of the global symmetry to $`U(2N_f)`$ (at $`\mu =0`$).
There is no antiunitary symmetry in the case of QCD with three or more colors with fundamental quarks. The Dirac operator is a complex matrix, thus $`\beta =2`$.
In this paper we shall study two classes of theories: $`\beta =1`$ and $`\beta =4`$, at finite chemical potential $`\mu `$. The main starting point of our analysis is the observation of the fact that in these theories the lowest lying baryons belong to the set of Goldstones of the spontaneously broken extended flavor symmetry $`SU(2N_f)`$. As a result the dependence on $`\mu `$ can be described in the Chiral Perturbation Theory framework .
We shall construct the effective Lagrangian governing the low-momentum modes in the theory. These modes are the Goldstone particles of the spontaneously broken global symmetries. The symmetry of the theory is largest at $`\mu =0`$, $`m=0`$, and so is the number of true Goldstone modes. A nonzero chemical potential $`\mu `$ and/or a bare quark mass $`m`$ removes part of the symmetry and some of the Goldstone modes acquire masses. Our main goal is to find the functional dependence of the masses of such pseudo-Goldstones on $`\mu `$ and $`m`$. As we proceed we learn about many other properties of the theory: condensates and vacuum alignment, phase transitions, etc.
The two cases: two-color fundamental quarks ($`\beta =1`$) and any-color adjoint quarks ($`\beta =4`$) can be analyzed in a similar way. We shall perform this analysis as follows. At each step we shall begin with $`\beta =1`$ case, explaining the concepts and ideas. Then we follow it immediately with the same analysis for the $`\beta =4`$ case with emphasis on the comparison between the two cases, which will help understand both cases better. We shall use similar notations for the objects which are conceptually the same in both cases. In many instances we need not rewrite the formulas, only changing the meaning of the notations suffices. All the formulas which use any properties specific to either $`\beta =1`$ or $`\beta =4`$ are explicitly tagged. The formulas without such tags are general and apply to both cases. As we shall quickly see the two cases naturally complement each other and are, in a certain sense, dual to each other.
## 3 Global symmetries at $`m=\mu =0`$
### 3.1 $`\beta =1`$
The fermionic part of the QCD Lagrangian with 2 fundamental colors is given by
$$=\overline{\psi }\gamma _\nu D_\nu \psi =i\left(\begin{array}{c}\psi _L^{}\\ \psi _R^{}\end{array}\right)^T\left(\begin{array}{cc}\sigma _\nu D_\nu & 0\\ 0& \sigma _\nu ^{}D_\nu \end{array}\right)\left(\begin{array}{c}\psi _L\\ \psi _R\end{array}\right).$$
(3)
We are working in Euclidean space with hermitian $`\gamma `$-matrices and we use spin matrices $`\sigma _\nu =(i,\sigma _k)`$. The quark flavor (as well as color and spin) indices are suppressed and the sum over $`N_f`$ flavors is implied. The symbol $`D_\nu `$ denotes color covariant derivative, $`_\nu +iA_\nu `$, an antihermitian operator, with $`A_\nu `$ being a matrix in color algebra, $`A_\nu =A_\nu ^a\tau _a/2`$. As usual $`\overline{\psi }=\psi ^{}\gamma _0=\psi ^T\gamma _0`$ and in the Euclidean partition function $`\psi `$ and $`\psi ^{}`$ are independent integration variables.
The fact that the Lagrangian (3) has a higher flavor symmetry than the apparent $`U(N_f)\times U(N_f)`$ is related to the pseudoreality of the two-color Dirac operator. In particular, the conjugate field $`\stackrel{~}{\psi }_R=\sigma _2\tau _2\psi _R^{}`$ transforms similarly (in the same color representation) to $`\psi _L`$. Rewriting the Lagrangian (3) we find
$$=i\left(\begin{array}{c}\psi _L^{}\\ \stackrel{~}{\psi }_R^{}\end{array}\right)^T\left(\begin{array}{cc}\sigma _\nu D_\nu & 0\\ 0& \sigma _\nu D_\nu \end{array}\right)\left(\begin{array}{c}\psi _L\\ \stackrel{~}{\psi }_R\end{array}\right)=i\mathrm{\Psi }^{}\sigma _\nu D_\nu \mathrm{\Psi },$$
(4)
where we have used the well-known properties of the Pauli matrices, $`\sigma _2\sigma _\nu ^{}\sigma _2=\sigma _\nu ^T`$ and $`\tau _2\tau _k\tau _2=\tau _k^T`$, taken into account anticommutativity of Grassman variables, dropped total derivatives, and introduced the spinor of dimension $`2N_f`$,
$$\mathrm{\Psi }\left(\begin{array}{c}\psi _L\\ \sigma _2\tau _2\psi _R^{}\end{array}\right)\left(\begin{array}{c}\psi _L\\ \stackrel{~}{\psi }_R\end{array}\right).(\beta =1)$$
(5)
In this form the $`U(2N_f)`$ symmetry becomes manifest. Due to the axial anomaly the symmetry in the corresponding quantum theory is only $`SU(2N_f)`$ (up to discrete symmetries).
### 3.2 $`\beta =4`$
QCD with quarks in the adjoint representation of the color group is described by the Lagrangian as in (3), but with different notations. The fields $`\psi `$ are now transforming according to the adjoint representation of the color group. The color covariant derivative is again given by $`D_\nu =_\nu +iA_\nu `$, but now $`A_\nu `$ is given by the antisymmetric matrix $`(A_\nu )^{bc}=A_\nu ^af_a^{bc}`$, where $`f_a^{bc}=f^{abc}`$ are the generators of the adjoint representation, i.e., the structure constants.
The antisymmetric property of the structure constants now replaces the property of the fundamental generators: $`\tau _2\tau _k\tau _2=\tau _k^Tf_a^{bc}=f_a^{cb}`$. Using this property we can again recast the Lagrangian using spinors of length $`2N_f`$ and obtain (4), but with the spinors $`\mathrm{\Psi }`$ (and $`\stackrel{~}{\psi }_R`$) defined by
$$\mathrm{\Psi }\left(\begin{array}{c}\psi _L\\ \sigma _2\psi _R^{}\end{array}\right)\left(\begin{array}{c}\psi _L\\ \stackrel{~}{\psi }_R\end{array}\right).(\beta =4)$$
(6)
The only difference from (5) is the absence of color $`\tau _2`$ matrix in the definition of $`\stackrel{~}{\psi }_R`$, and, of course, the fact that the spinors $`\psi `$ carry an adjoint, instead of a fundamental, color index. Similarly, in terms of the spinors (6) the $`SU(2N_f)`$ symmetry of the theory becomes manifest.
## 4 Spontaneous symmetry breaking and Goldstones
### 4.1 $`\beta =1`$
Let us now understand transformation properties of $`\overline{\psi }\psi `$, the order parameter of the chiral symmetry breaking, with respect to the $`SU(2N_f)`$ symmetry of the theory. We can rewrite $`\overline{\psi }\psi `$ as
$$\overline{\psi }\psi =\left(\begin{array}{c}\psi _L^{}\\ \psi _R^{}\end{array}\right)^T\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\left(\begin{array}{c}\psi _L\\ \psi _R\end{array}\right)=\frac{1}{2}\mathrm{\Psi }^T\sigma _2\tau _2\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\mathrm{\Psi }+\mathrm{h}.\mathrm{c}..(\beta =1)$$
(7)
The Pauli matrices $`\sigma _2`$ and $`\tau _2`$ ensure antisymmetrization in spin and color indices, to produce a spin and color singlet. We see that the condensate is not invariant under all $`SU(2N_f)`$ rotations. The subgroup which leaves (7) invariant is $`Sp(2N_f)`$. The Goldstone manifold is therefore given by $`SU(2N_f)/Sp(2N_f)`$ with $`N_f(2N_f1)1`$ independent degrees of freedom.
Another way of counting the total number of Goldstone modes is starting from the observation that the condensate (7) is a product of two fundamental $`SU(2N_f)`$ flavor representations, antisymmetric in flavor indices. Therefore, the condensate belongs to an antisymmetric tensor representation the dimension of which is $`N_f(2N_f1)`$. Condensation can occur in any of these $`N_f(2N_f1)`$ directions; the fluctuations along the remaining $`N_f(2N_f1)1`$ directions then become Goldstone modes.
The effective theory for the Goldstone modes can be written in terms of the fluctuations of the orientation of the chiral condensate, $`\mathrm{\Sigma }`$, which, according to the previous paragraph, is an antisymmetric unimodular ($`det\mathrm{\Sigma }=1`$) unitary matrix (exactly $`N_f(2N_f1)1`$ independent components). If we denote the equilibrium value of the orientation of the chiral condensate by $`\mathrm{\Sigma }_c`$, the Goldstone manifold given by $`SU(2N_f)/Sp(2N_f)`$ can be parameterized, according to the transformation of $`\mathrm{\Sigma }`$ under $`SU(2N_f)`$, by
$$\mathrm{\Sigma }=U\mathrm{\Sigma }_cU^T,$$
(8)
where
$`U=\mathrm{exp}\left({\displaystyle \frac{i\mathrm{\Pi }}{2F}}\right)\text{ and }\mathrm{\Pi }=\pi _a{\displaystyle \frac{X_a}{\sqrt{2N_f}}}.`$ (9)
The fields $`\pi _a`$ are the Goldstone modes. The implied sum (over $`a`$) is over the $`N_f(2N_f1)1`$ generators of the coset $`SU(2N_f)/Sp(2N_f)`$ and in order to simplify the algebra in later sections we are using the normalization
$$\mathrm{Tr}X_aX_b=2N_f\delta _{ab}.$$
(10)
The construction of the Goldstone manifold (9) corresponds to the classification of the $`(2N_f)^21`$ generators of the $`SU(2N_f)`$ with respect to a fixed antisymmetric antiunitary matrix, in our case given by $`\mathrm{\Sigma }_c`$, into $`T_k`$ and $`X_a`$ . The $`T_k`$ generators leave $`\mathrm{\Sigma }_c`$ invariant,
$$\mathrm{exp}(i\varphi _kT_k)\mathrm{\Sigma }_c\mathrm{exp}(i\varphi _kT_k)^T=\mathrm{\Sigma }_c,\text{ i.e. }T_k\mathrm{\Sigma }_c=\mathrm{\Sigma }_cT_k^T.$$
(11)
By definition, they are the generators of the symplectic group $`Sp(2N_f)`$. The remaining generators, $`X_a`$, form the coset $`SU(2N_f)/Sp(2N_f)`$. They obey the relation
$$X_a\mathrm{\Sigma }_c=\mathrm{\Sigma }_cX_a^T,\text{ and thus }U\mathrm{\Sigma }_cU^T=U^2\mathrm{\Sigma }_c.$$
(12)
The partition of generators in generators of the coset, $`X_a`$, and generators of the invariant subgroup $`Sp(2N_F)`$, $`T_i`$, depends on the matrix $`\mathrm{\Sigma }_c`$. The defining relations (12) are left unaltered by rotation of $`\mathrm{\Sigma }_c`$ according to
$`\mathrm{\Sigma }_cV\mathrm{\Sigma }_cV^T,`$ (13)
where $`V`$ is an $`SU(2N_f)`$ matrix, and a simultaneous rotation of the generators by
$`X_aVX_aV^{}.`$ (14)
This means that the set of broken generators, $`X_a`$, changes, if the matrix $`\mathrm{\Sigma }_c`$ is changed.
If we use, as we do below, the following choice for $`\mathrm{\Sigma }_c`$:
$$\mathrm{\Sigma }_c=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),(\beta =1),$$
(15)
the generators $`X_a`$ can be written in the block representation as
$$\mathrm{\Pi }=\left(\begin{array}{cc}P^T& Q\\ Q^{}& P\end{array}\right),$$
(16)
with $`N_f\times N_f`$ matrices $`P`$ and $`Q`$ such that $`\mathrm{Tr}P=0`$, $`P^{}=P`$, and $`Q^T=Q`$. It is then easy to see that the number of independent components in $`P`$ and $`Q`$ are,
$$N_P=N_f^21\text{ and }N_Q=N_f(N_f1).(\beta =1)$$
(17)
### 4.2 $`\beta =4`$
What are the transformation properties of the chiral condensate with respect to the $`SU(2N_f)`$ symmetry in the $`\beta =4`$ case? We can rewrite $`\overline{\psi }\psi `$ as
$$\overline{\psi }\psi =\left(\begin{array}{c}\psi _L^{}\\ \psi _R^{}\end{array}\right)^T\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\left(\begin{array}{c}\psi _L\\ \psi _R\end{array}\right)=\frac{1}{2}\mathrm{\Psi }^T\sigma _2\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\mathrm{\Psi }+\mathrm{h}.\mathrm{c}..(\beta =4)$$
(18)
The Pauli matrix $`\sigma _2`$ ensures the antisymmetrization with respect to spin indices, to produce a spin singlet. There is no antisymmetrization in color (unlike in the case of fundamental colors $`\beta =1`$) — the color singlet is a symmetric product of two adjoint representations. The Pauli principle now demands that the $`SU(2N_f)`$ flavor indices must be symmetric, and they are, as is evident in (18). Therefore the condensate belongs to the symmetric (as opposed to antisymmetric for $`\beta =1`$) second rank tensor representation of $`SU(2N_f)`$, which has dimension $`N_f(2N_f+1)`$. The condensation can occur in any of the $`N_f(2N_f+1)`$ directions. The fluctuations in the remaining $`N_f(2N_f+1)1`$ directions become Goldstone bosons. Alternatively, since the chiral condensate (18) is invariant under $`SO(2N_f)`$, the Goldstone manifold is $`SU(2N_f)/SO(2N_f)`$, which gives us the same number of Goldstone modes.
The Goldstone manifold $`SU(2N_f)/SO(2N_f)`$ should now be parameterized by symmetric unimodular unitary matrices $`\mathrm{\Sigma }`$. The Goldstone fields are introduced in the same way as in (8) and (9), i.e., $`\mathrm{\Sigma }=U\mathrm{\Sigma }_cU^T`$, where $`\mathrm{\Sigma }_c`$ is now also a symmetric unimodular unitary matrix and the implied sum in (9) is over the generators $`X_a`$ of the coset $`SU(2N_f)/SO(2N_f)`$.
The classification of the generators is also similar to $`\beta =1`$ case. The generators of the coset, $`X_a`$ obey the commutation relation (12) with a given symmetric unitary matrix, also denoted by $`\mathrm{\Sigma }_c`$, whereas the $`T_i`$ generators leave $`\mathrm{\Sigma }_c`$ invariant as in (11). The $`T_i`$ are now the generators of an $`SO(2N_f)`$ subgroup (as opposed to $`Sp(2N_f)`$ in the $`\beta =1`$ case). The remaining generators $`X_a`$ form a coset $`SU(2N_f)/SO(2N_f)`$, which is the Goldstone manifold in this case.
In this case the standard choice for the matrix $`\mathrm{\Sigma }_c`$ is
$$\mathrm{\Sigma }_c=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).(\beta =4)$$
(19)
With this choice, the matrix $`\mathrm{\Pi }=\pi _aX_a/\sqrt{2N_f}`$ in the coset can be split into 4 blocks of size $`N_f\times N_f`$ as in (16). The matrix $`P`$ is again hermitian and traceless, but the matrix $`Q`$ is now symmetric. Therefore, the counting of the independent degrees of freedom in the $`\beta =4`$ case is
$$N_P=N_f^21\text{ and }N_Q=N_f(N_f+1).(\beta =4)$$
(20)
In both cases, $`\beta =1`$ and $`\beta =4`$, the kinetic term of the effective Lagrangian describing the Goldstone modes should be invariant under the global $`SU(2N_f)`$ group and under Lorentz transformation. The corresponding $`SU(2N_f)`$ nonlinear sigma-model is given by
$$_{\mathrm{eff}}=\frac{F^2}{2}\mathrm{Tr}_\nu \mathrm{\Sigma }_\nu \mathrm{\Sigma }^{},$$
(21)
where $`F`$ is the pion decay constant.
## 5 Bare quark mass $`m`$
### 5.1 $`\beta =1`$
If the bare quark mass $`m`$ (in this paper, the same for all quarks) is not zero, an explicit $`SU(2N_f)`$ breaking term in the effective Lagrangian (21) appears. To determine its form we first rewrite the bare mass term using $`SU(2N_f)`$ notations,
$$m\overline{\psi }\psi =\frac{1}{2}m\mathrm{\Psi }^T\sigma _2\tau _2\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\mathrm{\Psi }+\mathrm{h}.\mathrm{c}.=\frac{1}{2}\mathrm{\Psi }^T\sigma _2\tau _2M\mathrm{\Psi }+\mathrm{h}.\mathrm{c}.,(\beta =1)$$
(22)
where the mass matrix $`M`$ is given by
$$M=m\widehat{M}\mathrm{and}\widehat{M}=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).(\beta =1)$$
(23)
We see explicitly that the bare mass term is only invariant under an $`Sp(2N_f)`$ subgroup of $`SU(2N_f)`$. The full $`SU(2N_f)`$ invariance can be restored if $`M`$ is also transformed together with $`\mathrm{\Psi }`$ according to
$$\mathrm{\Psi }V\mathrm{\Psi }\text{ and }MV^{}MV^{}.$$
(24)
This extended symmetry must be also manifest in the effective theory,
$$\mathrm{\Sigma }V\mathrm{\Sigma }V^T\text{ and }MV^{}MV^{}.$$
(25)
The lowest order term induced by the quark mass must therefore have the form
$$_{\mathrm{q}.\mathrm{mass}}=G\mathrm{ReTr}(M\mathrm{\Sigma })=mG\mathrm{ReTr}(\widehat{M}\mathrm{\Sigma }).$$
(26)
One can view $`\mathrm{ReTr}(\widehat{M}\mathrm{\Sigma })`$ as a generalized cosine of the angle between unitary matrices $`\mathrm{\Sigma }`$ and $`\widehat{M}^{}`$. It is maximal when $`\mathrm{\Sigma }`$ is aligned with $`\widehat{M}^{}`$. Therefore the direction of $`\mathrm{\Sigma }`$ minimizing (26) is given by
$$\mathrm{\Sigma }_c=\widehat{M}^{},$$
(27)
which leads us to our choice of $`\mathrm{\Sigma }_c`$ (15). The mass term comes with a phenomenological coefficient, which we denote by $`G`$. It is given by the derivative of the vacuum energy with respect to $`m`$ and is, therefore, proportional to the chiral condensate in the chiral limit $`m0`$ at $`\mu =0`$ (see Section 12),
$$G=\frac{\overline{\psi }\psi _0}{2N_f}.$$
(28)
The resulting Lagrangian with the mass term,
$$_{\mathrm{eff}}=\frac{F^2}{2}\mathrm{Tr}_\nu \mathrm{\Sigma }_\nu \mathrm{\Sigma }^{}mG\mathrm{ReTr}(\widehat{M}\mathrm{\Sigma }),$$
(29)
is the familiar Chiral Perturbation Theory Lagrangian at lowest order in the momentum expansion . Expanded to second order in the pion fields according to (8), (9), it yields a spectrum with $`N_f(2N_f1)1`$ degenerate (pseudo-)Goldstones with masses given by the usual Gell-Mann$``$Oakes$``$Renner relation
$$m_\pi ^2=\frac{mG}{F^2}.$$
(30)
We can use this relation to trade $`G`$ for another parameter, $`m_\pi `$, and write
$$_{\mathrm{eff}}=\frac{F^2}{2}[\mathrm{Tr}_\nu \mathrm{\Sigma }_\nu \mathrm{\Sigma }^{}2m_\pi ^2\mathrm{Re}\mathrm{Tr}(\widehat{M}\mathrm{\Sigma }).].$$
(31)
The symmetry of the theory is reduced from $`SU(2N_f)`$ to $`Sp(2N_f)`$ by the mass term. Since the chiral condensate does not break any more symmetries in this case we do not have any true Goldstones when $`m0`$.
### 5.2 $`\beta =4`$
To determine the form of the term in the effective Lagrangian induced by a small bare quark mass $`m`$ we rewrite the quark mass term as
$$m\overline{\psi }\psi =\frac{1}{2}m\mathrm{\Psi }^T\sigma _2\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\mathrm{\Psi }+\mathrm{h}.\mathrm{c}.=\frac{1}{2}\mathrm{\Psi }^T\sigma _2M\mathrm{\Psi }+\mathrm{h}.\mathrm{c}.,(\beta =4)$$
(32)
where we have used the spinors of length $`N_f`$ introduced in section 3. In this case the mass matrix is given by
$$M=m\widehat{M}\mathrm{with}\widehat{M}=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).(\beta =4)$$
(33)
Using the same arguments as in section 5.1 we find that the mass term in the effective Lagrangian as dictated by the extended flavor symmetry is given again by (26), with $`\mathrm{\Sigma }`$ now being a symmetric matrix and $`M`$ given by (33). Similarly, this term is minimized when $`\mathrm{\Sigma }=\widehat{M}^{}`$. With $`\widehat{M}`$ now taken from (33) we arrive at our choice of $`\mathrm{\Sigma }_c`$ (19).
The symmetry of the theory is reduced from $`SU(2N_f)`$ to $`SO(2N_f)`$ by the mass term. Since the chiral condensate does not break any more symmetries in this case we do not have any truly massless Goldstones when $`m0`$. The masses of all $`N_f(2N_f+1)`$ (pseudo-)Goldstones are equal and are given by the Gell-Mann$``$Oakes$``$Renner relation (30).
## 6 Chemical potential $`\mu `$
In this section we review the introduction of the chemical potential in the effective partition function following the approach of . This approach relies only on the $`SU(2N_f)`$ symmetry of the theory and the resulting $`\mu `$-dependent terms are the same for both cases, $`\beta =1`$ and $`\beta =4`$.
### 6.1 Global symmetries and $`\mu `$
At nonzero chemical potential the microscopic Lagrangian is given by
$$=\overline{\psi }\gamma _\nu D_\nu \psi \mu \overline{\psi }\gamma _0\psi +m\overline{\psi }\psi .$$
(34)
As was the case for the mass term, we can also rewrite the baryon charge density in terms of the $`SU(2N_f)`$ spinors.
$`\overline{\psi }\gamma _0\psi =\left(\begin{array}{c}\psi _L^{}\\ \psi _R^{}\end{array}\right)^T\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\left(\begin{array}{c}\psi _L\\ \psi _R\end{array}\right)=\mathrm{\Psi }^{}\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\mathrm{\Psi }\mathrm{\Psi }^{}B\mathrm{\Psi };`$ (43)
$`B\left(\begin{array}{cc}+1& 0\\ 0& 1\end{array}\right).`$ (46)
The physical meaning of $`+1`$ and $`1`$ in the baryon charge matrix $`B`$ is simple: they are the baryon charges of the quarks $`\psi _L`$ and conjugate quarks $`\stackrel{~}{\psi }_R`$. We see that this term is not a singlet under $`SU(2N_f)`$. It transforms in the adjoint representation of this group, in other words, the baryon charge is one of the $`(2N_f)^21`$ generators of this group. In terms of the spinors of length $`2N_f`$, the microscopic Lagrangian is thus given by
$`=i\mathrm{\Psi }^{}\sigma _\nu (D_\nu \mu B_\nu )\mathrm{\Psi }[{\displaystyle \frac{1}{2}}\mathrm{\Psi }^T\left\{\begin{array}{c}\sigma _2\tau _2\\ \sigma _2\end{array}\right\}M\mathrm{\Psi }+\mathrm{h}.\mathrm{c}.],`$ (49)
where the upper branch corresponds to $`\beta =1`$ and the lower branch to $`\beta =4`$. For reasons that will become clear in the next subsection, we have introduced the four-vector $`B_\nu =(B,\mathrm{𝟎})`$.
As in the case of the quark mass term, the chemical potential term, $`\mu \mathrm{\Psi }^{}B\mathrm{\Psi }`$ violates the $`SU(2N_f)`$ symmetry. Similarly, we can maintain this symmetry by accompanying the rotation of $`\mathrm{\Psi }`$ by a corresponding rotation of $`B`$,
$$\mathrm{\Psi }V\mathrm{\Psi }\text{ and }BVBV^{}.$$
(50)
Such an extended symmetry must be manifest in the effective Lagrangian,
$$\mathrm{\Sigma }V\mathrm{\Sigma }V^T\text{ and }BVBV^{}.$$
(51)
This restricts the lowest order nonderivative term in $`\mu `$ to a linear combination of
$$\mu ^2\mathrm{Tr}(\mathrm{\Sigma }B^T\mathrm{\Sigma }^{}B)\text{ and }\mu ^2\mathrm{Tr}(BB)$$
(52)
with arbitrary coefficients. Only the first of these terms contains a dependence on the Goldstone fields.<sup>1</sup><sup>1</sup>1Such type of symmetry breaking terms also occur in the context of the non-hermitian Random Matrix Theory .
At $`m=0`$ the chemical potential breaks the global symmetry of the theory from $`SU(2N_f)`$ down to the usual (as in three-color QCD) $`SU(N_f)_L\times SU(N_f)_R\times U(1)_B`$. If $`m0`$ also, the symmetry of the theory is only $`SU(N_f)_V\times U(1)_B`$.
### 6.2 Local symmetry and the coefficient of the $`\mu ^2`$ term
The coefficients of the terms (52) in the effective Lagrangian can be related to baryon number susceptibility, i.e., the second derivative of the vacuum energy with respect to $`\mu `$. However, unlike $`G`$, i.e. the chiral condensate $`\overline{\psi }\psi _0`$, these parameters are not independent. They are related to the pion decay constant $`F`$ by virtue of a local symmetry .
As was observed in , one can further extend the symmetry (50) of the microscopic Lagrangian (49) to include local $`SU(2N_f)`$ flavor transformations
$$\mathrm{\Psi }V\mathrm{\Psi }\text{ and }B_\nu VB_\nu V^{}\frac{1}{\mu }V_\nu V^{}.$$
(53)
We can also recover the Lorentz (Euclidean) symmetry by transforming $`B_\nu `$ as a four-vector.
To make such an extended local symmetry (and also Lorentz symmetry) manifest in the effective Lagrangian (31) we must replace the normal derivative by a flavor covariant derivative,
$`_\nu \mathrm{\Sigma }=_\nu \mathrm{\Sigma }\mu (B_\nu \mathrm{\Sigma }+\mathrm{\Sigma }B_\nu ^T);`$
$`_\nu \mathrm{\Sigma }^{}=_\nu \mathrm{\Sigma }^{}+\mu (\mathrm{\Sigma }^{}B_\nu +B_\nu ^T\mathrm{\Sigma }^{}).`$ (54)
Thus we arrive at the effective Lagrangian of lowest order in the momentum expansion
$`_{\mathrm{eff}}`$ $`=`$ $`{\displaystyle \frac{F^2}{2}}\left[\mathrm{Tr}_\nu \mathrm{\Sigma }_\nu \mathrm{\Sigma }^{}2m_\pi ^2\mathrm{ReTr}(\widehat{M}\mathrm{\Sigma })\right]`$ (55)
$`={\displaystyle \frac{F^2}{2}}\mathrm{Tr}_\nu \mathrm{\Sigma }_\nu \mathrm{\Sigma }^{}+2\mu F^2\mathrm{Tr}B\mathrm{\Sigma }^{}_0\mathrm{\Sigma }`$
$`F^2\mu ^2\mathrm{Tr}\left(\mathrm{\Sigma }B^T\mathrm{\Sigma }^{}B+BB\right)F^2m_\pi ^2\mathrm{ReTr}\left(\widehat{M}\mathrm{\Sigma }\right).`$
It is important to note that the dependence on $`\mu `$ comes with no additional parameters. It is completely fixed, by the local symmetry, in terms of an already existing parameter $`F`$.
## 7 Vacuum alignment
The static part of the Lagrangian (55) determines the vacuum alignment of the field $`\mathrm{\Sigma }`$ as well as the masses of the excitations. This part of the Lagrangian has the form
$`_{\mathrm{st}}(\mathrm{\Sigma })`$ $`=`$ $`F^2\mu ^2\mathrm{Tr}\left(\mathrm{\Sigma }B^T\mathrm{\Sigma }^{}B+BB\right)F^2m_\pi ^2\mathrm{ReTr}(\widehat{M}\mathrm{\Sigma })`$ (56)
$`=`$ $`{\displaystyle \frac{F^2m_\pi ^2}{2}}\left[{\displaystyle \frac{x^2}{2}}\mathrm{Tr}\left(\mathrm{\Sigma }B^T\mathrm{\Sigma }^{}B+BB\right)2\mathrm{R}\mathrm{e}\mathrm{T}\mathrm{r}(\widehat{M}\mathrm{\Sigma })\right],`$
where we introduced $`x=2\mu /m_\pi `$. The $`\mu ^2`$ and $`m_\pi ^2`$ terms in (56) compete for the direction of the condensation which we denote by $`\overline{\mathrm{\Sigma }}`$. For $`x=0`$ the orientation of $`\overline{\mathrm{\Sigma }}`$ is determined by the mass matrix $`\overline{\mathrm{\Sigma }}=\widehat{M}^{}`$. This value we denoted by $`\mathrm{\Sigma }_c`$. This is the orientation of the usual chiral condensate which carries no baryon charge. When $`x=\mathrm{}`$ the static Lagrangian is minimized on a manifold in the space of $`\mathrm{\Sigma }`$, from which we choose the following value and denote it by $`\mathrm{\Sigma }_d`$,
$$\mathrm{\Sigma }_d=\left(\begin{array}{cc}iI& 0\\ 0& iI\end{array}\right),(\beta =1)$$
(57)
where $`I`$ is an antisymmetric $`N_f\times N_f`$ matrix, which, written in $`(N_f/2)\times (N_f/2)`$ blocks, has the form
$$I=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).$$
(58)
This minimum (even for $`m0`$, unlike $`\mathrm{\Sigma }_c`$) has a degeneracy, which for $`N_f=2`$, is simply a rotation by the generator $`B`$. The condensate $`\mathrm{\Sigma }_d`$ breaks spontaneously the baryon number symmetry. The degeneracy leads to the appearance of a massless Goldstone boson. For $`N_f>2`$ this condensate breaks more than just the $`U(1)_B`$ symmetry. For $`m=0`$ it also breaks $`SU(N_f)_L\times SU(N_f)_R`$ down to $`Sp(N_f)_L\times Sp(N_f)_R`$ given by the rotations leaving the block matrix $`I`$ in (57) invariant . When $`m0`$ it breaks the $`SU(N_f)_V`$ symmetry down to $`Sp(N_f)_V`$.<sup>2</sup><sup>2</sup>2Similarly to the alignment of $`\mathrm{\Sigma }_c`$ to the bare quark mass matrix $`M^{}`$, the direction of $`\mathrm{\Sigma }_d`$ is determined by an external diquark source, $`J`$, which, as we shall see in section 11, breaks the degeneracy. When $`J`$ is zero, we can choose any orientation of $`\mathrm{\Sigma }_d`$ within the manifold of minima, and the results will not change due to the symmetry relating all such minima.
The situation is similar in the $`\beta =4`$ case. The only difference is that now the matrix $`\mathrm{\Sigma }`$ is a symmetric $`SU(2N_f)`$ matrix. It rotates between the value of $`\mathrm{\Sigma }_c`$ given by (19) and $`\mathrm{\Sigma }_d`$, which we chose as
$$\mathrm{\Sigma }_d=\left(\begin{array}{cc}i& 0\\ 0& i\end{array}\right).(\beta =4)$$
(59)
The discussion of the preceding paragraph carries over to the $`\beta =4`$ case with the already familiar substitution of $`Sp(N_f)`$ by $`SO(N_f)`$. The condensate $`\mathrm{\Sigma }_d`$ breaks $`SU(N_f)_{L,R}`$ flavor symmetries down to $`SO(N_f)_{L,R}`$, in addition to breaking $`U(1)_B`$.
At intermediate values of $`x`$ the orientation of the condensate $`\overline{\mathrm{\Sigma }}`$ rotates, as a function of $`x`$, from $`\mathrm{\Sigma }_c`$ to $`\mathrm{\Sigma }_d`$. We shall prove that it can always be written as the linear combination
$$\overline{\mathrm{\Sigma }}=\mathrm{\Sigma }_\alpha \mathrm{\Sigma }_c\mathrm{cos}\alpha +\mathrm{\Sigma }_d\mathrm{sin}\alpha ,$$
(60)
where $`\alpha =0`$ for $`x=0`$ and $`\pi /2`$ for $`x=\mathrm{}`$.
The angle $`\alpha `$ is a function of $`x`$. It is determined by substituting the value of $`\mathrm{\Sigma }`$ given by (60) into (56). This results in the static Lagrangian
$$_{\mathrm{st}}(\mathrm{\Sigma }_\alpha )=F^2m_\pi ^2N_f\left[\frac{x^2}{2}(\mathrm{cos}2\alpha 1)2\mathrm{cos}\alpha \right].$$
(61)
Minimizing with respect to $`\alpha `$ we find
$`\alpha =0,`$ $`\text{when }x<1;`$
$`\mathrm{cos}\alpha ={\displaystyle \frac{1}{x^2}},`$ $`\text{when }x>1.`$ (62)
Note that $`\alpha =0`$ is always an extremum of (61) but it becomes a maximum for $`x>1`$.
We find that the condensate is a non-analytic function of $`x=2\mu /m_\pi `$. There is a second order phase transition as a function of $`\mu `$ at $`x=1`$. For small $`\mu <m_\pi /2`$ the vacuum does not change. At $`\mu =m_\pi /2`$ a transition occurs and the direction of the condensate starts rotating from that of $`\mathrm{\Sigma }_c`$ to that of $`\mathrm{\Sigma }_d`$. A nonzero value of the projection on $`\mathrm{\Sigma }_d`$, proportional to $`\mathrm{sin}\alpha `$, means that the vacuum breaks the baryon number symmetry spontaneously. This happens at the value of $`\mu `$ equal to 1/2 of the mass of the lightest baryon in the system – the diquark, or, the baryonic pion.
## 8 Global minimum
In this section we show that the minimum of the static potential is given by $`\mathrm{\Sigma }_\alpha `$ defined in (60). The argument in this section shows that it is an absolute, or global, minimum of $`_{\mathrm{st}}(\mathrm{\Sigma })`$. In the next section we expand the Goldstone fields to second order about this minimum.
### 8.1 $`\beta =1`$
We decompose the antisymmetric unitary matrix $`\mathrm{\Sigma }`$ into 4 blocks of size $`N_f\times N_f`$,
$$\mathrm{\Sigma }=\left(\begin{array}{cc}A& C\\ C^T& B\end{array}\right).(\beta =1)$$
(63)
The antisymmetric matrices $`A`$ and $`B`$ satisfy the following unitarity constraints:
$$AA^{}+CC^{}=1,BB^{}+(C^{}C)^T=1,AC^{}=CB^{}.(\beta =1)$$
(64)
Taking this into account we can express $`_{\mathrm{st}}`$ entirely in terms of the matrix $`C`$,
$`_{\mathrm{st}}(\mathrm{\Sigma })=F^2m_\pi ^2\left[x^2\mathrm{Tr}\left(C{\displaystyle \frac{1}{x^2}}\right)\left(C^{}{\displaystyle \frac{1}{x^2}}\right)N_f\left(x^2+{\displaystyle \frac{1}{x^2}}\right)\right].`$ (65)
Ignore, for the moment, the constraints on the matrix elements of $`C`$. The trace in (65) can be viewed as the distance in the $`2N_f^2`$ dimensional space of real and imaginary parts of the matrix elements of $`C`$ from the point given by diagonal matrix $`1/x^2`$. For $`x>1`$ the absolute minimum is achieved when $`C=1/x^2`$. The matrices $`A`$ and $`B`$ can be chosen, for example, as $`A=B=iI\sqrt{11/x^4}`$, to satisfy all the constraints (64). If we define $`\mathrm{cos}\alpha =1/x^2`$, we observe that the resulting matrix in (63) is given by $`\mathrm{\Sigma }_\alpha `$ (60).
When $`x<1`$, we notice that unitarity constraints (64) demand that $`\mathrm{Tr}CC^{}1`$. This means that we have to consider the points in the space of $`C`$ only within the unit hypersphere around $`C=0`$. It is easy to see that the minimum distance within this sphere is at the surface point closest to $`1/x^2`$, i.e. at $`C=1`$. The constraints (64) can be satisfied only by $`A=B=0`$ and the resulting matrix in (63) is $`\mathrm{\Sigma }_c`$, the minimum of the static potential at $`\mu =0`$.
### 8.2 $`\beta =4`$
The same analysis applies in the $`\beta =4`$ case. Since $`\mathrm{\Sigma }`$ is now a symmetric matrix its block decomposition is given by
$$\mathrm{\Sigma }=\left(\begin{array}{cc}A& +C\\ C^T& B\end{array}\right),(\beta =4)$$
(66)
with symmetric matrices $`A`$ and $`B`$ and with unitarity constraints
$$AA^{}+CC^{}=1,BB^{}+(C^{}C)^T=1,AC^{}=CB^{}.(\beta =4)$$
(67)
The static part of the Lagrangian is the same as in (65). The minimum is given by $`C=1`$ for $`x<1`$ and $`C=1/x^2`$ for $`x>1`$. The matrices $`A`$ and $`B`$ are not determined uniquely by the unitarity constraints. One can check that $`A=B=i\sqrt{11/x^4}`$ (as in (60)) satisfy these constraints. <sup>3</sup><sup>3</sup>3 As in the case of $`\beta =1`$, the case of odd number of flavors has to be investigated separately. Let us consider $`N_f=1`$. Then a symmetric unitary matrix has two degrees of freedom and can be parameterized as $`\mathrm{\Sigma }=\left(\begin{array}{cc}i\mathrm{sin}\theta e^{i\varphi }& \mathrm{cos}\theta \\ \mathrm{cos}\theta & i\mathrm{sin}\theta e^{i\varphi }\end{array}\right).`$ (70) Notice that the determinant of $`\mathrm{\Sigma }`$ has to be equal to the determinant of the mass matrix $`\widehat{M}`$ (33). We thus have that the matrix $`C`$ in (66) is just a number parameterized according to $`C=\mathrm{cos}\theta `$, leading to the effective potential (up to constants) $`x^2\mathrm{Tr}\left(C^{}{\displaystyle \frac{1}{x^2}}\right)\left(C{\displaystyle \frac{1}{x^2}}\right)x^2{\displaystyle \frac{1}{x^2}}=x^2\mathrm{cos}^2\theta 2\mathrm{cos}\theta x^2.`$ The minimum is at $`\theta =0`$ for $`x<1`$ and at $`\mathrm{cos}\theta =1/x^2`$ for $`x>1`$. In the diquark condensation phase ($`x>1`$), we find one massless excitation (the $`\varphi `$-mode), a nonzero baryon density and nonzero chiral and diquark condensates.
## 9 Curvatures at the minimum
We have established that the global minimum of the effective Lagrangian is achieved at the value of the field $`\mathrm{\Sigma }`$ given by (60) and (7). The orientation of the condensate rotates as a function of $`x=2\mu /m_\pi `$ (nonanalytic at $`x=1`$). In this section we shall expand the effective Lagrangian up to the second order in the fluctuations of $`\mathrm{\Sigma }`$ which will help us to determine the (pseudo-)Goldstone masses and their dependence on $`x`$.
### 9.1 Normal phase: $`\mu <m_\pi /2`$
When $`x<1`$ the vacuum orientation of $`\mathrm{\Sigma }`$ does not depend on $`x`$ and is given by $`\overline{\mathrm{\Sigma }}=\mathrm{\Sigma }_c`$. Expanding $`\mathrm{\Sigma }`$ around $`\mathrm{\Sigma }_c`$ using the Goldstone fields defined in (9) according to
$`\mathrm{\Sigma }=U\mathrm{\Sigma }_cU^T=U^2\mathrm{\Sigma }_c=\left(1+{\displaystyle \frac{i\mathrm{\Pi }}{F}}{\displaystyle \frac{\mathrm{\Pi }^2}{2F^2}}+\mathrm{}\right)\mathrm{\Sigma }_c,`$ (71)
we find
$$_{\mathrm{st}}(\mathrm{\Sigma })=_{\mathrm{st}}(\mathrm{\Sigma }_c)+\frac{m_\pi ^2}{2}\left[\frac{x^2}{4}\mathrm{Tr}[B,\mathrm{\Pi }]^2+\mathrm{Tr}\mathrm{\Pi }^2\right]+\mathrm{},$$
(72)
where ellipsis denotes terms of higher powers of $`\mathrm{\Pi }`$. The commutator in this Lagrangian can be written as
$`[B,\mathrm{\Pi }]=b\mathrm{\Pi },`$ (73)
with $`b`$ the baryon charge of the pseudo-Goldstone modes. Since all our pseudo-Goldstone modes are quark-antiquark or (anti-)diquark states the values of $`b`$ are $`b=0,\pm 2`$. The curvature of the $`P`$, $`Q`$ and $`Q^{}`$ (16) modes, thus depends on $`\mu `$ through the baryon charge and is given by $`m_\pi ^2(b\mu )^2`$. For example, in the $`\beta =1`$ case of $`N_f=2`$, there are 3 pseudoscalar mesons, 1 diquark and 1 antidiquark. Using the block decomposition of the generators (16) we find
$$_{\mathrm{st}}(\mathrm{\Sigma })=_{\mathrm{st}}(\mathrm{\Sigma }_c)+m_\pi ^2\mathrm{Tr}\left[P^2+(1x^2)QQ^{}\right]+\mathrm{}.$$
(74)
Two comments are in order here. First, note that, the curvature (and the mass) of the meson modes $`P`$ does not depend on $`\mu `$ in the normal phase. Second, at $`x=1`$ the curvature of the diquark modes $`Q`$ vanish, which signals a phase transition and the onset of diquark condensation.
### 9.2 Diquark condensation phase: $`\mu >m_\pi /2`$.
In this phase the condensate begins to rotate according to (60), (7). This rotation can be also written as
$$\mathrm{\Sigma }_\alpha =V_\alpha \mathrm{\Sigma }_cV_\alpha ^T=V_\alpha ^2\mathrm{\Sigma }_c,\text{ where }V_\alpha ^2=e^{i\alpha X_2},$$
(75)
and $`X_2`$ is the generator that rotates $`\mathrm{\Sigma }_c`$ into $`\mathrm{\Sigma }_d`$. Comparing (75) (60) we find
$$\mathrm{\Sigma }_d=iX_2\mathrm{\Sigma }_c.$$
(76)
This generator belongs to the set of broken generators with respect to $`\mathrm{\Sigma }_c`$ (as well as with respect to $`\mathrm{\Sigma }_d`$) since it satisfies (12).
We could parameterize fluctuations of $`\mathrm{\Sigma }`$ around the vacuum value given by $`\mathrm{\Sigma }_\alpha `$ as
$$\mathrm{\Sigma }=U_\alpha \mathrm{\Sigma }_\alpha U_\alpha ^T,$$
(77)
where $`U_\alpha `$ are unitary matrices generated by rotated (in the sense of (14)) generators, $`V_\alpha X_aV_\alpha ^{}`$, instead of $`X_a`$ (8), (9). However, as we shall find, the meson mass matrix is diagonal in the basis given by the parametrization
$$\mathrm{\Sigma }=V_\alpha U\mathrm{\Sigma }_cU^TV_\alpha ^T,$$
(78)
where we have rotated the coset generators $`X_a`$ back to their $`\alpha =0`$ values using (13), (14) and (75).
Before expanding the static Lagrangian, we substitute (78) into (56) to obtain
$`_{\mathrm{st}}(\mathrm{\Sigma })={\displaystyle \frac{F^2m_\pi ^2}{2}}`$ $`[{\displaystyle \frac{x^2}{2}}\mathrm{Tr}(U^2\mathrm{\Sigma }_c\left(V_{\alpha }^{}{}_{}{}^{}BV_\alpha \right)^T\mathrm{\Sigma }_c^{}\left(U^{}\right)^2\left(V_{\alpha }^{}{}_{}{}^{}BV_\alpha \right)+BB)`$ (79)
$`2\mathrm{R}\mathrm{e}\mathrm{T}\mathrm{r}\left(\left(V_\alpha ^T\widehat{M}V_\alpha \right)U^2\mathrm{\Sigma }_c\right)].`$
The rotated values of $`B`$ and $`\mathrm{\Sigma }_c`$ (by angle $`\alpha `$ in the sense of (75) ) can be expressed as
$`V_\alpha ^{}BV_\alpha =V_\alpha BV_\alpha ^{}=B\mathrm{cos}\alpha +X_1\mathrm{sin}\alpha ;`$
$`V_\alpha ^T\widehat{M}V_\alpha =(V_\alpha \widehat{M}^{}V_\alpha ^T)^{}=\mathrm{\Sigma }_\alpha ^{}=\mathrm{\Sigma }_c^{}\mathrm{cos}\alpha \mathrm{\Sigma }_d^{}\mathrm{sin}\alpha ,`$ (80)
where we have introduced $`X_1`$ by
$$X_1=iBX_2.$$
(81)
This generator, similarly to $`X_2`$, belongs to the coset of broken generators with respect to $`\mathrm{\Sigma }_c`$ (12). It is the generator into which $`B`$ is rotated while $`\mathrm{\Sigma }_c`$ is rotated into $`\mathrm{\Sigma }_d`$.
Substituting (9.2) into (79) and expanding to second order in $`\mathrm{\Pi }`$ (71) we find
$`_{\mathrm{st}}(\mathrm{\Sigma })`$ $`=`$ $`_{\mathrm{st}}(\mathrm{\Sigma }_\alpha )+Fm_\pi ^2\left[{\displaystyle \frac{x^2}{2}}\mathrm{sin}2\alpha +\mathrm{sin}\alpha \right]\mathrm{Tr}\left(X_2\mathrm{\Pi }\right)`$ (82)
$`+{\displaystyle \frac{m_\pi ^2}{2}}\left[{\displaystyle \frac{x^2}{4}}\left(\mathrm{Tr}[B,\mathrm{\Pi }]^2\mathrm{cos}^2\alpha \mathrm{Tr}[X_1,\mathrm{\Pi }]^2\mathrm{sin}^2\alpha \right)+\mathrm{Tr}\mathrm{\Pi }^2\mathrm{cos}\alpha \right]+\mathrm{}.`$
As should be expected, the linear term vanishes due to (7) and we shall concentrate now on the quadratic term.
#### 9.2.1 $`\beta =1`$
At this point we need the explicit form of $`X_1`$ for $`\beta =1`$. According to (76) we have
$$X_2=i\mathrm{\Sigma }_d\mathrm{\Sigma }_c^{}=\left(\begin{array}{cc}0& I\\ I& 0\end{array}\right).(\beta =1)$$
(83)
For $`X_1`$ we therefore find
$$X_1=iBX_2=\left(\begin{array}{cc}0& iI\\ iI& 0\end{array}\right).(\beta =1)$$
(84)
We then use the block decomposition of $`\mathrm{\Pi }`$ (16) and find
$`_{\mathrm{st}}(\mathrm{\Sigma })=_{\mathrm{st}}(\mathrm{\Sigma }_\alpha )+m_\pi ^2\mathrm{Tr}[x^2(QQ^{}\mathrm{cos}^2\alpha (P_S^2+Q_RQ_R^{})\mathrm{sin}^2\alpha )`$
$`+(P^2+QQ^{})\mathrm{cos}\alpha ]+\mathrm{}.`$ $`(\beta =1)`$ (85)
We have introduced the following projections of $`P`$ and $`Q`$:
$`P_S={\displaystyle \frac{1}{2}}(P+IP^TI)\text{ and }P_A={\displaystyle \frac{1}{2}}(PIP^TI);`$
$`Q_R={\displaystyle \frac{1}{2}}(Q+IQ^{}I)\text{ and }Q_I={\displaystyle \frac{1}{2}}(QIQ^{}I).`$ $`(\beta =1)`$ (86)
These projections are orthogonal, $`\mathrm{Tr}P_SP_A=\mathrm{Re}\mathrm{Tr}Q_RQ_I^{}=0`$. Using this fact and the relation $`\mathrm{cos}\alpha =1/x^2`$ we obtain
$`_{\mathrm{st}}(\mathrm{\Sigma })=_{\mathrm{st}}(\mathrm{\Sigma }_\alpha )+m_\pi ^2x^2\mathrm{Tr}\left[Q_RQ_R^{}\mathrm{sin}^2\alpha +P_A^2\mathrm{cos}^2\alpha +P_S^2\right]+\mathrm{}.`$ $`(\beta =1)`$ (87)
From (87) we can now read off the curvatures for the different multiplets of the (pseudo-)Goldstone modes. We see that there is a true flat direction, $`Q_I`$, which describes true massless Goldstones of the diquark condensation phase. These fields are the phases (a single $`U(1)_B`$ phase for $`N_f=2`$ and a set of $`U(1)_B\times SU(N_f)_V/Sp(N_f)_V`$ phases for other $`N_f`$) of the diquark condensate. As we shall see in the next section, the linear derivative terms in the effective Lagrangian mix $`Q_I`$ and $`Q_R`$, thus the actual true Goldstone excitations are certain linear combinations of $`Q_I`$ and $`Q_R`$.
The degeneracies of the multiplets follow from the definitions (9.2.1) and are given by
$`N_{P_S}`$ $`=`$ $`{\displaystyle \frac{N_f(N_f+1)}{2}};N_{P_A}={\displaystyle \frac{N_f(N_f1)}{2}}1;`$
$`N_{Q_R}`$ $`=`$ $`N_{Q_I}={\displaystyle \frac{N_f(N_f1)}{2}}.(\beta =1)`$ (88)
They correspond to representations of the group $`Sp(N_f)`$, which is the residual symmetry remaining intact after spontaneous breaking. The number of the flat directions $`Q_I`$ exactly matches the number of the broken generators in $`SU(N_f)_V\times U(1)_BSp(N_f)_V`$.
#### 9.2.2 $`\beta =4`$
The derivation of the curvatures of the effective potential $`_{\mathrm{st}}`$ in the $`\beta =4`$ case follows the same lines as in the $`\beta =1`$ case. All the differences stem from the fact that $`\mathrm{\Sigma }`$ is a symmetric unitary matrix in this case. Consequently, $`\mathrm{\Sigma }_c`$ and $`\mathrm{\Sigma }_d`$ are given by (19) and (59) instead of (15) and (57). In particular, the generator which rotates $`\mathrm{\Sigma }_c`$ into $`\mathrm{\Sigma }_d`$ is now given by
$$X_2=i\mathrm{\Sigma }_d\mathrm{\Sigma }_c^{}=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),(\beta =4)$$
(89)
and the generator $`X_1`$, into which $`B`$ rotates and which is responsible for the splitting of both $`P`$ and $`Q`$ branches (in the block decomposition (16)), is given by
$$X_1=iBX_2=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right).(\beta =4)$$
(90)
It is now straightforward to rewrite the quadratic part of the Lagrangian (82) in terms of the $`PQ`$-block decomposition. The mass matrix becomes diagonal in terms of the following projections of $`P`$ and $`Q`$:
$`P_S={\displaystyle \frac{1}{2}}(P+P^T)\text{ and }P_A={\displaystyle \frac{1}{2}}(PP^T);`$
$`Q_R={\displaystyle \frac{1}{2}}(Q+Q^{})\text{ and }Q_I={\displaystyle \frac{1}{2}}(QQ^{}).`$ $`(\beta =4)`$ (91)
The static part of the effective Lagrangian in terms of these fields is given by
$`_{\mathrm{st}}(\mathrm{\Sigma })=_{\mathrm{st}}(\mathrm{\Sigma }_\alpha )+m_\pi ^2x^2\mathrm{Tr}\left[Q_RQ_R^{}\mathrm{sin}^2\alpha +P_S^2\mathrm{cos}^2\alpha +P_A^2\right]+\mathrm{},`$ $`(\beta =4)`$ (92)
which is similar to (87), but $`S`$ and $`A`$ labels exchange places.
The main difference is in the dimension of the degenerate multiplets which matches the representations of the corresponding residual symmetry group. By inspection one finds
$`N_{P_A}`$ $`=`$ $`{\displaystyle \frac{N_f(N_f1)}{2}};N_{P_S}={\displaystyle \frac{N_f(N_f+1)}{2}}1;`$
$`N_{Q_R}`$ $`=`$ $`N_{Q_I}={\displaystyle \frac{N_f(N_f+1)}{2}}.(\beta =4)`$ (93)
## 10 Spectrum
In order to complete our program and determine the spectrum of low-lying excitations we must take into account the derivative terms in the Lagrangian (55). They contribute in a non-trivial way to the quadratic form in the Goldstone fields. Due to the non-Lorentz invariant nature of the system we study at finite $`\mu `$, the dispersion laws<sup>4</sup><sup>4</sup>4In Euclidean field theory the dispersion relation is given by the poles of the propagator in the $`Eip_0`$-plane. do not have a simple form, $`E^2=𝒑^2+m^2`$. The mass, as measured on the lattice and given by the exponential fall-off of the propagator at large Euclidean time, is the value of $`E`$, i.e. $`ip_0`$, at the pole of the propagator at $`𝒑`$$`=0`$. These pole masses, or rest energies, we shall now evaluate.
### 10.1 Normal phase: $`\mu <m_\pi /2`$
Expanding the derivative terms in the Lagrangian (55) in the same way as we expanded the static part in the previous section we obtain
$$(\mathrm{\Sigma })=(\mathrm{\Sigma }_c)+\frac{1}{2}\mathrm{Tr}\left\{\left(_\nu \mathrm{\Pi }\mu [B_\nu ,\mathrm{\Pi }]\right)^2+m_\pi ^2\mathrm{\Pi }^2\right\}+\mathrm{},$$
(94)
where $`B_\nu =B\delta _{\nu 0}`$, as defined earlier in Section 6.2. Using the Fourier decomposition of $`\mathrm{\Pi }`$,
$$\mathrm{\Pi }(x)=\underset{p}{}\mathrm{\Pi }_p\mathrm{exp}(ipx)=\underset{p}{}\mathrm{\Pi }_p\mathrm{exp}(Ex_0i𝒑𝒙),$$
(95)
we find that the dispersion law has the generic form
$$(E+b\mu )^2=𝒑^2+m_\pi ^2,\text{ or }E=b\mu +\sqrt{𝒑^2+m_\pi ^2},$$
(96)
where $`b`$ is the baryon charge of the given excitation. This is in agreement with the fact that the effect of $`\mu `$ on each state is simply an energy shift by $`b\mu `$. In particular, the dispersion law of the $`P`$-type Goldstone modes, which carry no baryon charge, are not affected by $`\mu `$ at all. The rest energies of the diquarks $`Q`$ and antidiquarks $`Q^{}`$ are shifted according to their charges. To summarize, the dispersion laws in the normal phase are given by <sup>5</sup><sup>5</sup>5Note that though the mass of the diquarks $`Q`$ vanishes at the transition point, $`\mu =m_\pi /2`$, the dispersion law is not linear, but quadratic, $`E𝒑^2`$. This is related to the well-known critical slowing down at a second order phase transition.
$`P`$ $`:`$ $`E=\left(m_\pi ^2+𝒑^2\right)^{1/2};`$
$`Q^{}`$ $`:`$ $`E=\left(m_\pi ^2+𝒑^2\right)^{1/2}+2\mu ;`$
$`Q`$ $`:`$ $`E=\left(m_\pi ^2+𝒑^2\right)^{1/2}2\mu .`$ (97)
The masses of these excitations as well as their degeneracies are given in Tables 1 and 2. A schematic picture of the mass dependence is shown in Figures 1 and 2.
### 10.2 Diquark condensation phase: $`\mu >m_\pi /2`$.
In this phase the ground state changes and we need to expand around the rotated value of the condensate $`\mathrm{\Sigma }=\mathrm{\Sigma }_\alpha `$. We find<sup>6</sup><sup>6</sup>6A useful observation which helps to write this in such a compact form is that $`\mathrm{Tr}[X_1,\mathrm{\Pi }]_0\mathrm{\Pi }=0`$ for any $`\mathrm{\Pi }`$. This can be checked explicitly using the block decomposition of $`\mathrm{\Pi }`$, but it is easier to see that once you realize that the commutator of two $`X`$-like generators is a $`T`$-like generator, and its trace with another $`X`$-like generator, such as $`_0\mathrm{\Pi }`$ is zero.
$`(\mathrm{\Sigma })=(\mathrm{\Sigma }_\alpha )`$ $`+`$ $`{\displaystyle \frac{1}{2}}\mathrm{Tr}\{(_\nu \mathrm{\Pi }\mu [V_\alpha B_\nu V_\alpha ^{},\mathrm{\Pi }])(_\nu \mathrm{\Pi }\mu [V_\alpha ^{}B_\nu V_\alpha ,\mathrm{\Pi }])`$ (98)
$`+`$ $`m_\pi ^2\mathrm{\Pi }^2\mathrm{cos}\alpha \}+\mathrm{}.`$
Expanding the product we find
$$(\mathrm{\Sigma })=\frac{1}{2}\mathrm{Tr}\left\{\left(_\nu \mathrm{\Pi }\right)^22\mu \mathrm{cos}\alpha [B,\mathrm{\Pi }]_0\mathrm{\Pi }\right\}+_{\mathrm{stat}}(\mathrm{\Sigma })+\mathrm{}$$
(99)
with the static part given by (82). We observe, first of all, that the linear derivative term contains only the charged fields, i.e. $`Q`$ and $`Q^{}`$. The dispersion laws of the $`P`$ fields remain unaffected by the linear term. It remains Lorentz invariant, with mass given by the curvature of the static part of the Lagrangian. In order to determine the dispersion laws for the $`Q`$ and $`Q^{}`$ fields we need to solve a secular equation obtained by substituting Fourier decomposition of $`Q`$’s into
$`(\mathrm{\Sigma })`$ $`=`$ $`(\mathrm{\Sigma }_\alpha )+\mathrm{Tr}\{(_\nu Q_R^{}_\nu Q_R+_\nu Q_I^{}_\nu Q_I)`$ (100)
$`4\mu \mathrm{cos}\alpha (Q_I^{}_0Q_R+Q_R^{}_0Q_I)+4\mu ^2Q_R^{}Q_R\mathrm{sin}^2\alpha \}+\mathrm{},`$
where we only wrote the $`Q`$-dependent terms. This expression is the same for either $`\beta =1`$ or $`\beta =4`$; only the definition of $`Q`$’s is different: eqs. (9.2.1) or eqs. (9.2.2). The secular equation has the form
$$det\left(\begin{array}{cc}E^2𝒑^2& 4\mu E\mathrm{cos}\alpha \\ 4\mu E\mathrm{cos}\alpha & E^2𝒑^24\mu ^2\mathrm{sin}^2\alpha \end{array}\right)=0.$$
(101)
We see that $`Q_R`$ and $`Q_I`$ are mixed by the linear derivative term. The mixing is maximal when $`\alpha =0`$, and vanishes when $`\alpha =\pi /2`$, i.e. when $`\mu m_\pi `$. For any $`\alpha `$ there is a solution for which the rest energy $`E(0)=0`$ — the Goldstone boson branch, which we denote $`\stackrel{~}{Q}`$. For $`\alpha =0`$ it is given by $`\stackrel{~}{Q}=Q=Q_R+Q_I`$ (at rest, $`𝒑=0`$). For $`\alpha =\pi /2`$ it is entirely $`Q_I`$. The other solution of the secular equation, the linear combination of $`Q_R`$ and $`Q_I`$ orthogonal to $`\stackrel{~}{Q}`$, is massive. We denote it by $`\stackrel{~}{Q}^{}`$.
The dispersion laws for the meson multiplets $`P_S`$, $`P_A`$ and the mixed diquark-antidiquark multiplets $`\stackrel{~}{Q}`$ and $`\stackrel{~}{Q}^{}`$ are given by (in the $`\beta =1`$ case) the following relations:
$`P_S`$ $`:`$ $`E^2=𝒑^2+m_\pi ^2x^2;`$
$`P_A`$ $`:`$ $`E^2=𝒑^2+m_\pi ^2x^2\mathrm{cos}^2\alpha ;`$
$`\stackrel{~}{Q}^{}`$ $`:`$ $`E^2=𝒑^2+2\mu ^2(1+3\mathrm{cos}^2\alpha )+2\mu \sqrt{\mu ^2(1+3\mathrm{cos}^2\alpha )^2+4𝒑^2\mathrm{cos}^2\alpha };`$
$`\stackrel{~}{Q}`$ $`:`$ $`E^2=𝒑^2+2\mu ^2(1+3\mathrm{cos}^2\alpha )2\mu \sqrt{\mu ^2(1+3\mathrm{cos}^2\alpha )^2+4𝒑^2\mathrm{cos}^2\alpha }.(\beta =1)`$ (102)
The only difference in the dispersion laws for $`\beta =4`$ is the interchange of $`P_S`$ and $`P_A`$, and also that the degeneracy of the multiplets as given by (9.2.1) and (9.2.2) is different. Note that the dispersion relation for the lowest lying mode, i.e. $`\stackrel{~}{Q}`$, is linear: $`Ep`$. The slope is a function of $`\mu `$: it vanishes at the transition point, $`x=1`$, and approaches 1 for large $`\mu `$. The linear slope of low energy excitations is characteristic of superfluidity.
The pole masses of the excitations, i.e. the position of the pole at zero momentum,$`E(0)`$, are given in Table 1 for $`\beta =1`$ and Table 2 for $`\beta =4`$. The representations of the residual symmetry groups are denoted by their Young diagrams. Because mesons are quark-quark or quark-antiquark pairs they transform as rank two tensors (unless they are singlets). The representations can be uniquely identified by their dimensions found in (17), (9.2.1) and (20), (9.2.2). An explicit form of the representations is given in the Appendix. Figures 1 and 2 show schematic pictures these results, together with the residual symmetry groups in the different phases. In particular, one observes that the spectrum is continuous at the transition point $`x=1`$.
## 11 Diquark source
In three-color QCD the diquark condensate is not a color-singlet, therefore, one cannot study the phenomenon of diquark condensation by applying an external diquark source: such a source term will not be gauge invariant. In the theories which we study in this paper, e.g., in two-color QCD, the diquark condensate is colorless. Therefore one can add a gauge-invariant source $`j`$ to the theory. This source plays a role similar to the role the quark mass $`m`$ plays with respect to the chiral condensate. Such a non-zero source term $`j`$ is in fact used in lattice simulations in order to measure the value of the diquark condensate in the limit $`j0`$. In this section we show how the results of the previous sections are modified in the presence of a non-zero $`j`$.
### 11.1 $`\beta =1`$
Let us first rewrite the scalar diquark source term in the macroscopic theory using our $`SU(2N_f)`$ spinor notations. In the notations of section 3 we find
$`i{\displaystyle \frac{j}{2}}\psi ^TC\gamma _5\tau _2I\psi +\mathrm{h}.\mathrm{c}.={\displaystyle \frac{j}{2}}\mathrm{\Psi }^T\sigma _2\tau _2\left(\begin{array}{cc}iI& 0\\ 0& iI\end{array}\right)\mathrm{\Psi }+\mathrm{h}.\mathrm{c}.{\displaystyle \frac{1}{2}}\mathrm{\Psi }^T\sigma _2\tau _2J\mathrm{\Psi }+\mathrm{h}.\mathrm{c}.,`$ (105)
$$J=j\widehat{J}\mathrm{and}\widehat{J}=\left(\begin{array}{cc}iI& 0\\ 0& iI\end{array}\right),(\beta =1)$$
(106)
the antisymmetric matrix $`I`$ is defined in (58), and the summation over the $`N_f`$ flavor indices has been suppressed. Comparing with the quark mass term (22) we find that the two belong to the same multiplet (i.e. they transform into one another) under the $`SU(2N_f)`$ rotations.<sup>7</sup><sup>7</sup>7There is a freedom in the choice of the orientation of the diquark source $`\widehat{J}`$. It is precisely the same freedom as the one corresponding to the $`SU(N_f)\times U(1)/Sp(N_f)`$ degeneracy of the diquark condensation vacuum $`\mathrm{\Sigma }_d`$. An infinitesimal diquark source $`j`$ determines the vacuum orientation of $`\mathrm{\Sigma }_d`$. Our choices reflect this fact. Indeed $`\mathrm{\Sigma }_d=\widehat{J}^{}`$. We can write the sum of the diquark source and the bare mass term as
$$m\overline{\psi }\psi i\frac{j}{2}(\psi ^TC\gamma _5\tau _2I\psi +\mathrm{h}.\mathrm{c}.)=\frac{1}{2}\mathrm{\Psi }^T\sigma _2\tau _2M_\varphi \mathrm{\Psi },(\beta =1)$$
(107)
where
$$M_\varphi =m\widehat{M}+j\widehat{J}=\sqrt{m^2+j^2}(\widehat{M}\mathrm{cos}\varphi +\widehat{J}\mathrm{sin}\varphi )=\sqrt{m^2+j^2}\widehat{M}_\varphi .$$
(108)
We have used a notation similar to (60) with
$$\mathrm{tan}\varphi =\frac{j}{m}.$$
(109)
It is easy to see how the diquark source term modifies the effective Lagrangian (55): we need to replace the mass matrix $`mG\widehat{M}`$ in (29) by $`\sqrt{m^2+j^2}G\widehat{M}_\varphi `$. The Gell-Mann$``$Oakes$``$Renner relation at $`\mu =0`$ becomes
$$m_\pi ^2=\frac{\sqrt{m^2+j^2}G}{F^2}.$$
(110)
Using this fact we can write the effective Lagrangian for the theory with both mass term and diquark source (compare to (55)) as
$$_{\mathrm{eff}}(\mathrm{\Sigma })=\frac{F^2}{2}\left[\mathrm{Tr}_\nu \mathrm{\Sigma }_\nu \mathrm{\Sigma }^{}2m_\pi ^2\mathrm{ReTr}(\widehat{M}_\varphi \mathrm{\Sigma })\right].$$
(111)
Now we can repeat the steps we performed in section 9, but for the Lagrangian (111).
Since the source term we introduced favors the direction of $`\mathrm{\Sigma }_d`$, we expect that the minimum of the Lagrangian (111), $`\overline{\mathrm{\Sigma }}`$, is again given by a linear combination of $`\mathrm{\Sigma }_c`$ and $`\mathrm{\Sigma }_d`$ as in (60),
$`\overline{\mathrm{\Sigma }}=\mathrm{\Sigma }_\alpha =\mathrm{\Sigma }_c\mathrm{cos}\alpha +\mathrm{\Sigma }_d\mathrm{sin}\alpha =V_\alpha \mathrm{\Sigma }_cV_\alpha ^T,\mathrm{where}V_\alpha =e^{i\alpha X_2},`$ (112)
and $`X_2`$ is the generator that rotates $`\mathrm{\Sigma }_c`$ into $`\mathrm{\Sigma }_d`$ as before. The value of $`\alpha `$, as determined by the saddle-point equations, now depends on the value of the diquark source. The proof of Section 8 is not easily extendable to this case. We shall continue on the assumption that the minimum is global. We shall prove, however, that it is a minimum when we expand to second order in fluctuations of $`\mathrm{\Sigma }`$ and find no linear terms and a positive quadratic form. If we substitute this ansatz into the effective Lagrangian (111) we find
$`_{\mathrm{st}}(\mathrm{\Sigma }_\alpha )=F^2m_\pi ^2N_f\left[{\displaystyle \frac{x^2}{2}}(\mathrm{cos}2\alpha 1)2\mathrm{cos}(\alpha \varphi )\right].`$ (113)
The dependence of the angle $`\alpha `$ on $`x`$ follows from minimizing this Lagrangian and is given by
$$x^2\mathrm{cos}\alpha \mathrm{sin}\alpha =\mathrm{sin}(\alpha \varphi ),$$
(114)
which is different from (7), but coincides with it in the limit $`\varphi 0`$, of course. For finite $`\varphi `$ the angle $`\alpha `$ is already non-zero ($`\alpha =\varphi `$) at $`x=0`$ as it should be, since a diquark source drives a non-zero diquark condensate. At $`x=\mathrm{}`$, the value of $`\alpha `$ is $`\pi /2`$ independently of $`\varphi `$. The major difference from (7) is that the dependence of $`\alpha `$ on $`x`$ is analytic. There is no phase transition as a function of $`\mu `$ when the diquark source $`j`$ is nonzero. This is to be expected: the diquark source plays the role an external magnetic field plays in a ferromagnet. The diquark source explicitly breaks $`SU(N_f)_V\times U(1)_B`$ symmetry down to $`Sp(N_f)_V`$ and there is only one phase (the one with residual $`Sp(N_f)_V`$ symmetry) for all values of $`x`$.
Our next step is to expand in powers of fluctuations of $`\mathrm{\Sigma }`$ around $`\overline{\mathrm{\Sigma }}`$. This can be done following the algebra in section 9 with minimal modifications. The Goldstone manifold is again parameterized by
$`\mathrm{\Sigma }=V_\alpha U\mathrm{\Sigma }_cU^TV_\alpha ^T,`$ (115)
but the second identity in (9.2) is now given by
$`V_\alpha ^T\widehat{M}_\varphi V_\alpha =\left(V_\alpha \widehat{M}_\varphi ^{}V_\alpha ^T\right)^{}=\mathrm{\Sigma }_{\varphi \alpha }^{}=\mathrm{\Sigma }_c^{}\mathrm{cos}(\alpha \varphi )\mathrm{\Sigma }_d^{}\mathrm{sin}(\alpha \varphi ),`$ (116)
whereas the first identity is the same as before. Using this algebra, the effective Lagrangian to second order in the pion fields can be expressed as
$`(\mathrm{\Sigma })`$ $`=`$ $`_{\mathrm{st}}(\mathrm{\Sigma }_\alpha )+Fm_\pi ^2\left[x^2\mathrm{cos}\alpha \mathrm{sin}\alpha \mathrm{sin}(\alpha \varphi )\right]\mathrm{Tr}(X_2\mathrm{\Pi })`$ (117)
$`+`$ $`{\displaystyle \frac{1}{2}}\mathrm{Tr}\left\{\left(_\nu \mathrm{\Pi }\right)^22\mu \mathrm{cos}\alpha [B,\mathrm{\Pi }]_0\mathrm{\Pi }\right\}`$
$`+`$ $`{\displaystyle \frac{m_\pi ^2}{2}}\left[{\displaystyle \frac{x^2}{4}}\left(\mathrm{Tr}[B,\mathrm{\Pi }]^2\mathrm{cos}^2\alpha \mathrm{Tr}[X_1,\mathrm{\Pi }]^2\mathrm{sin}^2\alpha \right)+\mathrm{Tr}\mathrm{\Pi }^2\mathrm{cos}(\alpha \varphi )\right]+\mathrm{}`$
Inspecting the linear terms we find that they vanish if $`\alpha `$ is chosen according to (114). This means that our ansatz (112) is indeed an extremum.
In order to find the mass spectrum we have to study the term of second order in the $`\mathrm{\Pi }`$-fields in the Lagrangian (117), and determine the dispersion laws of the different (pseudo-)Goldstone fields. Using block decomposition of the generators $`\mathrm{\Pi }`$ (16), projections (9.2.1) of $`P`$ and $`Q`$ and the relationship (114) between $`x`$ and $`\alpha `$ we find
$`(\mathrm{\Sigma })`$ $`=`$ $`_{\mathrm{st}}(\mathrm{\Sigma }_\alpha )+\mathrm{Tr}\left[\left(_\nu Q_R^{}_\nu Q_R+_\nu Q_I^{}_\nu Q_I\right)4\mu \mathrm{cos}\alpha \left(Q_I^{}_0Q_R+Q_R^{}_0Q_I\right)\right]`$ (118)
$`+m_\pi ^2\mathrm{Tr}\left[Q_IQ_I^{}{\displaystyle \frac{\mathrm{sin}\varphi }{\mathrm{sin}\alpha }}+Q_RQ_R^{}\left(x^2\mathrm{sin}^2\alpha +{\displaystyle \frac{\mathrm{sin}\varphi }{\mathrm{sin}\alpha }}\right)\right]`$
$`+\mathrm{Tr}\left[_\nu P_A_\nu P_A+P_A^2m_\pi ^2\left(x^2\mathrm{cos}^2\alpha +{\displaystyle \frac{\mathrm{sin}\varphi }{\mathrm{sin}\alpha }}\right)\right]`$
$`+\mathrm{Tr}\left[_\nu P_S_\nu P_S+P_S^2m_\pi ^2\left(x^2+{\displaystyle \frac{\mathrm{sin}\varphi }{\mathrm{sin}\alpha }}\right)\right]+\mathrm{}.(\beta =1)`$
As in the case of a vanishing diquark source studied in the previous Section, the linear derivative term contains only the charged fields $`Q`$ and $`Q^{}`$. Therefore the dispersion relations for the neutral fields $`P`$ are not affected by this term: the masses of the $`P`$ fields are given by the curvature of the static part of the Lagrangian at the minimum. However the linear derivative term mixes the charged Goldstones. In order to obtain the dispersion laws for the $`Q`$ fields, as in Section 10, we have to solve a secular equation obtained by substituting the Fourier decomposition of these fields into the effective Lagrangian (118). The secular equation reads
$$det\left(\begin{array}{cc}E^2𝒑^2m_\pi ^2\frac{\mathrm{sin}\varphi }{\mathrm{sin}\alpha }& 4\mu E\mathrm{cos}\alpha \\ 4\mu E\mathrm{cos}\alpha & E^2𝒑^24\mu ^2\mathrm{sin}^2\alpha m_\pi ^2\frac{\mathrm{sin}\varphi }{\mathrm{sin}\alpha }\end{array}\right)=0.$$
(119)
The dispersion laws for the different Goldstone modes are therefore found to be given by
$`P_S`$ $`:`$ $`E^2=𝒑^2+m_\pi ^2\left(x^2+{\displaystyle \frac{\mathrm{sin}\varphi }{\mathrm{sin}\alpha }}\right);`$
$`P_A`$ $`:`$ $`E^2=𝒑^2+m_\pi ^2\left(x^2\mathrm{cos}^2\alpha +{\displaystyle \frac{\mathrm{sin}\varphi }{\mathrm{sin}\alpha }}\right);`$
$`\stackrel{~}{Q}^{}`$ $`:`$ $`E^2=𝒑^2+m_\pi ^2{\displaystyle \frac{\mathrm{sin}\varphi }{\mathrm{sin}\alpha }}+2\mu ^2(1+3\mathrm{cos}^2\alpha )`$
$`+2\mu \sqrt{\mu ^2(1+3\mathrm{cos}^2\alpha )^2+4\mathrm{cos}^2\alpha (𝒑^2+m_\pi ^2{\displaystyle \frac{\mathrm{sin}\varphi }{\mathrm{sin}\alpha }})};`$
$`\stackrel{~}{Q}`$ $`:`$ $`E^2=𝒑^2+m_\pi ^2{\displaystyle \frac{\mathrm{sin}\varphi }{\mathrm{sin}\alpha }}+2\mu ^2(1+3\mathrm{cos}^2\alpha )`$ (120)
$`2\mu \sqrt{\mu ^2(1+3\mathrm{cos}^2\alpha )^2+4\mathrm{cos}^2\alpha (𝒑^2+m_\pi ^2{\displaystyle \frac{\mathrm{sin}\varphi }{\mathrm{sin}\alpha }})};(\beta =1),`$
where $`\alpha `$ is an implicit function of $`x`$ and $`\varphi `$ given by (114). Note that the $`\stackrel{~}{Q}`$ are no longer true Goldstone modes, since the symmetry, which is broken spontaneously in the diquark condensation phase $`x>1`$ at $`j=0`$, is now broken explicitly by the diquark source term. The masses of the different multiplets are given by the value of $`E`$ at $`𝒑=0`$. The positivity of all masses shows that the minimum given by the saddle-point equation (114) is at least a local minimum. This mass spectrum at a small non-zero diquark source is depicted in Fig. 3.
### 11.2 $`\beta =4`$
In the $`\beta =4`$ case the diquark source term has the form<sup>8</sup><sup>8</sup>8As before, our previous choice of $`\mathrm{\Sigma }_d`$ corresponds to our choice of $`J`$, i.e., $`\mathrm{\Sigma }_d=\widehat{J}^{}`$.
$`i{\displaystyle \frac{j}{2}}\psi ^TC\gamma _5\psi +\mathrm{h}.\mathrm{c}.={\displaystyle \frac{j}{2}}\mathrm{\Psi }^T\sigma _2\left(\begin{array}{cc}i& 0\\ 0& i\end{array}\right)\mathrm{\Psi }+\mathrm{h}.\mathrm{c}.{\displaystyle \frac{1}{2}}\mathrm{\Psi }^T\sigma _2J\mathrm{\Psi }+\mathrm{h}.\mathrm{c}.,`$ (123)
$$J=j\widehat{J}\mathrm{and}\widehat{J}=\left(\begin{array}{cc}i& 0\\ 0& i\end{array}\right)(\beta =4)$$
(124)
The analysis follows the same lines as in the $`\beta =1`$ case. We again introduce the combined mass matrix
$`M_\varphi =m\widehat{M}+j\widehat{J}=\sqrt{m^2+j^2}\widehat{M}_\varphi ,`$ (125)
with $`\widehat{M}`$ and $`\widehat{J}`$ defined for $`\beta =4`$ in (33) and (123). In the expansion of the Lagrangian to second order in the $`\mathrm{\Pi }`$-fields for $`\beta =1`$ we have only used the general commutation properties of $`\widehat{M}_\varphi `$ which, with a proper redefinition of the generators $`X_1`$ and $`X_2`$, are the same in the case $`\beta =4`$. We thus obtain the same second order Lagrangian (117) as for $`\beta =1`$ resulting in the same saddle point equation (114) and the same mean field value of the free energy. However, the mass spectrum is slightly different. We express the quadratic part of the Lagrangian in terms of the block-decomposition (16) and the projections (9.2.2). The final result is
$`(\mathrm{\Sigma })`$ $`=`$ $`_{\mathrm{st}}(\mathrm{\Sigma }_\alpha )+\mathrm{Tr}\left[\left(_\nu Q_R^{}_\nu Q_R+_\nu Q_I^{}_\nu Q_I\right)4\mu \mathrm{cos}\alpha \left(Q_I^{}_0Q_R+Q_R^{}_0Q_I\right)\right]`$ (126)
$`+m_\pi ^2\mathrm{Tr}\left[Q_IQ_I^{}{\displaystyle \frac{\mathrm{sin}\varphi }{\mathrm{sin}\alpha }}+Q_RQ_R^{}\left(x^2\mathrm{sin}^2\alpha +{\displaystyle \frac{\mathrm{sin}\varphi }{\mathrm{sin}\alpha }}\right)\right]`$
$`+\mathrm{Tr}\left[_\nu P_S_\nu P_S+P_S^2m_\pi ^2\left(x^2\mathrm{cos}^2\alpha +{\displaystyle \frac{\mathrm{sin}\varphi }{\mathrm{sin}\alpha }}\right)\right]`$
$`+\mathrm{Tr}\left[_\nu P_A_\nu P_A+P_A^2m_\pi ^2\left(x^2+{\displaystyle \frac{\mathrm{sin}\varphi }{\mathrm{sin}\alpha }}\right)\right]+\mathrm{}.(\beta =4)`$
As in the $`\beta =1`$ case, the linear derivative term mixes the $`Q`$ and $`Q^{}`$ fields. In order to get the dispersion relations of these fields, the same secular equation as before (119) has to be solved. The dispersion laws are the same as in (11.1) with the familiar exchange of labels $`SA`$, and the degeneracy of the multiplets given by (9.2.2) instead of (9.2.1). The corresponding mass spectrum at a small non-zero $`j`$ is shown in Fig. 3.
## 12 Condensate and baryon charge densities
In this section we derive the dependence of the chiral condensate, the diquark condensate and the baryon charge densities. These densities are obtained by differentiating the vacuum energy density with respect to $`m`$, $`j`$ and $`\mu `$,
$$\overline{\psi }\psi =\frac{_{\mathrm{vac}}}{m};\psi \psi =\frac{_{\mathrm{vac}}}{j};n_B=\frac{_{\mathrm{vac}}}{\mu }.$$
(127)
The shorthand symbol $`\psi \psi `$ denotes the l.h.s. of the equation (105) in the $`\beta =1`$ case and of (123) in the $`\beta =4`$ case without factor $`j`$.
In the effective theory the vacuum energy density is given by the value of the effective Lagrangian (55) at the minimum,
$`_{\mathrm{vac}}=_{\mathrm{eff}}(\mathrm{\Sigma }_\alpha )`$ $`=`$ $`{\displaystyle \frac{F^2}{2}}\left[2\mu ^2\mathrm{Tr}\left(\mathrm{\Sigma }_\alpha B^T\mathrm{\Sigma }_\alpha ^{}B+BB\right)2m_\pi ^2\mathrm{ReTr}\left(\mathrm{\Sigma }_\varphi ^{}\mathrm{\Sigma }_\alpha \right)\right]`$ (128)
$`=4N_fF^2\mu ^2\mathrm{sin}^2\alpha 2N_fG(m\mathrm{cos}\alpha +j\mathrm{sin}\alpha ),`$
where we used (109) and (110). Differentiating<sup>9</sup><sup>9</sup>9The implicit dependence of $`\alpha `$ on $`m`$, $`j`$ and $`\mu `$ is not contributing because $`_{\mathrm{vac}}/\alpha =0`$., we find
$$\overline{\psi }\psi =2N_fG\mathrm{cos}\alpha ;\psi \psi =2N_fG\mathrm{sin}\alpha ;n_B=8N_fF^2\mu \mathrm{sin}^2\alpha ,$$
(129)
where the angle of the rotation of the condensate, $`\alpha `$, depends implicitly on $`\mu `$, $`m`$ and $`j`$ through the solution of (114). Note, in particular, that the sum $`\overline{\psi }\psi ^2+\psi \psi ^2`$ does not depend on $`\mu `$, $`m`$ or $`j`$ (to the order in Chiral Perturbation Theory we are working), which reflects the fact that the condensate rotates.
For $`j=0`$ the dependence of $`\alpha `$ on $`\mu `$ is simple: $`\alpha =0`$, for $`\mu <m_\pi /2`$, and $`\alpha =\mathrm{arccos}[m_\pi ^2/(4\mu ^2)]`$ otherwise (7). That means the densities (129) are constant for $`\mu <m_\pi /2`$, as they should since the vacuum state does not change until $`\mu `$ reaches $`m_\pi /2`$. The results for the condensate and baryon charge densities are summarized in Table 3.
For $`j0`$ the angle $`\alpha `$ varies between the values $`\varphi `$ and $`\pi /2`$ as a function of $`\mu `$ and $`m`$, determined by the saddle point equation (114) $`x^2\mathrm{sin}\alpha \mathrm{cos}\alpha =\mathrm{sin}(\alpha \varphi )`$. The saddle point is a quartic equation which, in principle, can be solved analytically for arbitrary $`j`$. The dependence of the densities (129) on $`\mu `$ is shown in Fig. 4 for $`j=0`$ and in Fig. 5 for $`j0`$.
## 13 Equation of state of the non-ideal diquark gas
We now take a different look at the system we study (quarks with two colors or any-color adjoint quarks) near the threshold, $`\mu m_\pi /2`$. First of all, we observe, that for $`\mu `$ very close to $`m_\pi /2`$ the density of baryons is very (arbitrarily) small. This means that we can describe the system by an almost ideal gas. It should be a Bose gas, since the baryons, i.e. diquarks, are bosons. The equation of state of a dilute non-ideal Bose gas is a textbook problem . The interaction between the bosons is crucial here (the dependence on its strength is non-analytic, e.g. the attractive Bose gas is unstable). The baryons in our system repel each other, as we shall see. Moreover, we shall also see that the strength of this repulsion is exactly the one that gives the correct equation of state (129), (7) at small densities:
$$n_B=32N_fF^2\left(\mu \frac{m_\pi }{2}\right)+\mathrm{}.$$
(130)
This is a very non-trivial check of the consistency of our effective Lagrangian!
Since we are working at very small densities, we can expand the Lagrangian (29) up to next (quartic order) in the Goldstone fields. The quartic terms will describe the interaction energy of the Goldstones due to two-body scattering,
$$_{\mathrm{int}}=\frac{1}{24F^2}\mathrm{Tr}[\mathrm{\Pi },_\mu \mathrm{\Pi }][\mathrm{\Pi },^\mu \mathrm{\Pi }]\frac{m_\pi ^2}{24F^2}\mathrm{Tr}\mathrm{\Pi }^4+\mathrm{}.$$
(131)
We are now working in Minkowski space and the extra minus sign in front of the first term, compared to (29), is a consequence of this. The presence of the time derivatives should not confuse the reader. Since we are dealing with a non-relativistic system, $`\mu m_\pi /2m_\pi `$, we only need the leading term in the time dependence of the relativistic fields, i.e. $`_0\mathrm{\Pi }\pm im_\pi \mathrm{\Pi }`$. The spatial derivatives give only a subleading contribution. The dominant effect is that of the $`s`$-wave amplitude.
We now consider the ground state of such a system with the baryon density fixed (unlike before, when we fixed $`\mu `$). Most of the particles will occupy the lowest energy, i.e., zero-momentum, level (the population of the excited levels is the next order effect). These particles form a condensate, whose amplitude we denote by $`\mathrm{\Pi }`$. The energy of the system is different from zero only because of the interaction, and, as a function of $`\mathrm{\Pi }`$, it is given to us by (131). Looking at the block decomposition of $`\mathrm{\Pi }`$ (16) we note that only $`Q`$ and $`Q^{}`$ (diquark and antidiquark) fields carry non-zero baryon number. Therefore, we set $`P=0`$ in our case. Thus, from (131), we find for the dependence of the energy on the magnitude of the $`Q`$ condensate
$$_{\mathrm{int}}=\frac{m_\pi ^2}{4F^2}\mathrm{Tr}(QQ^{})^2+\mathrm{},$$
(132)
where we neglected subleading non-relativistic corrections.
Now we need to calculate $`n_B`$ as a function of $`Q`$. Since $`Q`$ carries baryon charge 2, we can write<sup>10</sup><sup>10</sup>10Another way to see (133) is to use the fact that $`n_B=/\mu `$. Thus, from (55) we find that $`n_B=2F^2\mathrm{Tr}B\mathrm{\Sigma }^{}_0\mathrm{\Sigma }=2i\mathrm{Tr}B\mathrm{\Pi }_0\mathrm{\Pi }`$, which coincides with (133).
$$n_B=2\mathrm{Tr}(Q^{}i_0QQi_0Q^{})=4m_\pi \mathrm{Tr}QQ^{}+\mathrm{},$$
(133)
where we neglected subleading corrections. Since all flavors of diquarks are equivalent<sup>11</sup><sup>11</sup>11There is only one flavor of diquarks for $`\beta =1`$, $`N_f=2`$. For larger $`N_f`$ the flavor $`SU(N_f)`$ symmetry ensures the equivalence of all diquarks. the matrix $`Q`$ is proportional to $`I`$ for $`\beta =1`$, or to $`1`$ for $`\beta =4`$. Thus we can relate traces in (132) and (133) by $`\mathrm{Tr}(QQ^{})^2=(\mathrm{Tr}QQ^{})^2/N_f`$. Therefore, we find for the dependence on $`n_B`$ of the vacuum energy at fixed $`n_B`$,
$$_{\mathrm{vac}}=\frac{n_B}{2}m_\pi +_{\mathrm{int}}=\frac{n_B}{2}m_\pi +\frac{n_B^2}{64F^2N_f}+\mathrm{}.$$
(134)
The first term is the rest energy of the diquarks. This determines the equation of state
$$\mu \frac{_{\mathrm{vac}}}{n_B}=\frac{m_\pi }{2}+\frac{n_B}{32F^2N_f}+\mathrm{}.$$
(135)
The fact that the equations of state (130) and (135) are identical is a very nontrivial property of the effective Lagrangian (55). It is intimately related to the symmetries (global and local) which determine the form of this Lagrangian.
This calculation also shows that the diquarks Bose-condense , and thus form a superfluid. Since the diquarks are charged we can call this phase a superconducting phase.
## 14 Conclusions and discussion
In this paper we have derived the low-energy effective Lagrangian for certain QCD-like theories at finite chemical potential $`\mu `$. It contains the same number of phenomenological parameters as at $`\mu =0`$: the pion decay constant $`F`$ and the vacuum value of the chiral condensate $`\overline{\psi }\psi _0`$. Using this effective Lagrangian we have completely determined its low-energy properties. In particular, we have established the dependence of the ground state, the condensates, and the masses of the excitations on the chemical potential $`\mu `$, the bare quark mass $`m`$ and the diquark source $`j`$. The theories to which our analysis applies include: (i) two-color QCD with quarks in fundamental color representation, and (ii) any-color QCD with quarks in adjoint color representation. The unifying feature of these theories is that the low-energy excitations, i.e. the Goldstone bosons of spontaneously broken symmetries, include baryons in the form of diquark states. Since the chiral Lagrangian describes the dynamics of such low-energy excitations, we can use the effective theory to describe phenomena associated with the condensation of diquarks.
Since we are working only to the lowest order in Chiral Perturbation Theory, we must remember that our approach has a limited region of applicability. First of all, the momenta of the particles, described by the Lagrangian (55) must be smaller than the scale $`M`$ of the masses of the non-Goldstone excitations, such as vector mesons, for example. In other words, the expansion parameter of the chiral perturbation theory is $`p/M`$. The bare quark mass is also included to leading order in $`m`$. It is obvious from (55) that $`mm_\pi ^2p^2`$ in the momentum power counting. The chemical potential $`\mu `$ is an external parameter with the dimension of mass. The interesting behavior of our theory is at values of $`\mu m_\pi p`$, and $`\mu `$ is counted as order $`p`$ in the momentum power counting. Our effective theory thus applies as long as $`m_\pi `$ and $`\mu `$ are much smaller than $`M`$.
Our effective theory predicts a phase transition at a value of $`\mu =m_\pi /2`$. It is important to note that this transition occurs within the domain of applicability of the effective theory, as long as $`m_\pi `$ is small, as discussed above, unlike, for example, the transition in QCD with three colors. The transition at $`\mu =m_\pi /2`$ is a robust prediction of our effective theory. The existence of such a transition is confirmed by lattice Monte Carlo simulations .
The fact that our effective Lagrangian (55) predicts the transition at a value of $`\mu `$ equal to $`m_\pi /2`$ is a direct consequence of the symmetries of the theory. The coefficient of the $`\mu `$-dependent term in the effective Lagrangian, responsible for the transition, is not an independent phenomenological parameter, but is fixed rigidly by the local symmetry . This symmetry corresponds to the conservation of the baryon charge (among other generators of $`SU(2N_f)`$). This ensures that the charge of a composite state in the effective theory is an additive sum of the charges of the constituents (non-renormalization). Indeed, the factor two in the relation $`\mu _0=m_\pi /2`$ is the baryon charge of the diquark. This is in agreement with the value of $`\mu `$, below which no transition can occur at zero temperature, given by the minimum value of the mass per baryon number among all baryons in the theory (see, e.g., ). In three-color QCD this would correspond to a value of $`\mu _0`$ approximately equal to 1/3 of the nucleon mass (offset by the binding energy of nuclear matter).<sup>12</sup><sup>12</sup>12We remind the reader that we measure the baryon charge in units of the $`U(1)_B`$ charge of a single quark.
Another way of looking at the $`\mu `$-dependence in the theory is the following. In the microscopic theory, given by the Lagrangian (49), $`\mu `$ enters as (a timelike component of) an Abelian gauge potential. As such, it can be completely removed from the Lagrangian by the time-dependent gauge transformation
$$\psi e^{\mu \tau }\psi \text{ and }\psi ^{}e^{\mu \tau }\psi ^{},$$
(136)
where $`\tau `$ is the Euclidean time<sup>13</sup><sup>13</sup>13 Note that, in Euclidean formulation, the $`\psi `$ and $`\psi ^{}`$ are independent variables and the global flavor symmetry group is in fact $`Gl(2N_f)`$ .. However, this does not mean that the partition function does not depend on $`\mu `$. Let us consider first the case of finite temperature. In this case, the boundary conditions in the Euclidean time direction for the quarks change after (136): they are no longer antiperiodic. The dependence on $`\mu `$ in the partition function comes entirely from the boundary conditions on the fermion fields
$$\psi |_{\tau =1/T}=e^{\mu /T}\psi |_{\tau =0},\text{ and }\psi ^{}|_{\tau =1/T}=e^{\mu /T}\psi ^{}|_{\tau =0}.$$
(137)
The $`\mu `$-independent Lagrangian with such boundary conditions is completely equivalent to the Lagrangian (49) and usual antiperiodic boundary conditions. This fact can be conveniently used in lattice simulations. The corresponding effective theory can also be defined with all $`\mu `$-dependence in the boundary conditions. The boundary conditions for a given effective (composite) bosonic field $`\varphi `$ should then read $`\varphi (1/T)=\mathrm{exp}(b\mu /T)\varphi (0)`$, where $`b`$ is the baryon charge of the field $`\varphi `$ . In the case of the matrix valued field $`\mathrm{\Sigma }`$ we have
$$\mathrm{\Sigma }|_{\tau =1/T}=\mathrm{exp}\left(\frac{\mu }{T}B\right)\mathrm{\Sigma }|_{\tau =0}\mathrm{exp}\left(\frac{\mu }{T}B^T\right),$$
(138)
where $`B`$ is the baryon charge matrix (43).
In the limit of zero temperature, our intuition suggests that the dependence of the partition function on boundary conditions should weaken and disappear. This is, however, not completely true for the boundary conditions such as (137) due to their singular nature in the limit $`1/T\mathrm{}`$. There is indeed an interval of $`\mu `$: $`\mu _0<\mu <\mu _0`$, where the partition function does not depend on $`\mu `$. However, outside this interval, the singular nature of the boundary conditions starts playing a role: a phase transition occurs and a $`\mu `$-dependence appears.
Why is the approach of this paper not directly applicable to real three-color QCD? The main reason is that at $`\mu =0`$ the effective theory described by the chiral Lagrangian does not contain excitations with non-zero baryon number (pions do not carry baryon charge). It is easy to see that, as a consequence, applying either the method of local symmetry of Section 6.2 or the gauge transformation described in this section, one finds no dependence on $`\mu `$ in the effective chiral Lagrangian of QCD. Related to that is the fact that the value of $`\mu _0m_{\mathrm{nucleon}}/3`$ is large, and the approach of this paper, based on an expansion in small $`\mu `$, will not reach it.
As we emphasized, the theories we considered are related by the fact that the fermion representations are pseudo-real. This is a consequence of the antiunitary symmetries of the Dirac operator (1), (2). We identified two such symmetry classes, distinguished by the Dyson index $`\beta =1`$ and $`\beta =4`$. The effective theories for these two classes of theories are very similar and dual, or complementary, to each other, from the point of view of the residual global symmetries. In the diquark (finite density) phase, for the $`\beta =1`$ theories the residual flavor symmetry is given by $`Sp(N_f)`$, while for the $`\beta =4`$ case it is $`SO(N_f)`$. The excitations form multiplets which correspond to symmetric or antisymmetric second rank tensor representations of these groups (see Figs 1,2 and Tables 1,2). One can see that the $`\beta =1`$ and $`\beta =4`$ cases “mirror” each other with respect to $`Sp(N_f)SO(N_f)`$ and symmetric $``$ antisymmetric. On the other hand, the Dirac operator for three-color QCD with fundamental quarks does not have any antiunitary symmetries, and the fermion representations are complex. This case falls into the third remaining Dyson class with the index $`\beta =2`$.
It would be interesting to apply the approach of this paper to lattice theories with pseudoreal fermions, for example to two-color QCD with fundamental quarks. The symmetries of such theories are different from their continuum counterparts and were analyzed in . In particular, the transition to continuum limit may turn our to be nontrivial. This can be related to the antiunitary symmetry of the lattice Dirac operator for staggered fermions in fundamental representation
$$𝒟_{xy}=\frac{1}{2}\underset{\mu }{}\eta _{x,\mu }(U_{x,\mu }\delta _{x+\widehat{\mu },y}U_{x\mu ,\mu }^{}\delta _{x\widehat{\mu },y}),$$
(139)
where $`\eta _{x,\mu }=(1)^{x_1+\mathrm{}+x_{\mu 1}}`$, and $`U`$ are $`SU(2)`$ color matrices. The antiunitary symmetry of this lattice Dirac operator is given by $`\tau _2𝒟=𝒟^{}\tau _2`$, or $`[𝒟,\tau _2K]=0`$. Since $`(\tau _2K)^2=1`$, we conclude that such Dirac operator belongs to the class $`\beta =4`$. However, in the continuum limit, the Dirac operator must be in the class $`\beta =1`$. Such an observation was also made in from the point of view of global symmetries and their breaking, where it was pointed out that it is not yet known how the apparent pattern of $`SU(2N_f)O(2N_f)`$ breaking becomes $`SU(2N_f)Sp(2N_f)`$ in the continuum. Even though the symmetry may modify many of the details of our analysis, when it is applied to lattice theories, some features should be robust. Such would include the phase transition at $`\mu =m_\pi /2`$, the relation similar to $`\overline{\psi }\psi ^2+\psi \psi ^2=\mathrm{const}`$ between the chiral and the diquark condensates, the linear dependence of $`n_B`$ on $`\mu `$ near the transition, the existence of several branches in the spectrum, similar to Figures 1,2,3, and many other qualitative features.
Acknowledgements
S. Hands, H. Leutwyler, M.-P. Lombardo, S. Morrison, E. Shuryak and D.K. Sinclair are acknowledged for useful discussions. J.B.K. is supported in part by the National Science Foundation, NSF-PHY96-05199. D.T. and J.J.M.V. are partially supported by the US DOE grant DE-FG-88ER40388. D.T. is supported in part by “Holderbank”-Stiftung and by Janggen-Pöhn-Stiftung. A.Z. is supported, in part, by the National Science and Engineering Research Council of Canada (NSERC).
## Appendix
In this appendix we analyze the representations of the remaining symmetry groups shown in Table 1 for $`\beta =1`$. They are all different subgroups of $`Sp(2N_f)`$.
Under the symmetry group $`Sp(2N_f)`$ the Goldstone fields $`\mathrm{\Sigma }`$ transform as
$`\mathrm{\Sigma }V\mathrm{\Sigma }V^T.`$ (140)
With $`\mathrm{\Sigma }`$ parameterized as $`U\mathrm{\Sigma }_cU^T`$ (with $`USU(2N_f)`$) and $`V\mathrm{\Sigma }_cV^T=\mathrm{\Sigma }_c`$ (the $`2N_f\times 2N_f`$ antisymmetric unit matrix is denoted by $`\mathrm{\Sigma }_c`$) it follows that the generators transform according to
$`X\mathrm{\Sigma }_cVXV^1\mathrm{\Sigma }_c=VX\mathrm{\Sigma }_cV^T\mathrm{or}(X\mathrm{\Sigma }_c)_{ij}V_{ik}V_{jl}(X\mathrm{\Sigma }_c)_{kl}.`$ (141)
From the transposition relation (12) it follows that $`(X\mathrm{\Sigma }_c)^T=X\mathrm{\Sigma }_c`$. If the symmetry group is $`Sp(2N_f)`$ the generators transform according to an antisymmetric rank two representation of $`Sp(2N_f)`$. The degeneracy is thus given by $`2N_f^2N_f1`$. If the symmetry group is $`SU(N_f)\times U(1)`$ the symmetry transformation is given by
$`V=\left(\begin{array}{cc}U_1& 0\\ 0& U_1^{}\end{array}\right),`$ (144)
the $`P`$-type generators transform according to $`P^TU_1P^TU_1^1`$ and the $`Q`$-type generators as $`QU_1QU_1^T`$. Since the $`Q`$ are antisymmetric they transform according to an antisymmetric rank two representation of $`SU(N_f)`$. The fields $`Q^{}`$ transform according to the conjugate representation. The dimension of both representations is equal to $`N_f(N_f1)/2`$. Since the $`P`$-type generators are traceless, the degeneracy is given by $`N_f^21`$.
In the diquark condensation phase the symmetry group $`Sp(2N_f)`$ is with respect to the rotated antisymmetric unit matrix, $`\mathrm{\Sigma }_\alpha `$, with symplectic transformations defined by
$`V\mathrm{\Sigma }_\alpha V^T=\mathrm{\Sigma }_\alpha .`$ (145)
According to the argument at the beginning of this appendix, the transformation properties of the generators are given by
$`(X_\alpha \mathrm{\Sigma }_\alpha )_{ij}V_{ik}V_{jl}(X_\alpha \mathrm{\Sigma }_\alpha )_{kl},`$ (146)
where the $`X_\alpha `$ are the rotated generators defined by
$`X_\alpha =V_\alpha XV_\alpha ^1,`$ (147)
and $`V_\alpha =\mathrm{exp}i\alpha X_2/2`$ defined in (75).
With $`Sp(N_f)`$ as symmetry group the symmetry transformations are given by (144) with $`U_1^{}=IU_1I`$ (the $`N_f\times N_f`$ antisymmetric unit matrix is denoted by $`I`$). Since in this case $`\{V,X_2\}=0`$ and using that $`\mathrm{\Sigma }=V_\alpha U\mathrm{\Sigma }_cU^TV_\alpha ^T`$ (see eq. (78)) we find that the representations can be discussed in terms of the fields $`U\mathrm{\Sigma }_cU^T`$ with the familiar $`PQ`$-block structure of the generators. The $`P`$-type generators thus transform as $`P^TIU_1P^TIU_1^T`$. The symmetric and and antisymmetric components of $`PI`$ corresponding to $`P_SI`$ and $`P_AI`$, respectively, transform independently. The dimensions of the representations are given by $`N_f(N_f+1)/2`$ and $`N_f(N_f1)1`$, respectively. The generators $`Q`$ and $`IQ^{}I`$ transform in the same way, and thus the linear combinations $`\stackrel{~}{Q}`$ and $`\stackrel{~}{Q}^{}`$ transform in the same way. Since $`Q`$ is antisymmetric they transform according to an irreducible rank two representation with dimension equal to $`N_f(N_f1)/2`$.
Finally, for symmetry group $`Sp(N_f)\times Sp(N_f)`$ the symmetry transformation is given by
$`V=\left(\begin{array}{cc}U_1& 0\\ 0& U_2\end{array}\right),`$ (150)
with both $`U_1`$ and $`U_2`$ symplectic transformations. In this case it is imperative to consider the rotated generators. However, we only need the rotated generators for $`\mu \mathrm{}`$. In terms of the block structure of the unrotated generators, we obtain after a straightforward calculation for $`\alpha =\pi /2`$,
$`X_\alpha \mathrm{\Sigma }_\alpha =\left(\begin{array}{cc}IP_AI+iQ_II& iIP_S+Q_R\\ iP_SI+IQ_RI& P_AiIQ_I\end{array}\right).`$ (153)
Therefore the combination $`iIP_S+Q_R`$ transforms as
$`iIP_S+Q_RU_1(iIP_S+Q_R)U_2^T.`$ (154)
Indeed, this transformation corresponds to the Young tableaux given in upper right corner of Fig. 1. Notice that for $`\mu \mathrm{}`$ we have that $`\stackrel{~}{Q}^{}Q_R`$. The 11- and 22-blocks of the matrix (153) correspond to two different linear combinations of $`P_A`$ and $`Q_I`$. One combination transforms as a rank two tensor with respect to $`U_1`$, and the other combination as a rank two tensor with respect to $`U_2`$. Since both combinations are antisymmetric with respect to transposition, they transform according to the Young tableaux given in the lower right corner of Fig. 1. The additional singlet terms arise because the irreducible representations are traceless.
In case of $`Sp(N_f)`$ symmetry, $`U_2=U_1^{}=IU_1I`$, the transformation properties of $`P_S`$, $`P_A`$, $`Q_R`$ and $`Q_I`$ can be obtained from (153) by combining the transformation properties of the diagonal blocks and of the off-diagonal blocks. The results are in agreement with the discussion in the paragraph following eq. (147).
The same analysis can be performed for $`\beta =4`$. In fact, because the matrix $`I`$ is absent, this case is somewhat simpler, and we leave it as an exercise to the reader. |
warning/0001/math0001149.html | ar5iv | text | # Une formule du type Baker-Campbell-Hausdorff pour les groupoïdes de Lie
## 1. INTRODUCTION
Cet article est consacré a l’ étude locale d’un groupoïde de Lie $`G`$ dans des cartes convenablement choisies. En particulier on obtient dans le théorème 2.4 le developpement de la multiplication ce qui constitue un analogue de la formule de Baker-Campbell-Hausdorff du cas des groupes de Lie (cf. \[K\] par exemple). Ce résultat a été présenté brièvement dans \[Ra1\], \[LR\] ou il est utilisé pour demontrer que le groupoïde tangent $`\stackrel{~}{G}`$ associé à un groupoï\]de de Lie $`G`$ est lui aussi un groupoïde de Lie. Il intervient aussi de manière essentielle dans le calcul du commutateur dans l’algèbre de convolution du groupoïde tangent $`\stackrel{~}{G}`$, ce qui permet de quantifier la structure de Poisson canonique du dual de l’algébroïde de Lie associé à $`G`$.
L’article est structuré comme il suit. Après avoir fixés la terminologie et les notations, on rappelle pour le benefice du lecteur les différentes constructions de l’algébroïde de Lie $`𝒢`$ associé à un groupoïde de Lie $`G`$. Dans la suite on explicite la structure locale de $`G`$ dans une carte et on montre comment on associe à une carte de $`G`$ choisie convenablement, une carte de son algébroïde de Lie $`𝒢`$. On écrit la multiplication et l’inversion de $`G`$ dans ces cartes et on les developpe en sèries de Taylor pour obtenir l’analogue de la formule de Baker-Campbell-Hausdorff. Enfin, comme application, on calcule les fonctions de structure de $`𝒢`$.
Rappelons brièvement les principaux faits sur les groupoïdes de Lie et les algébroïdes de Lie associés. Pour une présentation detaillée de la théorie des groupoïdes de Lie on pourra se reporter aux ouvrages de A. Weinstein, P. Dazord et A. Coste \[CDW\], où K. Mackenzie \[M\]. Les notations et les définitions de la théorie des groupoïdes seront celles données dans \[Re\] par J. Renault. Par définition un groupoïde est un ensemble $`G`$ muni d’un produit $`G\times GG^{(2)}(g_1,g_2)g_1g_2G`$ défini sur un sous-ensemble $`G^{(2)}`$ de $`G\times G`$, et une application inverse $`Ggg^1G`$ vérifiant:
1. $`(g^1)^1=g`$
2. Si $`(g_1,g_2),(g_2,g_3)G^{(2)}`$ alors $`(g_1g_2,g_3),(g_1,g_2g_3)G^{(2)}`$ et $`(g_1g_2)g_3=g_1(g_2g_3)`$
3. $`(g^1,g)G^{(2)}`$. Si $`(g_1,g_2)G^{(2)}`$ alors $`g_1^1(g_1g_2)=g_2`$
4. $`(g,g^1)G^{(2)}`$. Si $`(g_2,g_1)G^{(2)}`$ alors $`(g_2g_1)g_1^1=g_2`$
$`G^{(2)}`$ s’appelle l’ensemble des paires composables et pour $`gG`$ on appelle $`s(g)=g^1g`$ le domaine de $`g`$ et $`r(g)=gg^1`$ l’image de $`g`$. La composition $`g_1g_2`$ est bien définie si et seulement si $`r(g_2)=s(g_1)`$. L’ensemble $`s(G)=r(G)`$, noté $`G^{(0)}`$, sera identifié à une partie de $`G`$ et appelé espace des unités. Pour $`xG^{(0)}`$, on notera $`G^x=r^1(x)`$, $`G_x=s^1(x)`$.
Un groupoïde de Lie est un groupoïde $`G`$ qui a une structure de variété différentiable compatible avec la structure algébrique :
1. $`G^{(0)}`$ est une sous-variété de $`G`$
2. $`r,s:GG^{(0)}`$ sont des submersions
3. la multiplication : $`G^{(2)}G`$ est différentiable
Comme consequences de la définition il faut noter que l’application $`i:GG`$, $`i(\gamma )=\gamma ^1`$ est un difféomorphisme (voir \[M\], p.85), et aussi le fait qu’en notant $`m=\text{dim}G`$ et $`n=\text{dim}G^{(0)}`$, pour tout $`xG^{(0)}`$, $`G^x`$ et $`G_x`$ sont des sous-variétés de $`G`$ de dimension $`mn`$.
On rappelle maintenant les differentes constructions de l’algébroïde de Lie associé à un groupoïde de Lie $`G`$ de base $`G^{(0)}`$. Les algébroïdes de Lie ont été introduits par J. Pradines \[P1\], et généralisent la notion d’algèbre de Lie dans le cadre de la théorie des groupoïdes de Lie.
Pour fixer les notations, pour toute application différentiable $`f`$ entre les variétés $`M`$ et $`N`$, $`Tf`$ designe l’application tangente et $`T_xf`$ l’application tangente en $`x`$ entre les espaces tangents $`T_xM`$ et $`T_{f(x)}N`$. Aussi pour $`E`$ fibré vectoriel de classe $`C^{\mathrm{}}`$ sur la variété $`M`$ on notera par $`C^{\mathrm{}}(M,E)`$ l’ensemble des sections de classe $`C^{\mathrm{}}`$ de $`E`$ sur $`M`$.
Par définition un algébroïde de Lie sur une variété $`M`$ est un triplet constitué d’un fibré vectoriel $`E`$ de base $`M`$ et classe $`C^{\mathrm{}}`$, une structure de -algèbre de Lie sur l’espace des sections $`C^{\mathrm{}}(M,E)`$, dont on note $`[,]`$ le crochet et un morphisme $`\rho :ETM`$ de fibrés vectoriels $`C^{\mathrm{}}`$, appelé ancre tels que:
(i) L’application induite entre les espaces des sections $`\overline{\rho }:C^{\mathrm{}}(M,E)C^{\mathrm{}}(M,TM)`$, $`\overline{\rho }(\xi )(x)=\rho (\xi (x))`$, $`\xi C^{\mathrm{}}(M,E)`$$`xM`$, est un morphisme d’algèbres de Lie : $`[\overline{\rho }(\xi ),\overline{\rho }(\eta )]=\overline{\rho }\left([\xi ,\eta ]\right)`$
(ii) Pour toute fonction $`fC^{\mathrm{}}(M)`$ et pour tout couple $`(\xi ,\eta )`$ de sections $`C^{\mathrm{}}`$ de $`E`$ ,
$`[\xi ,f\eta ]=f[\xi ,\eta ]+\overline{\rho }(\xi )(f)\eta `$
On fixe $`\{e_1,e_2,\mathrm{},e_p\}`$ un repère local sur $`UM`$ pour $`E`$ et $`(q_1,\mathrm{},q_n,\lambda _1,\mathrm{},\lambda _p)`$ coordonnées locales de $`E`$ avec les $`q_i`$ coordonnées locales pour la base $`M`$ et les $`\lambda _j`$ coordonnées dans les fibres associées au repère $`\{e_1,e_2,\mathrm{},e_p\}`$. Alors, localement, le fait que $`E`$ est un algébroïde de Lie implique l’existence des fonctions de structure $`c_{ijk},a_{ij}C^{\mathrm{}}(U)`$ telles que $`[e_i,e_j]={\displaystyle \underset{k}{}}c_{ijk}e_k`$ et $`\overline{\rho }(e_i)={\displaystyle \underset{j}{}}a_{ij}{\displaystyle \frac{}{q_j}}`$.
Remarquons tout d’abord que pour tout $`gG`$, $`R_g:G_{r(g)}G_{s(g)}`$, $`R_g(h)=hg`$ et $`L_g:G^{s(g)}G^{r(g)}`$, $`L_g(h)=gh`$ sont des difféomorphismes et que pour tout $`gG`$ on a $`T_gG^{r(g)}=KerT_gr`$ et $`T_gG_{s(g)}=KerT_gs`$. Cela permet de définir les champs invariants à gauche sur $`G`$ par $`L(G)=\{\xi C^{\mathrm{}}(G,TG)\xi KerTr,TL_g\xi =\xi L_g\}`$ et les champs invariants à droite par $`R(G)=\{\xi C^{\mathrm{}}(G,TG)\xi KerTs,TR_g\xi =\xi R_g\}`$. Il est facile a voir que $`L(G)`$ et $`R(G)`$ sont des algèbres de Lie.
Le fait que $`G^{(0)}`$ est une sous-variété de $`G`$ permet de considerer $`T_xG^{(0)}`$ sous-espace de $`T_xG`$, pour tout $`xG^{(0)}`$. On note $``$, $``$, respectivement $`𝒩`$, les fibrés vectoriels sur $`G^{(0)}`$ dont les fibres au dessus de $`uG^{(0)}`$ sont $`_u=KerT_ur`$, $`_u=KerT_us`$, respectivement $`𝒩_u=T_uG/T_uG^{(0)}`$.
###### LEMME 1.1.
On a les isomorphismes des espaces vectoriels $`L(G)C^{\mathrm{}}(G^{(0)},)`$ et $`R(G)C^{\mathrm{}}(G^{(0)},)`$.
Preuve. L’application $`\mathrm{\Phi }:L(G)C^{\mathrm{}}(G^{(0)},),\mathrm{\Phi }(\xi )=\xi _{_{G^{(0)}}}`$ est bien définie. Montrons que $`\mathrm{\Phi }`$ est injective. Soit $`\xi _{_{G^{(0)}}}=\eta _{_{G^{(0)}}}`$. Alors pour $`gG`$ on a $`\xi (g)=\xi (gu)=T_uL_g\xi (u)=T_uL_g\eta (u)=\eta (gu)=\eta (g)`$, où $`u=s(g)`$.
Montrons que $`\mathrm{\Phi }`$ est surjective. Soit $`\eta C^{\mathrm{}}(G^{(0)},)`$. On définit $`\xi :GTG,`$ $`\xi (g)=T_uL_g\eta (u)`$, pour $`u=s(g)`$. On voit que $`L_u=Id_{G^u}`$, d’où $`\xi (u)=\eta (u)`$ pour $`uG^{(0)}`$, donc $`\eta =\xi _{G^{(0)}}`$. Il reste à montrer que $`\xi L(G)`$. On a $`T_g^{}L_g\xi (g^{})=T_g^{}L_g\left[T_uL_g^{}\xi (u)\right]=T_uL_{gg^{}}\xi (u)=\xi (gg^{})`$, où $`u=s(g^{})`$. On a utilisé le fait que $`T_g^{}L_gT_uL_g^{}=T_u(L_gL_g^{})=T_uL_{gg^{}}`$. Comme $`\mathrm{\Phi }`$ est évidemment linéaire on a démontré le premier isomorphisme. La démonstration du second isomorphisme est analogue.
Les fibrés $``$, avec le crochet donné par $`[\xi ,\eta ]=\mathrm{\Phi }([\mathrm{\Phi }^1(\xi ),\mathrm{\Phi }^1(\eta )])`$ et l’ancre $`\rho :TG^{(0)}`$, $`\rho _u=T_us`$, respectivement $``$, avec le crochet défini de manière analogue à $``$ et l’ancre $`\mu :TG^{(0)}`$, $`\mu _u=T_ur`$, sont deux algébroïdes de Lie antiisomorphes par l’application tangente $`Ti`$ de l’inversion $`i`$ de $`G`$.
L’application $`I_xT_xr:T_xGKerT_xr`$, $`XXT_xrX`$, est bien définie et surjective. Son noyau est $`T_xG^{(0)}`$ et par factorisation on obtient l’isomorphisme d’espaces vectoriels $`𝒩_xKerT_xr`$. De la même manière, en considèrant $`I_xT_xs:T_xGKerT_xs`$, $`XXT_xsX`$ on démontre l’isomorphisme $`𝒩_xKerT_xs`$.
Ces deux isomorphismes définissent sur $`𝒩`$ deux structures d’algébroïde de Lie antiisomorphes. Le crochet de Lie sur $`C^{\mathrm{}}(G^{(0)},𝒩)`$ est défini en utilisant l’isomorphisme $`L(G)\xi \left[\xi _{G^{(0)}}\right]C^{\mathrm{}}(G^{(0)},𝒩)`$, où $`\left[X\right]`$ est l’image de $`XT_xG`$ dans $`T_xG/T_xG^{(0)}`$. L’ancre sur $`𝒩`$ est $`\rho :𝒩TG^{(0)}`$, $`\rho _x=T_xrT_xs`$.
Dans la suite on appellera algébroïde de Lie du groupoïde de Lie $`G`$ le fibré $``$ avec la structure d’algébroïde définie précédemment et pour mettre en évidence qu’il est l’algébroïde associé au groupoïde $`G`$ il sera noté $`𝒢`$.
## 2. La structure locale d’un groupoïde de Lie
### 2.1. Les cartes
Soit $`G`$ un groupoïde de Lie et $`𝒢`$ son algébroïde de Lie. On va expliciter dans cette section la structure de $`G`$ dans une carte convenablement choisie au voisinage d’un point $`x_0G^{(0)}G`$. Des cartes de ce genre ont été utilisées aussi dans \[NWX\].
Comme $`r`$ est une submersion au point $`x_0`$ appartenant à la sous-variété $`G^{(0)}`$ de $`G`$, il existe $`U`$ voisinage ouvert de 0 dans $`\text{}^n`$, $`V`$ voisinage ouvert de 0 dans $`\text{}^m`$ et les cartes $`\psi :U\times VG`$, $`\phi :UG^{(0)}`$ vérifiant:
(1)
$$\begin{array}{c}1.\psi (0,0)=x_0\hfill \\ 2.r(\psi (u,v))=\phi (u)\hfill \\ 3.\psi (U\times \{0\})=\psi (U\times V)G^{(0)}\hfill \end{array}$$
La deuxième condition revient au diagramme commutatif :
Des deux dernières conditions on déduit $`\phi (u)=\psi (u,0)`$, et en conséquence on pourra exprimer la structure de $`G`$ en utilisant seulement la carte $`\psi `$. Toutefois pour la simplicité des notations on gardera $`\phi =\psi (,0)`$.
A la carte $`\psi `$ de $`G`$ s’associe canoniquement une carte de l’algébroïde de Lie $`𝒢`$. Plus précisément on a :
###### LEMME 2.1.
L’application $`\theta :U\times \text{}^m𝒢`$, $`\theta (u,v)=(\phi (u),{\displaystyle \frac{\psi }{v}}(u,0)v)`$ est une carte de $`𝒢`$ au voisinage de la fibre $`𝒢_{x_0}`$ et la famille $`\{e_1,e_2,..,e_m\}`$ définie par $`e_i(\phi (u))=\theta (u,f_i)`$, $`i=\overline{1,m}`$, où $`\{f_1,f_2,\mathrm{},f_m\}`$ est la base canonique de $`\text{}^m`$, est un repère mobile de $`𝒢`$ sur $`\phi (U)`$.
Preuve. Pour tout $`uU`$ on a
$`G^{\phi (u)}\psi (U\times V)=\{\gamma \psi (U\times V)|r(\gamma )=\phi (u)\}=\psi (\{u\}\times V)`$
L’application $`\psi (u,):VG^{\phi (u)}`$ est alors une carte de la sous-variété $`G^{\phi (u)}`$, qui associe $`0V`$ à $`\phi (u)`$. On peut identifier $`\{u\}\times \text{}^m`$ avec $`T_{\phi (u)}G^{\phi (u)}`$ par l’isomorphisme $`{\displaystyle \frac{\psi }{v}}(u,0):T_{(u,0)}(\{u\}\times V)=\{u\}\times \text{}^mT_{\phi (u)}G^{\phi (u)}`$. L’image de $`\theta `$ est le voisinage $`\theta (U\times \text{}^m)=𝒢_{\phi (U)}`$ de la fibre $`𝒢_{x_0}`$.
Remarque. Dans un groupe de Lie il existe un voisinage de l’unité difféomorphe avec un voisinage de l’élément nul de l’algèbre de Lie associé. Le lemme précédent permet de donner la généralisation suivante pour les groupoïdes de Lie :
Pour tout $`x_0G^{(0)}`$ il existe un voisinage de $`x_0`$ dans $`G`$ qui est difféomorphe avec un voisinage de $`(x_0,0)`$ dans $`𝒢`$.
En effet avec les notations précedentes, $`\alpha =\psi \theta ^1`$ est un difféomorphisme entre le voisinage $`\theta (U\times V)`$ de $`(x_0,0)𝒢`$ et le voisinage $`\psi (U\times V)`$ de $`x_0`$ dans $`G`$. On remarque de plus que pour tout $`x\phi (U)`$, $`\alpha (𝒢_x)G^x`$.
Ce résultat n’est qu’un cas particulier de la proposition suivante qui est basée sur l’existence d’une application exponentielle pour tout groupoïde de Lie. Cette application exponentielle introduite par Pradines dans \[P2\] généralise à la fois l’exponentielle d’un groupe de Lie et l’exponentielle d’une variété munie d’une connexion.
###### PROPOSITION 2.2.
Soit $`G`$ un groupoïde de Lie et $`𝒢`$ son algébroïde de Lie. Il existe alors un voisinage $`V`$ de $`G^{(0)}`$ vu comme la section nulle $`\left\{(x,0)\right|xG^{(0)}\}`$ dans $`𝒢`$, un voisinage $`W`$ de $`G^{(0)}`$ dans $`G`$ et un difféomorphisme $`\alpha :VW`$ tel que $`\alpha (𝒢_xV)=G^xW`$ et $`\alpha _x^{}(0)`$ est l’identité de $`𝒢_x`$, où $`\alpha _x`$ est la restriction de $`\alpha `$ sur $`𝒢_xV`$.
L’idée de la démonstration est la suivante. Soit $``$ une connexion sur l’algébroïde de Lie $`𝒢`$. On associe à $``$ une connexion invariante à gauche sur $`G`$, dont la restriction à $`G^x`$ est une connexion linéaire $`_x`$. On peut alors définir fibre par fibre une application exponentielle, et prendre comme $`\alpha `$ cette exponentielle. Pour les détails voir \[L\] ou \[NWX\].
### 2.2. La multiplication et l’inversion
Pour exprimer le produit et l’inversion de $`G`$ dans la carte $`\psi `$ on a besoin de la forme de l’application source $`s`$ dans cette carte, forme qui est explicitée dans le lemme suivant.
###### LEMME 2.3.
Il existe une submersion $`\sigma :U\times VU`$ telle que $`s(\psi (u,v))=\phi (\sigma (u,v))`$. De plus $`\sigma (u,0)=u`$.
Preuve. En réduisant eventuellement $`V`$, on peut supposer que $`s(\psi (u,v))\phi (U)`$, pour $`(u,v)U\times V`$. Il existe alors un élément $`\sigma (u,v)U`$ tel que $`s(\psi (u,v))=\phi (\sigma (u,v))`$. Evidemment $`\sigma =\phi ^1s\psi `$ est une submersion, comme expression dans les cartes de la submersion $`s`$. Enfin $`\phi (\sigma (u,0))=s(\psi (u,0))=\psi (u,0)=\phi (u)`$, donc $`\sigma (u,0)=u`$.
On peut maintenant donner les développements dans la carte $`\psi `$ de la multiplication et de l’inversion de $`G`$. Le résultat suivant représente l’analogue de la formule de Baker-Campbell-Hausdorff pour les groupoïdes de Lie.
###### PROPOSITION 2.4.
(i) Pour $`u,u_1U`$ et $`v,wV`$ on a $`(\psi (u,v),\psi (u_1,w))G^{(2)}`$ si et seulement si $`u_1=\sigma (u,v)`$. Dans ce cas le produit est donné par $`\psi (u,v)\psi (\sigma (u,v),w)=\psi (u,p(u,v,w))`$$`p:U\times V\times VV`$ est une application différentiable qui a un développement de la forme $`p(u,v,w)=v+w+B(u,v,w)+O_3(u,v,w)`$ avec $`B`$ bilinéaire en $`(v,w)`$ et $`O_3(u,v,w)`$ de l’ordre de $`(v,w)^3`$.
(ii) Soit $`(u,v)U\times V`$ tel que $`\psi (u,v)^1\psi (U\times V)`$. Alors $`\psi (u,v)^1=\psi (\sigma (u,v),w)`$, où $`w`$ vérifie $`p(u,v,w)=0`$. De plus on a le développement $`w=v+B(u,v,v)+O_3(u,v)`$, avec $`O_3(u,v)`$ de degré d’homogénéité superieur à 3 en $`v`$.
Preuve. (i) Soit $`g=\psi (u,v)`$ et $`h=\psi (u_1,w)`$. On a $`s(g)=\phi (\sigma (u,v))`$ et $`r(h)=\phi (u_1)`$ ce qui montre que $`(g,h)G^{(2)}`$ si et seulement si $`u_1=\sigma (u,v)`$. De plus $`r(gh)=r(g)=\phi (u)`$ assure l’existence d’un unique $`p(u,v,w)V`$ tel que
$`\psi (u,v)\psi (\sigma (u,v),w)=\psi (u,p(u,v,w))`$
On définit ainsi l’application $`p:U\times V\times VV`$, qui vérifie en particulier $`p(u,0,w)=w`$ et $`p(u,v,0)=v`$. En effet
$`\psi (u,p(u,0,w))=\psi (u,0)\psi (\sigma (u,0),w)=\phi (u)\psi (u,w)=\psi (u,w)`$
$`\psi (u,p(u,v,0))=\psi (u,v)\psi (\sigma (u,v),0))=\psi (u,v)s(\psi (u,v))=\psi (u,v)`$.
Il s’ensuit que $`{\displaystyle \frac{p}{v}}(u,v,0)=I`$, $`{\displaystyle \frac{p}{w}}(u,0,w)=I`$, $`{\displaystyle \frac{^2p}{v^2}}(u,0,0)=0`$, $`{\displaystyle \frac{^2p}{w^2}}(u,0,0)=0`$, et par un développement de Taylor $`p(u,v,w)=p(u,0,0)+{\displaystyle \frac{p}{v}}(u,0,0)v+{\displaystyle \frac{p}{w}}(u,0,0)w+B(u,v,w)+O_3(u,v,w)=v+w+B(u,v,w)+O_3(u,v,w)`$, où $`B(u,v,w)`$ est pour chaque $`u`$ bilinéaire en $`(v,w)`$ et $`O_3(u,v,w)`$ est homogène d’un degré superieur à 3 en $`v`$ et $`w`$.
(ii) Soit $`g=\psi (u,v)`$. On cherche $`u_1`$ et $`w`$ tels que $`g^1=\psi (u_1,w)`$. D’une part $`r(g^1)=\phi (u_1)`$ et $`s(g)=\phi (\sigma (u,v))`$ impliquent $`u_1=\sigma (u,v)`$. D’autre part comme $`gg^1=r(g)`$, on déduit $`\psi (u,v)\psi (\sigma (u,v),w)=\psi (u,0)`$, donc $`p(u,v,w)=0`$. Le théorème des fonctions implicites assure l’existence d’un $`f`$ différentiable tel que $`w=f(u,v)`$. On développe
$`f(u,v)=f(u,0)+{\displaystyle \frac{f}{v}}(u,0)v+\mathrm{}=f_1(u,v)+f_2(u,v)+\mathrm{}`$
$`f_k(u,v)`$ est de degré d’homogénéité $`k`$ en $`v`$. On va déterminer $`f_1`$ et $`f_2`$. Pour cela on utilise $`p(u,v,w)=0`$. On a donc
$`0=p(u,v,w)=v+w+B(u,v,w)+\mathrm{}`$
$`=v+f_1(u,v)+f_2(u,v)+\mathrm{}+B(u,v,f_1(u,v)+f_2(u,v)+\mathrm{})+\mathrm{}`$
Le terme de degré d’homogénéité 1 dans le développement précédant est $`v+f_1(u,v)`$ donc $`f_1(u,v)=v`$. On remplace dans l’équation précédente et on trouve $`0=f_2(u,v)+B(u,v,v)+\mathrm{}`$. En identifiant le terme de degré 2 on obtient $`f_2(u,v)=B(u,v,v)`$.
## 3. Le calcul des fonctions de structure de $`𝒢`$
On se propose de calculer les fonctions de structure de l’algébroïde de Lie $`𝒢`$. On rappelle que $`a_{ij}=\overline{\rho }(e_i)(q_j)`$ et les $`c_{ijk}`$ sont données par $`[e_i,e_j]=c_{ijk}e_k`$, où $`\rho =Ts`$ est l’ancre de $`𝒢`$, $`\{e_1,\mathrm{},e_m\}`$ le repère mobile de $`𝒢`$ défini dans le lemme 2.1 et $`q_j=pr_j\phi ^1`$ sont les fonctions de coordonnées de $`G^{(0)}`$. On notera par $`B_1,..,B_m`$ les coordonnées de l’application $`B:U\times V\times VV`$ dans la base $`\{f_1,f_2,..,f_m\}`$ de $`\text{}^m`$.
###### PROPOSITION 3.1.
Pour tout $`uU`$, les fonctions de structure de l’algébroïde $`𝒢`$ sont données par $`a_{ij}(\phi (u))={\displaystyle \frac{\sigma _j}{v_i}}(u,0)`$ et $`c_{ijk}(\phi (u))=B_k(u,f_i,f_j)B_k(u,f_j,f_i)`$.
Preuve. (i) Le calcul de $`a_{ij}`$
On obtient la forme de $`a_{ij}`$ par le calcul suivant
$`a_{ij}(\phi (u))=\overline{\rho }(e_i)(q_j)(\phi (u))=\rho (e_i(\phi (u)))(q_j)=(T_{\phi (u)}s)(e_i(\phi (u)))(q_j)`$
$`=e_i(\phi (u))(q_js)=({\displaystyle \frac{\psi }{v}}(u,0)f_i)(pr_j\phi ^1s)`$
$`={\displaystyle \frac{}{v_i}}(pr_j\phi ^1s\psi )(u,0)={\displaystyle \frac{\sigma _j}{v_i}}(u,0)`$.
(ii) Le calcul de $`c_{ijk}`$
1. On explicite d’abord la forme d’une section $`\xi C^{\mathrm{}}(G^{(0)},𝒢)`$ dans les cartes de $`G^{(0)}`$ et $`𝒢`$ engendrées par $`\psi `$.
Cette forme est donnée par $`\mathrm{\Xi }(u)=\theta ^1\xi \phi (u)=(u,{\displaystyle \frac{\psi }{v}}(u,0)^1(\xi (\phi (u))))`$. On note $`\xi _0:U\text{}^m`$, $`\xi _0(u)={\displaystyle \frac{\psi }{v}}(u,0)^1(\xi (\phi (u)))`$ et on a $`\mathrm{\Xi }(u)=(u,\xi _0(u)`$.
2. Par le lemme 1.1, on associe à tout $`\xi C^{\mathrm{}}(G^{(0)},𝒢)`$ une section équivariante à gauche $`\mathrm{\Phi }^1(\xi )C^{\mathrm{}}(G,TG)`$, définie par $`\mathrm{\Phi }^1(\xi )(\gamma )=(T_{s(\gamma )}L_\gamma )\xi (s(\gamma ))`$ et le crochet de Lie sur $`C^{\mathrm{}}(G^{(0)},𝒢)`$ est donné par $`[\xi ,\eta ]=\mathrm{\Phi }\left([\mathrm{\Phi }^1(\xi ),\mathrm{\Phi }^1(\eta )]\right)`$.
On note $`\mathrm{\Omega }(\xi )=(\psi ^{})^1\mathrm{\Phi }^1(\xi )\psi `$ la forme de $`\mathrm{\Phi }^1(\xi )`$ dans les cartes.
On montre dans cette étape que
(2)
$$\mathrm{\Omega }(\xi )(u,v)=(0,\xi _0(\sigma (u,v))+B(u,v,\xi _0(\sigma (u,v))))$$
Par définition
$`\mathrm{\Phi }^1(\xi )(\psi (u,v))=\left(T_{s(\psi (u,v))}L_{\psi (u,v)}\right)\xi (s(\psi (u,v)))`$
$`=\left(T_{s(\psi (u,v))}L_{\psi (u,v)}\right)\left({\displaystyle \frac{\psi }{v}}(\sigma (u,v),0)\xi _0(\sigma (u,v))\right)`$
$`=T_0\left(L_{\psi (u,v)}\psi (\sigma (u,v),)\right)\xi _0(\sigma (u,v))`$
Mais $`L_{\psi (u,v)}\psi (\sigma (u,v),w)=\psi (u,v)\psi (\sigma (u,v),w)=\psi (u,p(u,v,w))`$, ce qui par dérivation conduit à $`T_0\left(L_{\psi (u,v)}\psi (\sigma (u,v),)\right)={\displaystyle \frac{\psi }{v}}(u,v){\displaystyle \frac{p}{w}}(u,v,0)`$, donc
(3)
$$\mathrm{\Phi }^1(\xi )(\psi (u,v))=\frac{\psi }{v}(u,v)\frac{p}{w}(u,v,0)\xi _0(\sigma (u,v))$$
Comme $`p(u,v,w)=v+w+B(u,v,w)+O_3(u,v,w)`$ on a
(4)
$$\frac{p}{w}(u,v,0)=I+\left(\begin{array}{cccc}B_1(u,v,f_1)& B_1(u,v,f_2)& \mathrm{}& B_1(u,v,f_m)\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ B_m(u,v,f_1)& B_m(u,v,f_2)& \mathrm{}& B_m(u,v,f_m)\end{array}\right)$$
En remplaçant 4 dans 3 on obtient
(5)
$$\mathrm{\Phi }^1(\xi )(\psi (u,v))=\frac{\psi }{v}(u,v)\xi _0(\sigma (u,v))+\frac{\psi }{v}(u,v)B(u,v,\xi _0(\sigma (u,v)))$$
La formule 2 résulte comme suit :
$`\mathrm{\Omega }(\xi )(u,v)=\psi ^{}(u,v)^1\left(\mathrm{\Phi }^1(\xi )(\psi (u,v))\right)`$
$`=\psi ^{}(u,v)^1{\displaystyle \frac{\psi }{v}}(u,v)\left(\xi _0(\sigma (u,v))+B(u,v,\xi _0(\sigma (u,v)))\right)`$
$`=(0,\xi _0(\sigma (u,v))+B(u,v,\xi _0(\sigma (u,v))))`$
3. La formule 2 montre en particulier que $`\mathrm{\Omega }(e_i)(u,v)=(0,f_i+B(u,v,f_i))`$. Pour les éléments de cette forme le crochet de Lie dans $`T(U\times V)`$ est donné par
$`[(0,a(u,v)),(0,b(u,v))]=(0,{\displaystyle \underset{k}{}}(a_k(u,v){\displaystyle \frac{b}{v_k}}(u,v)b_k(u,v){\displaystyle \frac{a}{v_k}}(u,v)))`$
En utilisant la bilinéarité de $`B`$ on obtient $`[\mathrm{\Omega }(e_i),\mathrm{\Omega }(e_j)](u,v)=`$
$`=(0,{\displaystyle \underset{k}{}}(\delta _{ik}+B_k(u,v,f_i))B(u,f_k,f_j){\displaystyle \underset{k}{}}(\delta _{jk}+B_k(u,v,f_j))B(u,f_k,f_i))`$
$`=(0,B(u,f_i,f_j)+B(u,B(u,v,f_i),f_j)B(u,f_j,f_i)B(u,B(u,v,f_j),f_i))`$
et en particulier $`[\mathrm{\Omega }(e_i),\mathrm{\Omega }(e_j)](u,0)=`$
$`=(0,B(u,f_i,f_j)B(u,f_j,f_i)+B(u,B(u,0,f_i),f_j)B(u,B(u,0,f_j),f_i))`$
$`=(0,B(u,f_i,f_j)B(u,f_j,f_i)).`$
On peut ainsi conclure
$`[e_i,e_j](\phi (u))=[\mathrm{\Phi }^1(e_i),\mathrm{\Phi }^1(e_j)](\psi (u,0))=\psi ^{}(u,0)[\mathrm{\Omega }(e_i),\mathrm{\Omega }(e_j)](u,0)`$
$`={\displaystyle \frac{\psi }{v}}(u,0){\displaystyle \underset{k}{}}\left(B_k(u,f_i,f_j)B_k(u,f_j,f_i)\right)f_k`$
$`={\displaystyle \underset{k}{}}(B_k(u,f_i,f_j)B_k(u,f_j,f_i))e_k(\phi (u))`$.
###### COROLLAIRE 3.2.
Pour toutes sections $`\xi ,\eta C^{\mathrm{}}(G^{(0)},𝒢)`$ on a
$`[\xi ,\eta ]_0(u)=B(u,\xi _0(u),\eta _0(u)B(u,\eta _0(u),\xi _0(u))`$
Preuve. On développe $`\xi =\xi _ie_i`$ et $`\eta =\eta _je_j`$ dans le repère mobile $`\{e_1,\mathrm{},e_m\}`$ et on s’en sert de la proposition précedente pour remplacer $`c_{ijk}`$ dans $`[\xi ,\eta ]=\xi _i\eta _jc_{ijk}e_k`$. Il ne reste qu’à utiliser la bilinéarité de $`B`$.
Exemple. Soit $`G`$ un groupe de Lie d’unité $`e`$. Dans ce cas $`U=\left\{0\right\}`$ et $`\psi :VG`$ est une carte vérifiant $`\psi (0)=e`$. On associe à $`\psi `$ la carte $`\theta =\psi ^{}(0):\text{}^m𝒢`$ de l’algèbre de Lie. Par la proposition 2.4 le produit dans $`G`$ est de la forme $`\psi (v)\psi (w)=\psi (p(v,w))`$, où $`p:V\times VV`$ est une application différentiable qui admet un développement de la forme $`p(v,w)=v+w+B(v,w)+O_3(v,w)`$ avec $`B`$ bilinéaire et l’inversion est donnée par $`\psi (v)^1=\psi (v+B(v,v)+O_3(v))`$. Comme $`\sigma =0`$ on retrouve $`a_{ij}=0`$ et la proposition 3.1 montre que les constantes de structure $`c_{ijk}`$ de l’algèbre de Lie sont données par $`c_{ijk}=B_k(f_i,f_j)B_k(f_j,f_i)`$.
On retrouve ainsi des résultats connus pour les groupes et algèbres de Lie qu’on peut trouver par exemple dans \[K\].
Exemple. Soit $`G=M\times M`$, le groupoïde principal transitif associé à la variété $`M`$, un élément $`(x,x)G^{(0)}`$ et $`\alpha :U_1M`$ une carte de $`M`$ telle que $`\alpha (0)=x`$. On peut prendre la carte $`\psi :U\times VG`$, donnée par $`\psi (u,v)=(\alpha (u),\alpha (u+v))`$, où $`U,V`$ sont des voisinages de $`0\text{}^m`$ telles que $`UU_1`$ et $`U+VU_1`$. La carte de $`TM`$ engendrée par $`\psi `$ est $`\theta :U\times \text{}^m𝒢`$, $`\theta (u,v)=(\alpha (u),\alpha ^{}(u)v)`$. Dans ce cas $`\sigma (u,v)=u+v`$, $`p(u,v,w)=v+w`$, $`B(u,v,w)=0`$. Pour les fonctions de structure de $`TM`$ on retrouve $`a_{ij}=\delta _{ij}`$ et $`c_{ijk}=0`$.
Remerciements Les résultats de cet article ont été obtenu durant mon sejour à l’Université d’Orléans et je tiens a exprimer ma gratitude à Jean Renault pour tout son soutien et ses très judicieuses remarques et à Claire Anantharaman qui au début de ce travail m’a fait comprendre le cas des groupes de Lie.
Académie Roumaine
Institut de Mathématiques
Calea Griviţei 21, P.O. Box 1-764,
Bucarest 70700, Romania
E-mail address: ramazan@pompeiu.imar.ro |
warning/0001/astro-ph0001186.html | ar5iv | text | # On the gamma-ray fluxes expected from Cassiopeia A
## 1 Introduction
The shell type supernova remnant (SNR) Cassiopeia A is a prominent source of nonthermal radiation in the Galaxy. It is the brightest and one of the best studied radio sources (e.g. Bell 1975; Tuffs 1986; Braun et al. 1987; Anderson et al. 1991; Kassim et al. 1995; etc.), with synchrotron radiation fluxes observed also at sub-millimeter wavelengths (Mezger et al. 1986), and probably even further in the near infrared (Tuffs et al. 1997) and at hard X-ray energies (Allen et al. 1997; Favata et al. 1997). The powerful radio emission of Cas A implies a total energy in relativistic electrons of order $`10^{48}\mathrm{erg}`$ or higher (e.g. Chevalier et al. 1978; Anderson et al. 1991). A synchrotron origin of hard X-rays implies that relativistic electrons are accelerated up to energies of a few tens of TeV (Allen et al. 1997; Favata et al. 1997). Thus, one should also expect the production of nonthermal $`\gamma `$-rays in Cas A, due to the bremsstrahlung and inverse Compton (IC) mechanisms, extending possibly beyond TeV energies. However, in the MeV $`\gamma `$-ray domain only the $`{}_{}{}^{44}\mathrm{Ti}`$ line emission at an energy of $`1.157\mathrm{MeV}`$ has been detected by COMPTEL (Iyudin et al. 1994). High energy $`\gamma `$-ray observations of Cas A by EGRET resulted in the flux upper limit $`I(>100\mathrm{MeV})1.2\times 10^7\mathrm{ph}/\mathrm{cm}^2\mathrm{s}`$ (Esposito et al. 1996). At TeV energies flux upper limits have been given by the Whipple (Lessard et al. 1999) and CAT (Goret et al. 1999) collaborations. A tentative detection of a weak signal at a significance level of $`5\sigma `$ has been recently reported by the HEGRA collaboration (Pühlhofer et al. 1999).
Gamma-ray fluxes expected from Cas A have been earlier calculated by Cowsik & Sarkar (1980) who have derived a lower limit to the mean magnetic field in the shell of Cas A, $`B_08\times 10^5\mathrm{G}`$, comparing the bremsstrahlung flux of radio emitting electrons with the upper flux limit $`I(>100\mathrm{MeV})1.1\times 10^6\mathrm{ph}/\mathrm{cm}^2\mathrm{s}`$ of SAS-2 (Fichtel et al. 1975) and COS B detectors. Calculations have been done in a commonly used ‘single-zone’ model approximation, which assumes a spatially homogeneous source containing magnetic field, relativistic electrons and gas. In the same single-zone approximation, broad-band fluxes of $`\gamma `$-rays expected from Cas A have been recently calculated by Ellison et al. (1999) and Goret et al. (1999).
In our previous work (Atoyan et al. 1999, hereafter Paper 1) we have carried out detailed modelling of the synchrotron radiation of Cas A from the radio to the X-ray bands in order to understand the relativistic electron content in the source. The principal feature of our study in Paper 1 was that we have proposed a spatially inhomogeneous model, consisting of two zones, that distinguishes between compact, bright steep-spectrum radio knots and the bright fragmented radio ring on the one hand, and the remainder of the shell - the diffuse ‘plateau’ - on the other hand. In this paper, using the model parameters derived from the interpretation of synchrotron fluxes observed, we calculate the fluxes of $`\gamma `$-rays predicted in the framework of the spatially inhomogeneous study of this SNR, and discuss implications which can be derived from future $`\gamma `$-ray observations of this SNR in different energy bands. In section 2, after a brief overview of the basic features of the model, we study the fluxes to be expected for $`\gamma `$-ray energies up to $`E10100\mathrm{GeV}`$. This emission is contributed mostly by bremsstrahlung of those electrons which are responsible for the synchrotron emission in the radio-to-infrared bands. In section 3 we discuss the fluxes expected in the very high energy region, $`E100\mathrm{GeV}`$, where both bremsstrahlung and inverse Compton radiation of relativistic electrons are important. In addition, a significant contribution to the total fluxes could be due to $`\pi ^0`$-decay $`\gamma `$-rays. These should originate from relativistic protons which presumably are also accelerated in Cas A. A summary and the conclusions are contained in section 4.
## 2 High energy bremsstrahlung of radio electrons
### 2.1 Overview of the model
In Paper 1 we have shown that the basic features of the observed radio emission, in particular,
(a) spectral turnover of the energy fluxes $`J_\nu \nu ^\alpha `$ at frequencies below 20 MHz,
(b) mean spectral index $`\alpha 0.77`$ of the total flux at higher frequencies up to $`\nu 30\mathrm{GHz}`$ (e.g. Baars et al. 1977) which then flattens to $`\alpha 0.65`$ in the submillimeter domain (Mezger et al. 1986),
(c) the spread of spectral indices of the radio knots from $`\alpha 0.6`$ to $`\alpha 0.9`$ (Anderson and Rudnick, 1996),
and other characteristics can be explained in the framework of a model that takes into account the effects of energy dependent propagation of relativistic electrons in a spatially nonuniform remnant. We have considered the simplest inhomogeneous model, consisting of two zones with essentially different spatial densities of relativistic electrons. The model allows a rather clear qualitative and quantitative distinction between compact bright radio structures in Cas A like the fragmented radio ring and the radio knots on the one hand, and the remainder of the radio emitting shell – the diffuse radio ‘plateau’ region enclosed between the radio ring at $`R_{\mathrm{ring}}1.7\mathrm{pc}`$ and the blast wave at $`R_02.5\mathrm{pc}`$ (for a distance $`d=3.4\mathrm{kpc}`$ to Cas A, cf. Reed et al. 1995) – on the other. The model predicts that the spatial density of relativistic electrons in the compact radio structures, which all together constitute what we termed zone 1, is much higher than in the surrounding plateau region, termed zone 2. Thus, zone 1 with a total volume $`V_1`$ much smaller that the volume $`V_0`$ of the shell is actually merged into zone 2 with a volume $`V_2=V_0V_1V_0`$.
In Paper 1 we have derived, by spatial integration of the diffusion equation for relativistic electrons, the set of Leaky-box type equations for the overall energy distribution functions of the electrons $`N_1(E,t)`$ and $`N_2(E,t)`$ in zones 1 and 2, respectively. We do not consider the acceleration process itself, but assume that the accelerated electrons are injected into zone 1 with some rate $`Q(E,t)`$, and study the effects of the energy-dependent diffusive propagation of the particles on their spectra $`N_{1,2}(E,t)`$. Such a phenomenological approach is justified if particle acceleration occurs in a volume significantly smaller than the volume where the bulk of the emission of zone 1 is produced. The equation for zone 1 then reads:
$$\frac{N_1}{t}=\frac{(P_1N_1)}{E}\frac{N_1}{\tau _{\mathrm{esc}}}+\frac{V_1N_2}{V_2\tau _{\mathrm{dif}}}+Q,$$
(1)
where $`P_1`$ represents the energy loss rate of the electrons in zone 1. The second term on the right hand side of Eq.(1) describes the escape of the electrons into the plateau region on timescales
$$\tau _{\mathrm{esc}}(E)=\left[\frac{1}{\tau _{\mathrm{dif}}(E)}+\frac{1}{\tau _\mathrm{c}}\right]^1.$$
(2)
The term $`\tau _\mathrm{c}2R_1/u_1`$, where $`R_1(0.030.1)\mathrm{pc}`$ is a typical radius of the compact radio components and $`u_1`$ is the fluid speed in zone 1, corresponds to the convective escape time which does not depend on the electron energy. The principal energy dependent term in Eq.(2) is the diffusive escape time:
$$\tau _{\mathrm{dif}}(E)=\tau _{}(E/E_{})^\delta +\tau _{\mathrm{min}},$$
(3)
where normalization to a typical energy $`E_{}=1\mathrm{GeV}`$ of radio emitting electrons is used, and $`\tau _{\mathrm{min}}=2R_1/c`$ takes into account that the escape time cannot be shorter than the light travel time.
The third term on the right hand side of Eq.(1) takes into account that the diffusive propagation and escape of particles is effectively possible only in the directions opposite to the gradients of their spatial density, i.e. from zone 1 with a high concentration of radio electrons into zone 2. In a standard power-law approximation for the above source function
$$Q(E)=Q_0E^{\beta _{\mathrm{acc}}}\mathrm{exp}(E/E_\mathrm{c}),$$
(4)
the spectrum of the accelerated electrons can be hard, $`\beta _{\mathrm{acc}}2.2`$. However the energy distribution $`N_1(E)`$ of radio electrons in zone 1 becomes steeper than $`Q(E)`$ on the timescale $`\tau _{\mathrm{esc}}(E)`$ of electron escape from zone 1 into zone 2. It is important that, as shown in Paper 1, the degree of steepening at energy $`E`$ essentially depends on the ratio (or ‘the gradient’) of the energy densities of the electrons $`N_1(E)/V_1`$ to $`N_2(E)/V_2`$ in these two zones. The maximum steepening, corresponding to the spectrum $`N_1(E)Q(E)\times \tau _{\mathrm{esc}}(E)E^{\beta _{\mathrm{max}}}`$ with $`\beta _{\mathrm{max}}=\beta _{\mathrm{acc}}+\delta `$, is reached only if this gradient is very high, and there is no steepening at all if the energy densities become equal. Therefore, although the two-zone model is also rather simplified, as compared with the reality, because it assumes the same electron density for all components in zone 1, it allows a qualitative explanation for the variations of the power-law indices of individual radio knots from $`\alpha _{\mathrm{min}}0.6`$ to $`\alpha _{\mathrm{max}}0.90.95`$ (Anderson & Rudnick 1996; see also Wright et al. 1999) assuming an efficient electron acceleration with the same hard power-law spectrum $`\beta _{\mathrm{acc}}=1+2\alpha _{\mathrm{min}}2.2`$ and a diffusive escape time with a parameter $`\delta 0.60.7`$. These variations do not then necessarily imply that particle acceleration significantly varies across the remnant (Wright et al. 1999), but can be explained rather by differences in the local gradients, with respect to the surrounding plateau, of the concentration of radio electrons for different members of zone 1.
In Paper 1 we have shown that a high contrast of electron densities needed for interpretation of the radio observations can be reached if the zone 1 components correspond to the sites of efficient electron acceleration. In principle, the interpretation of the radio data of Cas A does not suggest an efficient acceleration of the electrons in zone 2. We note however that such an acceleration, with similarly hard power-law index $`\beta _{\mathrm{acc}}`$ as in zone 1, is not excluded, provided that the amount of radio electrons accelerated directly in the plateau region does not significantly exceed the number of electrons leaking into zone 2 from zone 1. Otherwise the gradients in the spatial densities of electrons would be reduced resulting then in a strong suppression of spectral modifications of $`N_1(E,t)`$ in the zone 1 regions.
The equation for the electron distribution $`N_2(E,t)`$ in zone 2 is similar to Eq.(1), where the source function $`Q`$ is substituted by $`Q_2=N_1/\tau _{\mathrm{esc}}`$ (taking into account that effectively the escape of relativistic electrons from zone 1 corresponds to their injection into zone 2, and neglecting possible acceleration in zone 2), and the sign of the term $`V_1N_2/V_2\tau _{\mathrm{dif}}`$ which describes the diffusive escape of electrons from zone 2 into zone 1 is changed. Note that a two-zone model does not assume any other escape of electrons from zone 2, e.g. into regions outside of the shell, because otherwise such an escape would correspond to a three-zone model (see Sect. 3 below).
Numerical calculations are done by the method of iterations using general analytical solutions to the Leaky-box type equations (e.g. see Atoyan & Aharonian 1999). For example, the solution for Eq.(1) can be presented as
$`N_1(E,t)`$ $`=`$ $`{\displaystyle \frac{1}{P_1(E)}}{\displaystyle _0^t}P_1(\zeta _t)Q_{\mathrm{eff}}(\zeta _t,t_1)\times `$ (5)
$`\mathrm{exp}\left({\displaystyle _{t_1}^t}{\displaystyle \frac{\mathrm{d}x}{\tau _{\mathrm{esc}}(\zeta _x)}}\right)\mathrm{d}t_1,`$
where $`Q_{\mathrm{eff}}(E,t)=Q+N_2/\tau _{\mathrm{esc}}`$. The variable $`\zeta _t`$ corresponds to the energy of an electron at instant $`t_1t`$ which has the energy $`E`$ at the instant $`t`$, and is determined from the equation
$$tt_1=_{E_\mathrm{e}}^{\zeta _t}\frac{\mathrm{d}E}{P_1(E)}.$$
(6)
At the first step of the iteration procedure a function $`N_1^{(1)}(E)`$ corresponding to the first approximation of $`N_1`$ is calculated for an injection function $`Q`$ in the form of Eq.(4), i.e. neglecting the term $`N_2/\tau _{\mathrm{esc}}`$ in $`Q_{\mathrm{eff}}(E)`$, or else - neglecting in Eq.(1) the term describing the diffusive ‘return’ flux of the electrons from zone 2 into zone 1. Then the energy distribution of the electrons $`N_2^{(1)}(E)`$ in zone 2 in the first approximation, corresponding to an effective injection function of the electrons in this zone $`N_1^{(1)}(E)/\tau _{\mathrm{esc}}(E)`$, is found. After that $`N_2^{(1)}(E)`$ is used for calculations of $`N_1^{(2)}(E)`$ in zone 1 in the second approximation, which now takes into account the ‘return’ flux of the electrons from zone 2. Calculations show that such a simple iteration procedure is quickly converging, so typically several iterations are sufficient to reach a good accuracy in the final electron distributions $`N_1`$ and $`N_2`$.
The characteristics of the broad band synchrotron radiation of Cas A can be best explained for the mean magnetic fields $`B_1(12)\mathrm{mG}`$ and $`B_2(0.30.5)\mathrm{mG}`$ in zones 1 and 2, respectively. The total energy content in relativistic electrons in each of these zones is of order $`10^{48}\mathrm{erg}`$, depending on the magnetic fields $`B_1`$ and $`B_2`$. The X-ray fluxes observed above $`10\mathrm{keV}`$ can be explained by synchrotron radiation of electrons in zone 1 with an exponential cutoff energy $`E_\mathrm{c}(1530)\mathrm{TeV}`$ in Eq.(4).
A possible fit to the broad band synchrotron fluxes of Cas A is shown in Fig. 1. Magnetic fields $`B_1=1.2\mathrm{mG}`$ and $`B_2=0.3\mathrm{mG}`$ in the compact bright radio structures and in the diffuse plateau region, respectively, are assumed. The total magnetic field energy in the shell is then $`W_{\mathrm{B2}}=3.8\times 10^{48}\mathrm{erg}`$, and in the compact zone 1 regions it is $`W_{\mathrm{B1}}=2\times 10^{47}\mathrm{erg}`$. The total energy of electrons in zones 1 and 2 is $`W_{\mathrm{e1}}=1.6\times 10^{48}\mathrm{erg}`$ and $`W_{\mathrm{e2}}=1.8\times 10^{48}\mathrm{erg}`$, respectively.
The energy distributions of the electrons $`N_1(E)`$ and $`N_2(E)`$ formed in zones 1 and 2, as well as the overall distribution $`N_{\mathrm{tot}}(E)=N_1(E)+N_2(E)`$, are shown in Fig. 2. As expected, $`N_{\mathrm{tot}}(E)`$ shows a pure power-law behavior with $`\beta =\beta _{\mathrm{acc}}`$ until the so-called ‘radiative break’ energy $`500\mathrm{GeV}`$. At this energy the synchrotron cooling time of the electrons in zone 2
$$t_\mathrm{s}300(E/1\mathrm{TeV})^1(B_2/0.2\mathrm{mG})^2\mathrm{yr}$$
(7)
becomes equal to the age of the source, and at higher energies the radiative losses steepen the spectrum of zone 2 electrons to the power-law index $`\beta =\beta _{\mathrm{acc}}+1`$. Meanwhile, at $`E100\mathrm{MeV}`$ the energy distribution of radio electrons in zone 1 is steepened to $`\beta _1\beta _{\mathrm{acc}}+\delta `$ because of diffusive escape of the electrons from that zone where their spatial density $`n_1(E)=N_1(E)/V_1`$ is much higher than $`n_2(E)=N_2(E)/V_2`$ in the surrounding plateau region.
### 2.2 Gamma-ray emission
The photon flux $`I(E)=J(E)/E`$ of the bremsstrahlung $`\gamma `$-rays at energies $`Em_\mathrm{e}c^2`$ practically repeats the power-law spectral shape of the parent electrons. Therefore, one might expect that $`I(E)`$ would have the same power-law index $`\alpha _{\mathrm{dif}}=\alpha +1\beta _{\mathrm{acc}}`$ as $`N_{\mathrm{tot}}(E)`$. However because of the different spectral shapes of the radio electrons in zones 1 and 2 where, most probably, the gas densities are also different, the total bremsstrahlung spectrum of Cas A can somewhat deviate from the power-law behavior of the injection spectrum.
In Fig.3 the thin solid line corresponds to the total flux of the bremsstrahlung photons produced by the electrons shown in Fig.2, and the dashed and dot-dashed lines represent the fluxes from zone 1 and zone 2, respectively. For the calculations we adopt a mean gas density $`n_{\mathrm{H},2}=15\mathrm{cm}^3`$ in terms of ‘H-atoms’ (i.e. the nucleons) in zone 2. This corresponds to a total mass of about $`15M_{}`$ (e.g. Fabian et al. 1980; Jansen et al. 1988; Vink et al. 1998), in the volume $`V1.2\times 10^{57}\mathrm{cm}^3`$ of the shell between $`R_{\mathrm{ring}}`$ and $`R_0`$. The mean value of the parameter $`C_\mathrm{Z}=Z(Z+1)/A`$ in the oxygen-rich ($`Z=8,A=16`$) gas in the shell (zone 2) of Cas A derived by Cowsik & Sarkar (1980) is $`\overline{C_\mathrm{Z}}=4.3`$. For zone 1 we assume the same atomic $`\overline{C_{\mathrm{Z},1}}=4.3`$, and the gas density $`n_{\mathrm{H},1}=4n_{\mathrm{H},2}`$. We note however that these parameters in the compact radio structures are not known. In particular, the radio knots show practically no optical line emission which would allow conclusions about the density and mass composition of the gas there. Therefore, for different values of the gas parameters in zone 1, the relative contribution of zones 1 and 2 to the total bremsstrahlung flux would change. As shown in Fig.4, this results in some uncertainty of the model predictions for the steepness of the total $`\gamma `$-ray fluxes at energies $`E100\mathrm{MeV}`$. At both smaller ($`10\mathrm{MeV}`$) and higher ($`1\mathrm{GeV}`$) energies, where the bremsstrahlung flux is contributed mainly either by the first or by the second zone, the spectral indices are much less affected by the uncertainty in the gas parameters in zone 1.
Because of the steep decline of the energy distribution of radio electrons in the compact zone 1 structures, the intensity of the $`\gamma `$-ray flux at $`E1\mathrm{GeV}`$ is dominated by the flat-spectrum bremsstrahlung of zone 2 (see Fig.3). It is also important that the contribution of the IC radiation, discussed in detail in the next section, at this energy is not yet significant. Therefore, measurements of the differential flux $`I(E)`$ of high energy $`\gamma `$-rays from Cas A by the future GLAST detector should allow a rather accurate determinations of a number of important parameters in Cas A. In particular, the spectral index of $`I(E)`$ at GeV energies could give rather accurate information about the spectrum of electrons in zone 2, $`\alpha _{\mathrm{dif}}\beta _2`$ (see Figs. 3 and 4). At these energies the energy distribution of electrons in zone 2 has a power law distribution $`N_2(E)E^{\beta _2}`$ with the spectral index of accelerated particles, $`\beta _2=\beta _{\mathrm{acc}}`$ (see Fig.2).
Bremsstrahlung $`\gamma `$-rays with energies $`E_\gamma 1\mathrm{GeV}`$ are produced by electrons with characteristic energies $`E_\mathrm{e}2E_\gamma `$. The mean frequency of the synchrotron photons emitted by the same electrons is $`\nu 10^3B_{\mathrm{mG}}(E_\mathrm{e}/m_\mathrm{e}c^2)^2`$ where $`B_{\mathrm{mG}}=B/1\mathrm{mG}`$ (e.g. Ginzburg 1979). In terms of the $`\gamma `$-ray energy this relation is reduced to
$$\nu 4B_{\mathrm{mG}}(E_\gamma /1\mathrm{GeV})^2\mathrm{GHz}.$$
(8)
For the expected mean magnetic field in the shell of Cas A of about (0.3–0.5) mG (Paper 1), electrons responsible for bremsstrahlung of (1-2) GeV $`\gamma `$-rays are also producing synchrotron photons with $`\nu `$ few GHz. Thus, a comparison of the radiation intensities in these two regions will give almost model independent information about the mean magnetic field $`B_2`$ in the shell.
Indeed, for the power-law distribution of electrons $`N(E)E^\beta `$, the intensity of synchrotron radiation can be written in the form (see Paper 1):
$`J_\nu `$ $``$ $`1.5\times 10^3\left({\displaystyle \frac{N_{}}{10^{50}}}\right)\left({\displaystyle \frac{B}{1\mathrm{mG}}}\right)^{\frac{1+\beta }{2}}\times `$ (9)
$`\left({\displaystyle \frac{\nu }{10\mathrm{GHz}}}\right)^{\frac{1\beta }{2}}\left({\displaystyle \frac{d}{3.4\mathrm{kpc}}}\right)^2\mathrm{Jy},`$
where $`N_{}E_{}N(E_{})`$, with $`E_{}=1\mathrm{GeV}`$, is the characteristic total number of 1 GeV electrons. The bremsstrahlung intensity I(E) (see e.g. Blumenthal and Gould 1970) can be presented in the form of an energy flux $`f(E)=E^2I(E)`$ as
$`f(E)`$ $``$ $`7.5\times 10^{15}n_\mathrm{H}\overline{C_\mathrm{Z}}\mathrm{\hspace{0.17em}2}^{2\beta }\left({\displaystyle \frac{N_{}}{10^{50}}}\right)\left({\displaystyle \frac{E}{1\mathrm{GeV}}}\right)^{2\beta }`$ (10)
$`\times \left[\mathrm{ln}\left({\displaystyle \frac{E}{1\mathrm{GeV}}}\right)+8.4\right]\left({\displaystyle \frac{d}{3.4\mathrm{kpc}}}\right)^2{\displaystyle \frac{\mathrm{erg}}{\mathrm{cm}^2\mathrm{s}}},`$
where $`\overline{C_\mathrm{Z}}=\overline{Z(Z+1)/A}`$, as previously. Then, because the total intensity of $`\gamma `$-rays at (1-2) GeV is dominated by the bremsstrahlung of radio electrons in the shell (zone 2) responsible also for the diffuse ‘radio plateau’ emission at $`\nu 5\mathrm{GHz}`$ with the known flux $`J_2(\nu )`$ <sup>1</sup><sup>1</sup>1 The measurements of Tuffs (1986) taken in 1978 have shown that $`50\%`$ of the total flux $`J(\nu )800\mathrm{Jy}`$ of Cas A was due to plateau emission. Taken into account the secular decline of the fluxes, the intensity of diffuse radio emission presently can be estimated as $`J_2(5\mathrm{GHz})350\mathrm{Jy}`$., and because the value of $`n_{\mathrm{H},2}\overline{C_{\mathrm{Z},2}}`$ in the shell is known with sufficient accuracy, comparison of equations (9) and (10) will result in a rather accurate estimate of the magnetic field $`B_2`$. For example, in the case of $`\alpha _{\mathrm{dif}}\beta _{\mathrm{acc}}2.2`$ as expected, the mean magnetic field in the shell would be
$$B_20.16\left[\frac{J_2(5\mathrm{GHz})}{350\mathrm{Jy}}\right]^{0.63}\left[\frac{f(1\mathrm{GeV})}{10^{11}\mathrm{erg}/\mathrm{cm}^2\mathrm{s}}\right]^{0.63}\mathrm{mG}$$
(11)
This knowledge of $`B_2`$ could then help to estimate the mean magnetic field also in compact bright radio structures as $`B_14B_2`$ which is needed for interpretation of the radio data (see Paper 1).
At energies $`100\mathrm{MeV}`$ the overall flux of $`\gamma `$-rays is dominated by the bremsstrahlung of relativistic electrons in the compact zone 1 region. As follows from Eq.(8), the radio counterpart of the $`E100\mathrm{MeV}`$ $`\gamma `$-rays is the region of frequencies $`\nu 40\mathrm{MHz}`$, where the synchrotron radiation is dominated by zone 1 emission (see Fig.1), and the fluxes are not yet affected by synchrotron self-absorption. Although the expected flux sensitivity of the GLAST detector drops significantly below 100 MeV, in principle it may be able to detect $`\gamma `$-ray fluxes down to energies $`10\mathrm{MeV}`$ (Bloom 1996). Detailed modelling of the flux $`I(E)`$ would help to extract the intensity of 100 MeV bremsstrahlung $`\gamma `$-rays produced in zone 1, i.e. in the bright radio ring and radio knots. Then, for known $`I_1(E100\mathrm{MeV})`$, $`J_1(\nu =40\mathrm{MHz}`$), and $`\beta =\beta _1(2.72.8)`$ in zone 1 (as predicted in Paper 1), the equations (9) and (10) can be used to derive the unknown product $`(n_\mathrm{H}\times \overline{C_\mathrm{Z}})`$ there. Note that in principle the gas parameters can be different in the radio ring and radio knots. A possibility of separation of the contributions of these two different sub-components of zone 1 to the overall flux $`I_1(E)`$ could enable determination of the product $`n_\mathrm{H}\overline{C_\mathrm{Z}}`$ in each of them.
Important additional information about the gas parameters in compact components could be derived from future observations of Cas A with high angular and energy resolution by the X-ray telescopes Chandra, XMM and Astro-E in the keV region. Even for the case of non-detection of the line features from radio knots, the intensity of their thermal X-ray emission will help to disentangle the density $`n_\mathrm{H}`$ from the mean abundance $`\overline{C_\mathrm{Z}}`$. Because the intensity of thermal emission $`Q_\mathrm{T}\overline{Z^2}n_\mathrm{Z}n_\mathrm{e}\overline{C_\mathrm{Z}}n_\mathrm{H}n_\mathrm{e}`$, and because for the ionized gas the thermal electron density $`n_\mathrm{e}=Zn_\mathrm{Z}n_\mathrm{H}/2`$ (except for a hydrogen-dominated medium where $`n_\mathrm{e}n_\mathrm{H}`$), the fluxes of thermal X-rays will in principle allow determination of the parameter $`\overline{C_\mathrm{Z}}n_\mathrm{H}^2`$ in the zone 1 structures. When combined with the knowledge of a different product of parameters $`\overline{C_\mathrm{Z}}`$ and $`n_\mathrm{H}`$ found from the comparison of $`\gamma `$-ray with synchrotron measurements, each of these two parameters could be found.
## 3 Gamma-ray emission at very high energies
### 3.1 Emission of relativistic electrons
If the fluxes of hard X-rays observed at $`E10\mathrm{keV}`$ have a synchrotron origin (Allen et al. 1997, Favata et al. 1997; – but see also Laming 1998), relativistic electrons in Cas A should be accelerated to energies up to tens of TeV. These electrons should then produce very high energy (VHE, $`E100\mathrm{GeV}`$) $`\gamma `$-rays. Along with bremsstrahlung, at these energies the principal mechanism for $`\gamma `$-ray production is the inverse Compton (IC) scattering of the electrons in the ambient soft photon field. In principle, the photon field is contributed by the synchrotron photons, the thermal dust emission with $`T=97\mathrm{K}`$ (Mezger et al. 1986) in the far infrared (FIR), the optical/IR line photons, and the 2.7 K cosmic microwave background radiation.
As shown in Fig.5, the most important target photon field for production of IC $`\gamma `$-rays in Cas A is the FIR radiation which is responsible for $`80\%`$ of the IC flux in the VHE region. Note that only due to the high density of the FIR radiation in Cas A (which has not been taken into account in recent calculations by Ellison et al. 1999), the flux of IC $`\gamma `$-rays becomes comparable at TeV energies with the bremsstrahlung flux. Because the photon fields should have practically the same density in both compact and diffuse zones, irrespective of where they are produced, and since the VHE electrons reside mostly in zone 2 (see Fig. 2), IC radiation is contributed mostly by this zone. For an assumed mean magnetic field in the shell $`B_2=0.3\mathrm{mG}`$ (Fig.5), the energy density of the magnetic field is $`w_\mathrm{B}=2.5\times 10^3\mathrm{eV}/\mathrm{cm}^3`$, whereas the density of the FIR radiation calculated for the luminosity $`L_{\mathrm{FIR}}3.6\times 10^{37}\mathrm{erg}/\mathrm{s}`$ (Mezger et al 1986) is only $`w_{\mathrm{rad}}2\mathrm{eV}/\mathrm{cm}^3`$. Because the ratio of synchrotron to IC (in the Thomson limit) emissivities
$$q_{\mathrm{sy}}/q_{\mathrm{IC}}=w_\mathrm{B}/w_{\mathrm{rad}}B^2,$$
(12)
the IC fluxes of VHE $`\gamma `$-rays are $`10^3`$ times lower than the fluxes of synchrotron radiation produced by the same parent electrons in the UV/X-ray region.
The fluxes of IC $`\gamma `$-rays to be expected from Cas A depend very sensitively on the mean magnetic field $`B_2`$ in the shell. As follows from Eq.(9), an increase of $`B_2`$ by a factor $`a`$ reduces the total number of electrons needed for explanation of the radio fluxes observed by a factor $`a^{(1+\beta _2)/2}`$. This results in a strong decrease of the bremsstrahlung flux by the same factor, i.e. $`f_{\mathrm{br}}1/a^{(1+\beta _2)/2}`$. The same dependence on the magnetic field also holds for the intensity of IC emission which is produced by the electrons below the radiative break energy (the ‘knee’ around 500 GeV in Fig.2) for which the synchrotron cooling time is larger than the age $`t_0`$ of the source. For these electrons the energy distribution $`N(E)Q(E)\times t_0`$ repeats the injection spectrum $`Q(E)`$. However, VHE $`\gamma `$-rays are produced by electrons of higher energies. In that case the spectrum of electrons is $`N(E)Q(E)\times t_\mathrm{s}`$. Then, besides of a less powerful injection rate $`Qa^{(1+\beta _2)/2}`$ needed for explanation of the radio fluxes, the spectral intensity of these electrons is additionally suppressed by a factor $`a^2B_2^2`$. Therefore the intensity of VHE radiation depends on the magnetic field in zone 2 as $`fB_2^{(5+\beta _2)/2}`$, which is significantly stronger than at GeV energies. This is the result of emission of VHE electrons in the “saturation” regime, when all the power injected in TeV electrons is channeled into (mostly) synchrotron and IC fluxes in the proportion defined by Eq.(12).
Thus, the flux of IC $`\gamma `$-rays could be significantly increased assuming smaller values of $`B_2`$, and hence also of $`B_1`$ because the ratio of magnetic fields in two zones should be approximately at the level $`B_1/B_24`$ in order to explain the radio data (see Paper 1). On the other hand, $`B_2`$ cannot be significantly less than 0.3 mG, otherwise the bremsstrahlung flux would then exceed the radiation flux observed at $`E100\mathrm{keV}`$, and the flux upper limit of EGRET at $`E100\mathrm{MeV}`$ (see Fig. 3). Note that this constraint on the magnetic field in the shell of Cas A imposed by the EGRET data are model independent since $`E100\mathrm{MeV}`$ $`\gamma `$-rays are produced by the same electrons with energies $`E_\mathrm{e}1\mathrm{GeV}`$ which are responsible for the observed radio fluxes. On the other hand, the constraints imposed by the hard X-ray fluxes are to some extent conditional, being based on a (resonable) assumption that the power-law injection spectrum starts from sub-MeV energies.
It is worthwhile to compare the constraints for the mean magnetic field in the shell, i.e. $`B_20.3\mathrm{mG}`$, with the constraints imposed in the framework of a spatially homogeneous (i.e. a single-zone) model commonly used. Because radiative losses cannot steepen the spectrum of the radio electrons in Cas A, both the acceleration spectrum and the overall energy distribution of electrons $`N(E)`$ at GeV energies should be steep, with a power-law index $`\beta _{\mathrm{acc}}2.5`$ for the mean power-law index of the observed radio fluxes $`\alpha 0.77`$. This is obviously much steeper than the index $`\beta _{\mathrm{acc}}2.2`$ predicted by the two-zone model. As a result, the overall number of electrons predicted in the framework of such single-zone approach at low energies is much higher than for the two-zone model, and hence the lower limit for the mean magnetic field in the shell of Cas A consistent with the X-ray fluxes and upper flux limit increases to $`B_{\mathrm{min}}1\mathrm{mG}`$ (e.g. Ellison et al. 1999). Both the steepness of the spectrum and the higher magnetic field predicted in the framework of a single-zone model significantly reduce the fluxes of $`\gamma `$-rays to be expected at very high energies. Note that the impact of the steepness of the acceleration spectrum $`Q(E)`$ can be significantly reduced in the framework of a more elaborated model where $`Q(E)`$ becomes significantly flatter at higher electron energies, as in calculations by Ellison et al. (1999). However, there is no way to reduce the impact of the high mean magnetic field $`1\mathrm{mG}`$ which essentially reduces the number of multi-TeV electrons needed for production of a given flux of synchrotron X-rays. Therefore, in the framework of a single-zone model the fluxes of VHE radiation which would not contradict the observed hard X-ray flux will be significantly lower than the fluxes to be expected in the framework of a spatially inhomogeneous two-zone model (compare Fig. 5 with the results of Ellison et al. 1999, and Goret et al. 1999).
For the same values of the mean magnetic fields in zones 1 and 2, the fluxes of IC $`\gamma `$-rays expected from Cas A at TeV energies could be higher than in Fig. 5 if we consider a more structured model for the magnetic field distribution in the shell than the two-zone model. Namely, we may assume that the magnetic field in the shell of Cas A may decrease from the highest value $`B_1`$ in the compact zone 1 (the acceleration sites) to a lower value $`B_2`$ in the surrounding zone 2, and further on to some $`B_3B_2`$ in zone 3. The chain of equations describing the energy distributions of particles in such 3-zone model is easily derived in the same way as described in Paper 1 for the 2-zone case. Note that the 3-zone modelling may be more adequate to the radio pattern of Cas A which shows significant variations in the brightness of the diffuse emission of the shell. Zone 3 would then represent the regions of the shell with relatively low magnetic field, as well as possibly the regions adjacent to the shell. Relativistic electrons could then escape from zone 2 into zone 3, as they do from zone 1 into zone 2. Because the spatial sizes of zone 2 are significantly larger than the sizes of the compact zone 1 structures, the characteristic timescales for the electron escape from zone 2 into zone 3 should be significantly larger than that the escape time from zone 1 into zone 2.
In Fig.6 we show the synchrotron radiation spectra calculated in the framework of the 3-zone model assuming that zones 2 and 3 have approximately the same volume filling factors in the shell. Note that zone 3 may include also regions not contained by the shell, in particular, the volume interior to the reverse shock if the diffusion coefficient there is sufficiently high (e.g. because of a low magnetic field there) so that relativistic particles would be able to significantly penetrate upstream of the freely expanding ejecta. In any case, zone 3 gives a principal possibility for high energy electrons to escape from zone 2 into a region with a lower magnetic field, resulting in some steepening of $`N_2(E)`$ at $`E10\mathrm{GeV}`$. Then it becomes possible to assume a power-law injection spectrum of the accelerated particles (in zone 1) harder than $`\beta _{\mathrm{acc}}2.2`$, without an excess of the radiation fluxes measured at $`6\mu \mathrm{m}`$ (Tuffs et al., 1997).
In Fig. 7 we show the integral fluxes of IC $`\gamma `$-rays, in terms of $`E\times I(>E)`$, calculated in the framework of the 3-zone model, assuming 2 different values for the magnetic field in zone 1, $`B_1=1\mathrm{mG}`$ (solid line) and $`B_1=1.6\mathrm{mG}`$ (dashed line), the magnetic field $`B_2=B_1/4`$ in zone 2, and $`B_3=0.1\mathrm{mG}`$ in zone 3. Although $`B_1`$ and $`B_2`$ in these 2 cases change only by a factor 1.6, the fluxes of TeV $`\gamma `$-rays drop dramatically, by a factor of 6-7. Note that an assumption of the magnetic field $`B_3`$ for zone 3 smaller than 0.1 mG does not result in a further increase of TeV $`\gamma `$-ray fluxes, because for such low magnetic fields all electrons up to several TeV are not in the “saturation” regime (see Eq.7), therefore variations of $`B_3`$ do not affect the electron energy distribution $`N_3(E)`$.
It should be noted in connection with the 3-zone model, as compared with the 2-zone model,that it allows an increase of the $`\gamma `$-ray fluxes at energies above 1 TeV (where the IC radiation dominates the overall flux from electrons) by a factor of 3 assuming the same magnetic fields $`B_1`$ and $`B_2`$ for zones 1 and 2. Because the model parameters in zone 3 (as the magnetic field $`B_3`$, the volume $`V_3`$ of zone 3 with this low field, and the escape time from zone 2 into zone 3) are difficult to deduce from radio observations, this introduces an additional uncertainty in the model predictions for the fluxes of VHE radiation. However, the mean magnetic fields $`B_2`$ and $`B_1`$ can be deduced rather accurately from future $`\gamma `$-ray detections of Cas A at lower energies (where bremsstrahlung dominates) as discussed in Sect. 2 . Then in principle the spectral measurements of VHE $`\gamma `$-ray fluxes could give an important information about the parameters of zone 3 with low magneic field.
The fluxes of TeV $`\gamma `$-rays on the level down to few per cent of the Crab nebula flux are in principle accessible for a system of Imaging Atmospheric Cherenkov Telescopes (IACTs) like the presently operating HEGRA. In this respect, the recent report of the HEGRA collaboration (Pühlhofer et al. 1999) about possible detection, at $`5\sigma `$ significance level, of a weak signal from Cas A appears very interesting, and needs a further confirmation. It should be said that the extraction of a signal from Cas A would require special care, because the fluxes of IC $`\gamma `$-rays are expected to decline already at energies $`E0.5\mathrm{TeV}`$ much faster than the flux of the Crab Nebula, which is generally considered as a “standard candle”. This can be seen in Fig. 7, as well as in Fig. 4 where the spectral index of the differential flux is plotted (to be compared with $`\alpha _{\mathrm{dif}}1.51.6`$ for the Crab Nebula). In Fig. 7 we also show the expected range of flux sensitivities of the forthcoming IACT arrays. It predicts that a significant flux of VHE $`\gamma `$-rays should be observed by the future VERITAS array (which is to operate in the northern hemisphere), if the observed fluxes of hard X-rays above 10 keV are indeed of synchrotron origin (Allen et al. 1997; Favata et al. 1997). Detection of the VHE $`\gamma `$-rays would allow a rather robust estimate of the mean magnetic fields $`B_1`$ and $`B_2`$. This could be done by a modelling, rather than by a direct comparison of the X-ray and TeV $`\gamma `$-ray fluxes, because they are produced in different zones – in compact zone 1 and the diffuse ‘radio plateau’, respectively.
### 3.2 $`\pi ^0`$-decay gamma-rays
The TeV $`\gamma `$-radiation in Cas A could also be efficiently produced by relativistic protons and nuclei which should be accelerated simultaneously with the electrons. Our study of the spectral and morphological characteristics of the radio emission of Cas A does not favour the blast wave as an efficient accelerator of an electrons competitive with electron acceleration in the compact structures. However, there are no comparable observational constraints on the proton acceleration sites. In particular, the protons could plausibly be accelerated also at the blast wave. The standard theory of diffusive shock acceleration suggests that the hadronic component of cosmic rays could be accelerated much more copiously than the electrons, which generally explains the up to two orders of magnitude overabundance of the nucleonic component in the observed galactic cosmic rays.
In Fig.7 we show by full dots the $`\gamma `$-ray fluxes resulting from the decay of $`\pi ^0`$-mesons produced by relativistic protons in the inelastic interactions with surrounding gas in the shell of Cas A. A power-law injection spectrum of relativistic protons in the form of Eq.(4), with $`\beta _{\mathrm{acc}}=2.15`$ as for the electrons, is supposed. For the characteristic maximum energy of accelerated protons we have assumed $`E_\mathrm{c}=200\mathrm{TeV}`$. Given the very high magnetic field in the acceleration region, $`B_11\mathrm{mG}`$, even for a rather young age of Cas A the relativistic protons, not being affected by synchrotron losses like the electrons, can easily reach such high values of $`E_\mathrm{c}`$. Indeed, using an estimate for the standard ‘parallel shock’ acceleration efficiency, the characteristic maximum energy of the protons accelerated during time $`t_0`$ (e.g. see Lagage & Cesarsky, 1983) can be written as
$`E_\mathrm{c}`$ $``$ $`450\left({\displaystyle \frac{B}{1\mathrm{mG}}}\right)\left({\displaystyle \frac{t_0}{100\mathrm{yr}}}\right)\times `$ (13)
$`\left({\displaystyle \frac{u_\mathrm{s}}{3000\mathrm{km}/\mathrm{s}}}\right)^2\eta ^1\mathrm{TeV},`$
where $`u_\mathrm{s}3000\mathrm{km}/\mathrm{s}`$ is a typical shock speed in Cas A (in particular, of the reverse shock), and $`\eta 1`$ is the so called gyrofactor (the ratio of the mean free path of a particle to its gyroradius). Thus, in the case of shock acceleration in the regime close to the Bohm diffusion limit, $`\eta 1`$, the protons in Cas A could reach the energy of order 500 TeV during an acceleration time as short as $`t_0100\mathrm{yr}`$.
The total energy of relativistic protons assumed for the calculation of $`\pi ^0`$-decay $`\gamma `$-ray fluxes in Fig. 7 is $`W_\mathrm{p}=2\times 10^{49}\mathrm{erg}`$, which corresponds to the mean injection power of the protons during $`t_0300\mathrm{yr}`$ of about $`L_\mathrm{p}2\times 10^{39}\mathrm{erg}/\mathrm{s}`$. On the other hand, for the magnetic field $`B_1=1\mathrm{mG}`$ the acceleration power of the electrons is $`L_\mathrm{e}=5.5\times 10^{38}\mathrm{erg}/\mathrm{s}`$, and $`L_\mathrm{e}=2.5\times 10^{38}\mathrm{erg}/\mathrm{s}`$ for $`B_1=1.6\mathrm{mG}`$. Thus, the protons are supposed to be accelerated only by a factor of $`48`$ more effectively than the electrons. This ratio cannot be increased further by more than a factor of 2, because otherwise the fluxes of TeV $`\pi ^0`$-decay $`\gamma `$-rays would exceed the flux upper limits shown in Fig. 7. Therefore the ratio of proton to electron acceleration efficiencies in Cas A is significantly smaller than, or else have not yet reached, the high value $`40`$ which is usually supposed for a typical source of the Galactic cosmic rays.
It is worth noticing that the upper limit to $`W_\mathrm{p}`$ imposed by the non-detection of TeV radiation from Cas A by Whipple and CAT detectors, is by one order of magnitude lower than the limitation imposed by the flux upper limit of the EGRET detector (see Fig. 7). Most probably, the total energy in accelerated protons in Cas A can be already now limited by a value not significantly exceeding $`10^{49}\mathrm{erg}`$. Of course, we cannot exclude that relativistic protons are accelerated only to energies below the TeV range. This would then imply that the electrons as well are not accelerated to these energies, and therefore that the hard X-ray fluxes detected from Cas A are not of a synchrotron origin (which cannot be still ruled out, see Laming 1998). Thus, we expect that rather important conclusions could be derived already from the fact of, hopefully, further confirmation of the possible HEGRA detection of TeV $`\gamma `$-rays Cas A.
For both hadronic and electronic origin of the TeV $`\gamma `$-rays, the fluxes are to be produced in the extended shell of Cas A enclosed between angular radii 1.5 and 2.5 arcmin, or perhaps even in a bit larger region, if the VHE electrons would escape from the shell (this may compose a part of zone 3, as discussed in previous section). In terms of diameter, this makes an appreciable size of order of 5 arcmin. The IACT arrays are able to reconstruct the direction of individual $`\gamma `$-rays with an accuracy of several arcminutes which, in combination with a significant statistics of the $`\gamma `$-rays, may result in a principal possibility to localize relatively strong point sources with an accuracy of about 1 arcmin (Aharonian et al. 1997). Therefore IACT arrays operating at energies $`100\mathrm{GeV}`$ might be able in principle to see some structure (possible ‘hot spots’) in the expected generally circular VHE $`\gamma `$-ray image of the source.
In this regard it might be worthwhile to note that radio observations of Kassim et al. (1995) at low frequencies show definite signs of the presence of a substantial mass (at least several $`M_{}`$) of cold gas in the central $`1^{}`$ region of the Cas A. If TeV particles would be able to penetrate into that region (by diffusive propagation upstream of the freely expanding ejecta – in the case of high diffusion coefficient), then a central ‘hot spot’ would appear, associated with the bremsstrahlung and $`\pi ^0`$-decay $`\gamma `$-rays produced by TeV electrons and protons illuminating the cold unshocked ejecta. The mean gas density of this ejecta in the central $`1^{}`$ region of Cas A can be estimated as $`n_\mathrm{H}10\mathrm{cm}^3`$ per each solar mass of the ejecta (to be compared with the mean $`n_\mathrm{H}15\mathrm{cm}^3`$ in the shell). Therefore, if the multi-TeV particles could really diffuse into that region, so that, say, 5 per cent of the total amount of these particles (protons and/or electrons) would be concentrated in that region, and the mass of freely expanding ejecta there could be about $`6\mathrm{M}_{}`$, then one could expect that up to $`20\%`$ of the total flux of VHE $`\gamma `$-rays would come from the centre of Cas A. Note that at energies $`100\mathrm{GeV}`$ where the sensitivity of IACT arrays is very high, most of the $`\gamma `$-ray flux produced by the electrons has a bremsstrahlung origin (see Fig.5).
Another possible ‘hot spot’ of the IC origin could appear (at higher energies) because of the known asymmetric spatial distribution of the target IR photons near the spectral peak around $`30\mu \mathrm{m}`$ (see $`25\mu \mathrm{m}`$ map by Dwek et al. 1987). If TeV electrons are distributed rather homogeneously in the shell, then the expected overall IC flux would not significantly change. However the image of Cas A at TeV energies will shift towards the position of maximum IR emission in the northern shell. If, on the other hand, the density of TeV electrons were enhanced in the northern shell (where the radio brightness is also enhanced) then we could expect an enhancement of the overall IC flux as well.
## 4 Conclusions
Although not yet conclusively detected, $`\gamma `$-rays should inevitably be produced in Cas A by relativistic electrons which are responsible for the synchrotron radiation in the radio to submillimeter wavelengths. Our model calculations predict (Fig. 3) that the bremsstrahlung flux of these electrons should be observed at least at energies between 10 MeV and 30 GeV by the future GLAST detector. It is also possible that instruments like INTEGRAL or Astro-E will be able to detect the flux of bremsstrahlung photons at energies $`100\mathrm{keV}`$ (produced predominantly in the compact radio structures) from Cas A. In that case one expects to see a profound hardening of the radiation spectrum above 100 keV. But this effect could be detected only if the mean magnetic field $`B_1`$ would not significantly exceed 1 mG, since otherwise the fluxes of soft $`\gamma `$-rays would fall below the limits of detectability of these instruments.
An important prediction of the spatially inhomogeneous model is that the spectra of accelerated electrons in Cas A correspond to a source function with a rather hard power-law index $`\beta _{\mathrm{acc}}2.2`$ (or even slightly harder), in agreement with predictions for efficient shock acceleration, and not a value of $`\beta _{\mathrm{acc}}2.5`$ as presently often assumed on the basis of the mean spectral index of the observed radio fluxes $`\alpha 0.77`$. This prediction of the model can be directly tested by future detection of the $`\gamma `$-ray fluxes by GLAST in the energy region around $`1\mathrm{GeV}`$. At these energies the total flux is dominated by the bremsstrahlung emission of electrons which escape from the compact acceleration regions into the surrounding shell, resulting in a power law distribution of electrons in the shell with an index $`\beta _2=\beta _{\mathrm{acc}}`$. Therefore, the spectral shape of the radiation detected at these energies will give direct information on the spectrum of acceleration in Cas A. The intensity of this radiation will give a rather robust estimate of the mean magnetic field $`B_2`$ in the shell regions responsible for the diffuse ‘radio plateau’ emission. This would allow also a rather good estimate of the mean magnetic field in the compact radio structures, $`B_14B_2`$.
Very important information on the magnetic field and chemical abundance of the gas in the compact radio structures can be derived by measuring the radiation fluxes at energies below 100 MeV. Although the expected angular resolution of GLAST may be insufficient to resolve the radio ring or radio knots, our study shows that most of the flux at these energies should be produced in these compact structures. Therefore, in combination with the known radio fluxes at low frequencies $`\nu 40\mathrm{GHz}`$, the $`\gamma `$-ray spectrum of Cas A at energies 10-100 MeV will provide an important information on the product of the gas parameters $`\overline{C_{\mathrm{Z},1}}`$ and $`n_{\mathrm{H},1}`$ in those structures. Adding also the information to be expected from the future high angular resolution measurements of the fluxes of thermal X-rays from the compact radio structures, even in the case of absence of line emission features it would be possible to disentangle $`n_{\mathrm{H},1}`$ and $`\overline{C_{\mathrm{Z},1}}`$.
A synchrotron origin of the X-ray fluxes above 10 keV implies acceleration of the electrons beyond 10 TeV. If so, we could expect also noticeable fluxes of VHE $`\gamma `$-rays. These fluxes depend very sensitively on the mean magnetic field $`B_1`$ in the compact regions of particle acceleration – the bright radio ring which is presumably connected with the reverse shock, and the radio knots –, and the field $`B_2`$ in the diffuse radio plateau. The model prediction for these fields (based on interpretation of synchrotron fluxes only) corresponds to $`B_1(12)\mathrm{mG}`$ and $`B_2B_1/4`$. For the case of low magnetic fields, $`B_11\mathrm{mG}`$, the integral flux $`I(>E)`$ above 1 TeV makes about $`5\%`$ of the Crab Nebula flux, and about $`(78)\%`$ above $`300\mathrm{GeV}`$. The $`\gamma `$-ray spectrum produced by electrons due to bremsstrahlung and IC radiation is rather unusual: it is moderately hard at $`E100\mathrm{GeV}`$, with a photon index $`\alpha _{\mathrm{dif}}2.12.3`$, and quickly steepens to $`\alpha _\gamma >3`$ already at $`E1\mathrm{TeV}`$ (see Fig.4). An increase of the magnetic field by only a factor of 2 results in a dramatic drop of the TeV $`\gamma `$-ray fluxes by one order of magnitude. Nevertheless, due to their high flux sensitivity and low detection energy threshold $`E_{\mathrm{th}}100\mathrm{GeV}`$ (Aharonian et al. 1997), the IACT arrays should be able to detect the fluxes of VHE radiation of Cas A even in the case of high magnetic fields $`B_12\mathrm{mG}`$. The preliminary detection of Cas A by the less sensitive IACT system of HEGRA, reported by Pühlhofer et al (1999), strengthens this expectation.
In this regard, it is significant that in Cas A the expected $`\gamma `$-ray fluxes of hadronic origin should readily be distinguished from $`\gamma `$-rays of electronic origin due to the much steeper electronic spectra at energies above a few hundred GeV (see Fig.4). The spectral information on the VHE $`\gamma `$-ray flux can be obtained by the future VERITAS array, and might be also accessible to the HEGRA IACT system (depending on the magnetic fields $`B_1`$ and $`B_2`$) for $`\gamma `$-rays of the electronic origin, or if the total energy of relativistic protons in Cas A is not less than $`W_\mathrm{p}^{(\mathrm{max})}2\times 10^{49}\mathrm{erg}`$. In both cases the predicted flux above 500 GeV is not much less than $`(0.050.1)`$ ‘Crab’ (see Fig. 7). Note that the current upper flux limits reported by CAT (Goret et al. 1999) and Whipple (Lessard et al. 1999) telescopes put an upper limit to the total energy of relativistic protons in Cas A of $`W_\mathrm{p}^{(\mathrm{max})}5\times 10^{49}\mathrm{erg}`$ (assuming acceleration of protons beyond TeV energies). A non-detection of $`\gamma `$-ray fluxes in the region above 100 GeV by future IACT arrays would imply that the efficiency of hadron acceleration in Cas A does not exceed the efficiency of electron acceleration.
###### Acknowledgements.
The work of AMA has been supported through the Verbundforschung Astronomie/Astrophysik of the German BMBF under the grant No. 05-2HD66A(7). |
warning/0001/hep-lat0001005.html | ar5iv | text | # Perturbative Renormalization Factors of Baryon Number Violating Operators for Improved Quark and Gauge Actions in Lattice QCD
## I Introduction
While nucleon decay is one of the most exciting prediction from grand unified theories (GUTs) with and without supersymmetry, none of the decay modes have been experimentally detected up to now. Furthermore, the ongoing Super-Kamiokande experiment is now pushing the lower limit on the partial lifetimes of the nucleon by an order of magnitude from the previous measurements. In principle this would give a strong constraint on (SUSY-)GUTs, however, the uncertainties in the theoretical prediction of lifetimes due to poor knowledge of quantum effects at low energy such as hadron or SUSY scales obscure the direct impact of the experimental lifetime bound on the physics at the GUT scale. In particular, one of the main uncertainties has been found in the evaluation of the hadron matrix elements for the nucleon decays, for which various QCD models have given estimate differing by a factor of ten. Therefore precise determination of the nucleon decay matrix elements from first principles is required, for which Lattice QCD can play a crucial role.
Recently we carried out a model-independent calculation of the nucleon decay matrix elements employing the Wilson quark action and the plaquette gauge action in the quenched approximation. Although one naively expects the error in the discretization and quenching approximation which are the two main systematic errors, it would be desirable to reduce these unknown systematic errors for high precision calculation. As a step toward this goal we are required to reduce the scaling violation effects by improving quark and gauge actions. In this article we present perturbative results for renormalization factors of baryon number violating operators for improved quark and gauge actions: the $`O(a)`$-improved “clover” action originally proposed by Sheikholeslami and Wohlert and the gauge action improved by addition of six link loops to the plaquette term in the Symanzik approach and in the Wilson’s renormalization group approach. Values of the one-loop coefficients of the renormalization factors are numerically evaluated for combinations of general values of the clover coefficients in the quark action and some specific values of the coefficients of the six-link loop terms in the gauge action.
This paper is organized as follows. In Sec.II we give the improved quark and gauge actions on the lattice and their Feynman rules relevant for their calculation. Our calculational procedure of the renormalization factors for the baryon number violating operators is described in Sec.III, where we present expressions and numerical values for the one-loop coefficients of the renormalization factors. Our conclusions are summarized in Sec.IV.
The physical quantities are expressed in lattice units, and the lattice spacing $`a`$ is suppressed unless necessary. Throughout this paper we use the same notation for quantities defined on the lattice and their counterparts in the continuum. In case of any possibility of confusion, however, we shall make a clear distinction between them.
## II Actions and Feynman rules
For the gauge action we consider the following general form including the standard plaquette term and six-link loop terms:
$$S_{\mathrm{gauge}}=\frac{1}{g^2}\left\{c_0\underset{\mathrm{plaquette}}{}\mathrm{Tr}U_{pl}+c_1\underset{\mathrm{rectangle}}{}\mathrm{Tr}U_{rtg}+c_2\underset{\mathrm{chair}}{}\mathrm{Tr}U_{chr}+c_3\underset{\mathrm{parallelogram}}{}\mathrm{Tr}U_{plg}\right\}$$
(1)
with the normalization condition
$$c_0+8c_1+16c_2+8c_3=1,$$
(2)
where six-link loops are $`1\times 2`$ rectangle, a bent $`1\times 2`$ rectangle (chair) and a three-dimensional parallelogram obtained by multiplying the link variables
$$U_{n,\mu }=\mathrm{exp}\left(ig\underset{a}{}T^aA_\mu ^a(n+\widehat{\mu }/2)\right).$$
(3)
The free gluon propagator is obtained in Ref.:
$`D_{\mu \nu }(k)={\displaystyle \frac{1}{(\widehat{k}^2)^2}}\left[(1A_{\mu \nu })\widehat{k}_\mu \widehat{k}_\nu +\delta _{\mu \nu }{\displaystyle \underset{\sigma }{}}\widehat{k}_\sigma ^2A_{\nu \sigma }\right]`$ (4)
with
$`\widehat{k}_\mu `$ $`=`$ $`2\mathrm{s}\mathrm{i}\mathrm{n}\left({\displaystyle \frac{k_\mu }{2}}\right),`$ (5)
$`\widehat{k}^2`$ $`=`$ $`{\displaystyle \underset{\mu =1}{\overset{4}{}}}\widehat{k}_\mu ^2.`$ (6)
The matrix $`A_{\mu \nu }`$ satisfies
$`(\mathrm{i})`$ $`A_{\mu \mu }=0\mathrm{for}\mathrm{all}\mu ,`$ (7)
$`(\mathrm{ii})`$ $`A_{\mu \nu }=A_{\nu \mu },`$ (8)
$`(\mathrm{iii})`$ $`A_{\mu \nu }(k)=A_{\mu \nu }(k).`$ (9)
$`(\mathrm{iv})`$ $`A_{\mu \nu }(0)=1\mathrm{for}\mu \nu ,`$ (10)
and its expression is given by
$`A_{\mu \nu }(k)`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Delta }_4}}[(\widehat{k}^2\widehat{k}_\nu ^2)(q_{\mu \rho }q_{\mu \tau }\widehat{k}_\mu ^2+q_{\mu \rho }q_{\rho \tau }\widehat{k}_\rho ^2+q_{\mu \tau }q_{\rho \tau }\widehat{k}_\tau ^2)`$ (15)
$`+(\widehat{k}^2\widehat{k}_\mu ^2)(q_{\nu \rho }q_{\nu \tau }\widehat{k}_\nu ^2+q_{\nu \rho }q_{\rho \tau }\widehat{k}_\rho ^2+q_{\nu \tau }q_{\rho \tau }\widehat{k}_\tau ^2)`$
$`+q_{\mu \rho }q_{\nu \tau }(\widehat{k}_\mu ^2+\widehat{k}_\rho ^2)(\widehat{k}_\nu ^2+\widehat{k}_\tau ^2)+q_{\mu \tau }q_{\nu \rho }(\widehat{k}_\mu ^2+\widehat{k}_\tau ^2)(\widehat{k}_\nu ^2+\widehat{k}_\rho ^2)`$
$`q_{\mu \nu }q_{\rho \tau }(\widehat{k}_\rho ^2+\widehat{k}_\tau ^2)^2(q_{\mu \rho }q_{\nu \rho }+q_{\mu \tau }q_{\nu \tau })\widehat{k}_\rho ^2\widehat{k}_\tau ^2`$
$`q_{\mu \nu }(q_{\mu \rho }\widehat{k}_\mu ^2\widehat{k}_\tau ^2+q_{\mu \tau }\widehat{k}_\mu ^2\widehat{k}_\rho ^2+q_{\nu \rho }\widehat{k}_\nu ^2\widehat{k}_\tau ^2+q_{\nu \tau }\widehat{k}_\nu ^2\widehat{k}_\rho ^2)],`$
with $`\mu \nu \rho \tau `$ the Lorentz indices. $`q_{\mu \nu }`$ and $`\mathrm{\Delta }_4`$ are written as
$`q_{\mu \nu }`$ $`=`$ $`(1\delta _{\mu \nu })\left[1(c_1c_2c_3)(\widehat{k}_\mu ^2+\widehat{k}_\nu ^2)(c_2+c_3)\widehat{k}^2\right],`$ (16)
$`\mathrm{\Delta }_4`$ $`=`$ $`{\displaystyle \underset{\mu }{}}\widehat{k}_\mu ^4{\displaystyle \underset{\nu \mu }{}}q_{\nu \mu }+{\displaystyle \underset{\mu >\nu ,\rho >\tau ,\{\rho ,\tau \}\{\mu ,\nu \}=\mathrm{}}{}}\widehat{k}_\mu ^2\widehat{k}_\nu ^2q_{\mu \nu }(q_{\mu \rho }q_{\nu \tau }+q_{\mu \tau }q_{\nu \rho }).`$ (17)
In the case of the standard plaquette action, the matrix $`A_{\mu \nu }`$ is simplified as
$$A_{\mu \nu }^{\mathrm{plaquette}}=1\delta _{\mu \nu }.$$
(18)
For the quark action we consider the $`O(a)`$-improved quark action:
$`S_{\mathrm{quark}}`$ $`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle \frac{1}{2}}{\displaystyle \underset{\mu }{}}\left\{\overline{\psi }_n(r+\gamma _\mu )U_{n,\mu }\psi _{n+\widehat{\mu }}+\overline{\psi }_n(r\gamma _\mu )U_{n\widehat{\mu },\mu }^{}\psi _{n\widehat{\mu }}\right\}+(m_0+4r){\displaystyle \underset{n}{}}\overline{\psi }_n\psi _n`$ (20)
$`c_{\mathrm{SW}}{\displaystyle \underset{n}{}}{\displaystyle \underset{\mu ,\nu }{}}ig{\displaystyle \frac{r}{4}}\overline{\psi }_n\sigma _{\mu \nu }P_{\mu \nu }(n)\psi _n,`$
where we define the Euclidean gamma matrices in terms of the Minkowski matrices in the Bjorken-Drell convention: $`\gamma _j=i\gamma _{BD}^j`$ $`(j=1,2,3)`$, $`\gamma _4=\gamma _{BD}^0`$, $`\gamma _5=\gamma _{BD}^5`$ and $`\sigma _{\mu \nu }=\frac{1}{2}[\gamma _\mu ,\gamma _\nu ]`$. The field strength $`P_{\mu \nu }`$ in the “clover” term is given by
$`P_{\mu \nu }(n)`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{i=1}{\overset{4}{}}}{\displaystyle \frac{1}{2ig}}\left(U_i(n)U_i^{}(n)\right),`$ (21)
$`U_1(n)`$ $`=`$ $`U_{n,\mu }U_{n+\widehat{\mu },\nu }U_{n+\widehat{\nu },\mu }^{}U_{n,\nu }^{},`$ (22)
$`U_2(n)`$ $`=`$ $`U_{n,\nu }U_{n\widehat{\mu }+\widehat{\nu },\mu }^{}U_{n\widehat{\mu },\nu }^{}U_{n\widehat{\mu },\mu },`$ (23)
$`U_3(n)`$ $`=`$ $`U_{n\widehat{\mu },\mu }^{}U_{n\widehat{\mu }\widehat{\nu },\nu }^{}U_{n\widehat{\mu }\widehat{\nu },\mu }U_{n\widehat{\nu },\nu },`$ (24)
$`U_4(n)`$ $`=`$ $`U_{n\widehat{\nu },\nu }^{}U_{n\widehat{\nu },\mu }U_{n+\widehat{\mu }\widehat{\nu },\nu }U_{n,\mu }^{}.`$ (25)
From the quark action (20) we obtain the free quark propagator
$$S_q^1(p)=i\underset{\mu }{}\gamma _\mu \mathrm{sin}(k_\mu )+m_0+r\underset{\mu }{}(1\mathrm{cos}(p_\mu )).$$
(26)
In order to calculate renormalization factors of the baryon number violating operators up to one-loop level, we need the following vertices,
$`V_{1\mu }^a(p,q)=igT^a\left\{\gamma _\mu \mathrm{cos}\left({\displaystyle \frac{p_\mu +q_\mu }{2}}\right)ir\mathrm{sin}\left({\displaystyle \frac{p_\mu +q_\mu }{2}}\right)\right\},`$ (27)
$`V_{c1\mu }^a(p,q)=gT^ac_{\mathrm{SW}}{\displaystyle \frac{r}{2}}\left({\displaystyle \underset{\nu }{}}\sigma _{\mu \nu }\mathrm{sin}(p_\nu q_\nu )\right)\mathrm{cos}\left({\displaystyle \frac{p_\mu q_\mu }{2}}\right)`$ (28)
with $`p_\mu `$ incoming quark momentum and $`q_\mu `$ outgoing quark momentum. The first vertex originates from the Wilson quark action and the second one is the interaction due to the clover term. In the present calculation of the vertex corrections for the baryon number violating operators, the two-gluon vertices with quarks give no contribution.
We should note that the baryon number violating operators contain a charge conjugated field, whose action is obtained from eq.(20) with the replacement of
$$T^a(T^a)^T,$$
(29)
where the superscript $`T`$ means the transposed matrix. According to this change, the Feynman rule of the quark-gluon vertices in eqs.(27) and (28) should be modified for the charge conjugated field.
## III Renormalization factors for baryon number violating operators
### A Calculational procedure
We consider the following baryon number violating operators in the continuum and on the lattice:
$`\left(𝒪_{X,Y}^{\mathrm{cont}}\right)_\delta `$ $`=`$ $`ϵ^{abc}\left[(\overline{\psi }_1^c)^a\mathrm{\Gamma }_X(\psi _2)^b\right]\left[\mathrm{\Gamma }_Y(\psi _3)^c\right]_\delta ,`$ (30)
$`\left(𝒪_{X,Y}^{\mathrm{latt}}\right)_\delta `$ $`=`$ $`ϵ^{abc}[\{1+rm_0(1z)\}(\overline{\psi }_1^c)^a\mathrm{\Gamma }_X(\psi _2)^b`$ (34)
$`+z{\displaystyle \frac{r}{2}}\{(\overline{\psi }_1^c\stackrel{}{D/})^a\mathrm{\Gamma }_X(\psi _2)^b(\overline{\psi }_1^c)^a\mathrm{\Gamma }_X(\stackrel{}{D/}\psi _2)^b\}`$
$`z^2{\displaystyle \frac{r^2}{4}}(\overline{\psi }_1^c\stackrel{}{D/})^a\mathrm{\Gamma }_X(\stackrel{}{D/}\psi _2)^b]`$
$`\times [\{1+{\displaystyle \frac{r}{2}}m_0(1z)\}\mathrm{\Gamma }_Y(\psi _3)^cz{\displaystyle \frac{r}{2}}\mathrm{\Gamma }_Y(\stackrel{}{D/}\psi _3)^c]_\delta `$
with
$`\stackrel{}{D/}\psi _n`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\mu }{}}\gamma _\mu \left(U_{n,\mu }\psi _{n+\widehat{\mu }}U_{n\widehat{\mu },\mu }^{}\psi _{n\widehat{\mu }}\right),`$ (35)
$`\overline{\psi }_n^c\stackrel{}{D/}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\mu }{}}\left(\overline{\psi }_{n+\widehat{\mu }}^cU_{n,\mu }^T\overline{\psi }_{n\widehat{\mu }}^cU_{n\widehat{\mu },\mu }^{}\right)\gamma _\mu ,`$ (36)
where $`\overline{\psi }^c=\psi ^TC`$ with $`C=\gamma _4\gamma _2`$ is a charge conjugated field of $`\psi `$. Dirac structures are represented by $`\mathrm{\Gamma }_X\mathrm{\Gamma }_Y=P_RP_R,P_RP_L,P_LP_R,P_LP_L`$ with the right- and left-handed projection operators $`P_{R,L}=(1\pm \gamma _5)/2`$. The summation over repeated color indices $`a,b,c`$ is assumed.
Ultraviolet divergences of composite operators are regularized by the cutoff $`a^1`$ in the lattice regularization scheme, while this is achieved by a reduction of the space-time dimension from four in some continuum regularization schemes, where we consider the naive dimensional regularization (NDR) scheme and the dimensional reduction (DRED) scheme. Operators defined in different regularization schemes can be related by renormalization factors:
$$𝒪_{X,Y}^{\mathrm{cont}}(\mu )=Z_{\mathrm{diag}}(\mu a)𝒪_{X,Y}^{\mathrm{latt}}(a)+Z_{\mathrm{mix}}\stackrel{~}{𝒪}_{X,Y}^{\mathrm{latt}}(a)$$
(37)
with $`\mu `$ the continuum renormalization scale. The explicit chiral symmetry breaking due to the Wilson term in the quark action (20) causes the mixing between operators with different chiral structures, which is denoted by $`\stackrel{~}{𝒪}_{X,Y}^{\mathrm{latt}}`$.
Since QCD is a asymptotically free theory, $`Z_{\mathrm{diag},\mathrm{mix}}`$ are expected to be perturbatively calculable in terms of the coupling constant $`g^2/(16\pi ^2)`$ at high energy scales, which gives the following expressions,
$`Z_{\mathrm{diag}}(\mu a)`$ $`=`$ $`1+{\displaystyle \frac{g^2}{16\pi ^2}}\left[{\displaystyle \frac{3}{2}}\left((1z)r\mathrm{\Sigma }_0{\displaystyle \frac{8}{3}}\mathrm{ln}(\mu a)+\mathrm{\Delta }_\psi \right)+8\mathrm{ln}(\mu a)+\mathrm{\Delta }_{V,\mathrm{diag}}\right],`$ (38)
$`Z_{\mathrm{mix}}`$ $`=`$ $`{\displaystyle \frac{g^2}{16\pi ^2}}\mathrm{\Delta }_{V,\mathrm{mix}},`$ (39)
where $`\mathrm{\Sigma }_0`$ denotes the additive mass renormalization on the lattice, $`\mathrm{\Delta }_\psi `$ is a contribution from the wavefunction and $`\mathrm{\Delta }_V`$ is from the vertex function. $`\mathrm{\Sigma }_0`$ and $`\mathrm{\Delta }_\psi `$ are obtained by calculating the continuum and lattice quark self-energies $`\mathrm{\Sigma }^{\mathrm{cont},\mathrm{latt}}`$ and $`\mathrm{\Delta }_V`$ is from the continuum and lattice vertex functions $`\mathrm{\Lambda }^{\mathrm{cont},\mathrm{latt}}`$.
The quark self-energies in the continuum and on the lattice are defined through the inverse full quark propagators:
$`(S_q^{\mathrm{cont}})^1(p)`$ $`=`$ $`ip/+m{\displaystyle \frac{g^2}{16\pi ^2}}\mathrm{\Sigma }^{\mathrm{cont}}(p),`$ (40)
$`(S_q^{\mathrm{latt}})^1(p)`$ $`=`$ $`i{\displaystyle \underset{\mu }{}}\gamma _\mu \mathrm{sin}(p_\mu )+m_0+r{\displaystyle \underset{\mu }{}}(1\mathrm{cos}(p_\mu )){\displaystyle \frac{g^2}{16\pi ^2}}\mathrm{\Sigma }^{\mathrm{latt}}(p),`$ (41)
where we consider massless quark. The difference between $`\mathrm{\Sigma }^{\mathrm{cont}}`$ and $`\mathrm{\Sigma }^{\mathrm{latt}}`$, which originates from the regularization scheme dependence of the self-energy, gives the quark wavefunction renormalization factor,
$$\frac{8}{3}\mathrm{ln}(\mu a)+\mathrm{\Delta }_\psi =\frac{\mathrm{\Sigma }^{\mathrm{cont}}(p)}{ip/}|_{p=0}\frac{\mathrm{\Sigma }^{\mathrm{latt}}(p)}{ip/}|_{p=0}.$$
(42)
The additive quark mass renormalization is expressed as
$$m_0\frac{g^2}{16\pi ^2}\mathrm{\Sigma }_0$$
(43)
with
$$\mathrm{\Sigma }_0=\mathrm{\Sigma }^{\mathrm{latt}}(p=0).$$
(44)
Notice that the renormalization factor (39) is given for the case of $`m_0=g^2/(16\pi ^2)\mathrm{\Sigma }_0`$, where we consider the renormalization for massless quark.
The vertex functions up to one-loop level in the continuum and on the lattice are expressed in the following way,
$`\mathrm{\Lambda }_{X,Y}^{\mathrm{cont}}`$ $`=`$ $`ϵ^{abc}\mathrm{\Gamma }_X\mathrm{\Gamma }_Y+{\displaystyle \frac{g^2}{16\pi ^2}}ϵ^{abc}\mathrm{\Gamma }_X\mathrm{\Gamma }_YV_{X,Y}^{\mathrm{cont}},`$ (45)
$`\mathrm{\Lambda }_{X,Y}^{\mathrm{latt}}`$ $`=`$ $`ϵ^{abc}\mathrm{\Gamma }_X\mathrm{\Gamma }_Y+{\displaystyle \frac{g^2}{16\pi ^2}}ϵ^{abc}[\mathrm{\Gamma }_X\mathrm{\Gamma }_Y(V_{X,Y}^{\mathrm{latt}}+zV_{X,Y}^{\mathrm{latt}}z^2V_{X,Y}^{\prime \prime \mathrm{latt}})`$ (47)
$`+\stackrel{~}{\mathrm{\Gamma }}_X\stackrel{~}{\mathrm{\Gamma }}_Y(\stackrel{~}{V}_{X,Y}^{\mathrm{latt}}+z\stackrel{~}{V}_{X,Y}^{\mathrm{latt}}z^2\stackrel{~}{V}_{X,Y}^{\prime \prime \mathrm{latt}})],`$
where the number of prime in the superscript of the lattice vertex corrections denotes the number of covariant derivative applied to the quark fields at the vertex. $`\stackrel{~}{\mathrm{\Gamma }}_X\stackrel{~}{\mathrm{\Gamma }}_Y`$ term represents the mixing contribution. The difference between $`\mathrm{\Lambda }_{X,Y}^{\mathrm{cont}}`$ and $`\mathrm{\Lambda }_{X,Y}^{\mathrm{latt}}`$ leads to
$`8\mathrm{ln}(\mu a)+\mathrm{\Delta }_{V,\mathrm{diag}}`$ $`=`$ $`V_{X,Y}^{\mathrm{cont}}\left(V_{X,Y}^{\mathrm{latt}}+zV_{X,Y}^{\mathrm{latt}}z^2V_{X,Y}^{\prime \prime \mathrm{latt}}\right),`$ (48)
$`\mathrm{\Delta }_{V,\mathrm{mix}}`$ $`=`$ $`\left(\stackrel{~}{V}_{X,Y}^{\mathrm{latt}}+z\stackrel{~}{V}_{X,Y}^{\mathrm{latt}}z^2\stackrel{~}{V}_{X,Y}^{\prime \prime \mathrm{latt}}\right).`$ (49)
We note that the lattice quark-self energy and the lattice vertex corrections are general function of the clover coefficient $`c_{\mathrm{SW}}`$ in the quark action and the six-link loop parameters $`c_{1,2,3}`$ in the gauge action. Calculation of $`\mathrm{\Delta }_\psi `$ was already carried out in Ref. employing the general values for $`c_{\mathrm{SW}}`$. For $`c_{1,2,3}`$ they choose some specific values: $`c_1=1/12,c_2=c_3=0`$ in the tree-level Symanzik improvement, $`c_1=0.252,c_2=0,c_3=0.17`$ suggested by Wilson based on renormalization group improvement and $`c_1=0.331,c_2=c_3=0`$ and $`c_1=0.27,c_2+c_3=0.04`$ by Iwasaki. According to Ref. we evaluate $`\mathrm{\Delta }_V`$ for general values of $`c_{\mathrm{SW}}`$ and for the specific values of $`c_{1,2,3}`$ that they employed.
### B Vertex corrections
We calculate the vertex corrections of the operators in eqs.(30) and (34) in the Feynman gauge employing the massless quarks and and the massless charge conjugated quark with momenta $`p_1=p_2=p_3=0`$ as external states. The infrared singularities are regularized by a fictitious gluon mass $`\lambda `$ introduced in the gluon propagator.
One-loop vertex corrections on the lattice are illustrated in Fig.1. We find that the lattice vertex corrections is a second polynomial function of the clover coefficients $`c_{\mathrm{SW}}`$. The relevant diagrams for $`V_{X,Y}^{\mathrm{latt}}`$ and $`\stackrel{~}{V}_{X,Y}^{\mathrm{latt}}`$ are Figs.1(a)$``$(i), the sum of which gives
$`\mathrm{\Gamma }_{R/L}\mathrm{\Gamma }_YV_{R/L,Y}+\stackrel{~}{\mathrm{\Gamma }}_{R/L}\stackrel{~}{\mathrm{\Gamma }}_Y\stackrel{~}{V}_{R/L,Y}`$ (50)
$`=`$ $`\mathrm{\Gamma }_{R/L}\mathrm{\Gamma }_Y\left(C_{B/}6\mathrm{ln}\left|{\displaystyle \frac{1}{\lambda ^2a^2}}\right|+{\displaystyle \underset{i=0,1,2}{}}c_{\mathrm{SW}}^{}{}_{}{}^{(i)}v_{\mathrm{diag}}^{(i)}\right)`$ (52)
$`+\left(\mathrm{\Gamma }_{L/R}\mathrm{\Gamma }_Y\pm {\displaystyle \frac{1}{4}}{\displaystyle \underset{\mu }{}}\gamma _\mu \gamma _5\mathrm{\Gamma }_Y\gamma _\mu \right){\displaystyle \underset{i=0,1,2}{}}c_{\mathrm{SW}}^{}{}_{}{}^{(i)}v_{\mathrm{mix}}^{(i)}`$
with $`C_{B/}=(N+1)/(2N)`$ in the SU($`N`$) group. The explicit forms of $`v_{\mathrm{diag}}^{(i)}`$ and $`v_{\mathrm{mix}}^{(i)}`$ are given by
$`v_{\mathrm{diag}}^{(0)}`$ $`=`$ $`C_{B/}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{\pi ^2}}\left[{\displaystyle \frac{1}{F_0^2G_0}}\left\{48(\mathrm{\Delta }_3+r^2\mathrm{\Delta }_1^2)^2+8I_a+16I_b\right\}\theta (\pi ^2k^2){\displaystyle \frac{6}{(k^2)^2}}\right]`$ (54)
$`+C_{B/}6\mathrm{ln}|\pi ^2|,`$
$`v_{\mathrm{diag}}^{(1)}`$ $`=`$ $`C_{B/}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{\pi ^2}}\left[{\displaystyle \frac{1}{F_0^2G_0}}\left\{8r^2\mathrm{\Delta }_1I_a\right\}\right],`$ (55)
$`v_{\mathrm{diag}}^{(2)}`$ $`=`$ $`C_{B/}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{\pi ^2}}\left[{\displaystyle \frac{1}{F_0^2G_0}}\left\{4r^4\mathrm{\Delta }_1^2I_a\right\}\right],`$ (56)
$`v_{\mathrm{mix}}^{(0)}`$ $`=`$ $`C_{B/}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{\pi ^2}}\left[{\displaystyle \frac{1}{F_0^2G_0}}\left\{16r^2\mathrm{\Delta }_1^2(\mathrm{\Delta }_34\mathrm{\Delta }_{1,0}^\mu )\right\}\right],`$ (57)
$`v_{\mathrm{mix}}^{(1)}`$ $`=`$ $`C_{B/}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{\pi ^2}}\left[{\displaystyle \frac{1}{F_0^2G_0}}\left\{8r^2\mathrm{\Delta }_1I_a\right\}\right],`$ (58)
$`v_{\mathrm{mix}}^{(2)}`$ $`=`$ $`C_{B/}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{\pi ^2}}\left[{\displaystyle \frac{1}{F_0^2G_0}}\left\{4r^2\mathrm{\Delta }_3I_a\right\}\right],`$ (59)
where
$`F_0`$ $`=`$ $`{\displaystyle \underset{\mu }{}}\mathrm{sin}^2(k_\mu )+{\displaystyle \frac{r^2}{4}}(\widehat{k}^2)^2,`$ (60)
$`G_0`$ $`=`$ $`(\widehat{k}^2)^2,`$ (61)
$`I_a`$ $`=`$ $`\mathrm{\Delta }_{1,1}^\mu 4\mathrm{\Delta }_3^2+(16\mathrm{\Delta }_34\mathrm{s}\mathrm{i}\mathrm{n}^2(k_\mu ))\mathrm{\Delta }_{1,0}^\mu ,`$ (62)
$`I_b`$ $`=`$ $`\mathrm{\Delta }_{1,1}^\mu +\mathrm{\Delta }_3^2+4(\mathrm{\Delta }_3+\mathrm{sin}^2(k_\mu ))\mathrm{\Delta }_{1,0}^\mu ,`$ (63)
$`\mathrm{\Delta }_1`$ $`=`$ $`{\displaystyle \frac{1}{4}}\widehat{k}^2,`$ (64)
$`\mathrm{\Delta }_3`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{\mu }{}}\mathrm{sin}^2(k_\mu ),`$ (65)
$`\mathrm{\Delta }_{1,1}^\mu `$ $`=`$ $`{\displaystyle \underset{\nu }{}}(\delta _{\mu \nu }+A_{\mu \nu })\mathrm{sin}^2(k_\mu )\mathrm{sin}^2(k_\nu ),`$ (66)
$`\mathrm{\Delta }_{1,0}^\mu `$ $`=`$ $`{\displaystyle \underset{\nu }{}}(\delta _{\mu \nu }+A_{\mu \nu })\mathrm{cos}^2\left({\displaystyle \frac{k_\mu }{2}}\right)\mathrm{sin}^2\left({\displaystyle \frac{k_\nu }{2}}\right).`$ (67)
We do not take the sum over the index $`\mu `$ for $`\mathrm{\Delta }_{1,1}^\mu `$ and $`\mathrm{\Delta }_{1,0}^\mu `$. In a similar way we obtain the expressions of $`V_{X,Y}^{\mathrm{latt}}`$ and $`\stackrel{~}{V}_{X,Y}^{\mathrm{latt}}`$ from Figs.1(a)$``$(i):
$`\mathrm{\Gamma }_{R/L}\mathrm{\Gamma }_YV_{R/L,Y}^{}+\stackrel{~}{\mathrm{\Gamma }}_{R/L}\stackrel{~}{\mathrm{\Gamma }}_Y\stackrel{~}{V}_{R/L,Y}^{}`$ (68)
$`=`$ $`\mathrm{\Gamma }_{R/L}\mathrm{\Gamma }_Y{\displaystyle \underset{i=0,1,2}{}}c_{\mathrm{SW}}^{}{}_{}{}^{(i)}v_{\mathrm{diag}}^{(i)}`$ (70)
$`+\left(\mathrm{\Gamma }_{L/R}\mathrm{\Gamma }_Y\pm {\displaystyle \frac{1}{4}}{\displaystyle \underset{\mu }{}}\gamma _\mu \gamma _5\mathrm{\Gamma }_Y\gamma _\mu \right){\displaystyle \underset{i=0,1,2}{}}c_{\mathrm{SW}}^{}{}_{}{}^{(i)}v_{\mathrm{mix}}^{(i)},`$
where
$`v_{\mathrm{diag}}^{(0)}`$ $`=`$ $`C_{B/}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{\pi ^2}}\left[{\displaystyle \frac{1}{F_0^2G_0}}\left\{16r^2\mathrm{\Delta }_1I_a+32r^2\mathrm{\Delta }_1I_b\right\}\right]`$ (72)
$`+C_F{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{\pi ^2}}\left[{\displaystyle \frac{1}{F_0G_0}}\left\{12r^2\mathrm{\Delta }_1(\mathrm{\Delta }_34\mathrm{\Delta }_{1,0}^\mu )\right\}\right],`$
$`v_{\mathrm{diag}}^{(1)}`$ $`=`$ $`C_{B/}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{\pi ^2}}\left[{\displaystyle \frac{1}{F_0^2G_0}}\left\{8(r^2\mathrm{\Delta }_3r^4\mathrm{\Delta }_1^2)I_a\right\}\right]`$ (74)
$`+C_F{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{\pi ^2}}\left[{\displaystyle \frac{1}{F_0G_0}}\left\{3r^2I_a\right\}\right],`$
$`v_{\mathrm{diag}}^{(2)}`$ $`=`$ $`C_{B/}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{\pi ^2}}\left[{\displaystyle \frac{1}{F_0^2G_0}}\left\{8r^4\mathrm{\Delta }_1\mathrm{\Delta }_3I_a\right\}\right],`$ (75)
$`v_{\mathrm{mix}}^{(0)}`$ $`=`$ $`C_{B/}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{\pi ^2}}[{\displaystyle \frac{1}{F_0^2G_0}}\{32r^2\mathrm{\Delta }_1\mathrm{\Delta }_3(\mathrm{\Delta }_34\mathrm{\Delta }_{1,0}^\mu )\}`$ (77)
$`+{\displaystyle \frac{1}{F_0G_0}}\left\{8r^2\mathrm{\Delta }_1(\mathrm{\Delta }_34\mathrm{\Delta }_{1,0}^\mu )\right\}],`$
$`v_{\mathrm{mix}}^{(1)}`$ $`=`$ $`C_{B/}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{\pi ^2}}\left[{\displaystyle \frac{1}{F_0^2G_0}}\left\{8(r^2\mathrm{\Delta }_3r^4\mathrm{\Delta }_1^2)I_a\right\}+{\displaystyle \frac{1}{F_0G_0}}\left\{2r^2I_a\right\}\right],`$ (78)
$`v_{\mathrm{mix}}^{(2)}`$ $`=`$ $`C_{B/}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{\pi ^2}}\left[{\displaystyle \frac{1}{F_0^2G_0}}\left\{8r^4\mathrm{\Delta }_1\mathrm{\Delta }_3I_a\right\}\right]`$ (79)
with $`C_F=(N^21)/(2N)`$ in the SU($`N`$) group. The sum of Figs.1(a)$``$(k), which include the tadpole diagrams at the vertex, yields $`V_{X,Y}^{\prime \prime \mathrm{latt}}`$ and $`\stackrel{~}{V}_{X,Y}^{\prime \prime \mathrm{latt}}`$,
$`\mathrm{\Gamma }_{R/L}\mathrm{\Gamma }_YV_{R/L,Y}^{\prime \prime }+\stackrel{~}{\mathrm{\Gamma }}_{R/L}\stackrel{~}{\mathrm{\Gamma }}_Y\stackrel{~}{V}_{R/L,Y}^{\prime \prime }`$ (80)
$`=`$ $`\mathrm{\Gamma }_{R/L}\mathrm{\Gamma }_Y{\displaystyle \underset{i=0,1,2}{}}c_{\mathrm{SW}}^{}{}_{}{}^{(i)}v_{\mathrm{diag}}^{\prime \prime (i)}`$ (82)
$`+\left(\mathrm{\Gamma }_{L/R}\mathrm{\Gamma }_Y\pm {\displaystyle \frac{1}{4}}{\displaystyle \underset{\mu }{}}\gamma _\mu \gamma _5\mathrm{\Gamma }_Y\gamma _\mu \right){\displaystyle \underset{i=0,1,2}{}}c_{\mathrm{SW}}^{}{}_{}{}^{(i)}v_{\mathrm{mix}}^{\prime \prime (i)},`$
where
$`v_{\mathrm{diag}}^{\prime \prime (0)}`$ $`=`$ $`C_{B/}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{\pi ^2}}\left[{\displaystyle \frac{1}{F_0^2G_0}}\left\{8r^4\mathrm{\Delta }_1^2I_a16r^4\mathrm{\Delta }_1^2I_b\right\}\right],`$ (83)
$`v_{\mathrm{diag}}^{\prime \prime (1)}`$ $`=`$ $`C_{B/}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{\pi ^2}}\left[{\displaystyle \frac{1}{F_0^2G_0}}\left\{8r^4\mathrm{\Delta }_1\mathrm{\Delta }_3I_a\right\}\right],`$ (84)
$`v_{\mathrm{diag}}^{\prime \prime (2)}`$ $`=`$ $`C_{B/}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{\pi ^2}}\left[{\displaystyle \frac{1}{F_0^2G_0}}\left\{4r^4\mathrm{\Delta }_3^2I_a\right\}\right],`$ (85)
$`v_{\mathrm{mix}}^{\prime \prime (0)}`$ $`=`$ $`C_{B/}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{\pi ^2}}[{\displaystyle \frac{1}{F_0^2G_0}}\left\{4r^2(16\mathrm{\Delta }_3^2\mathrm{\Delta }_{1,0}^\mu +8r^2\mathrm{\Delta }_1^2\mathrm{\Delta }_3^2+4r^4\mathrm{\Delta }_1^4\mathrm{\Delta }_3)\right\}`$ (87)
$`+{\displaystyle \frac{1}{F_0G_0}}\{8r^2(4\mathrm{\Delta }_3\mathrm{\Delta }_{1,0}^\mu +r^2\mathrm{\Delta }_1^2\mathrm{\Delta }_3)\}+{\displaystyle \frac{1}{G_0}}\left\{4r^2\mathrm{\Delta }_{1,0}^\mu \right\}],`$
$`v_{\mathrm{mix}}^{\prime \prime (1)}`$ $`=`$ $`C_{B/}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{\pi ^2}}\left[{\displaystyle \frac{1}{F_0^2G_0}}\left\{8r^4\mathrm{\Delta }_1\mathrm{\Delta }_3I_a\right\}+{\displaystyle \frac{1}{F_0G_0}}\left\{2r^4\mathrm{\Delta }_1I_a\right\}\right],`$ (88)
$`v_{\mathrm{mix}}^{\prime \prime (2)}`$ $`=`$ $`C_{B/}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{\pi ^2}}\left[{\displaystyle \frac{1}{F_0^2G_0}}\left\{4r^6\mathrm{\Delta }_1^2\mathrm{\Delta }_3I_a\right\}\right].`$ (89)
We present numerical values of $`v_{\mathrm{diag},\mathrm{mix}}^{(i)}`$ in Table I, $`v_{\mathrm{diag},\mathrm{mix}}^{(i)}`$ in Table II and $`v_{\mathrm{diag},\mathrm{mix}}^{\prime \prime (i)}`$ in Table III, which are evaluated with $`r=1`$ using the Monte Carlo integration routine BASES for specific values of $`c_1`$ and $`c_2+c_3`$. The numerical accuracy is better than $`0.01\%`$. Numerical values for $`v_{\mathrm{diag},\mathrm{mix}}^{(i)}`$, $`v_{\mathrm{diag},\mathrm{mix}}^{(i)}`$ and $`v_{\mathrm{diag},\mathrm{mix}}^{\prime \prime (i)}`$ can be also obtained by using the results for vertex corrections of bilinear operators in Ref., in which the numerical values are evaluated in a different way. This is used as a check of our calculation. We note that a special case of $`c_{\mathrm{SW}}=0`$ and $`c_1=c_2+c_3=0`$ represents combination of the Wilson quark action and the plaquette gauge action, for which perturbative renormalization factors for the baryon number violating operators has been already calculated. Comparison of our results with theirs gives us another check of our calculation. Numerical values in Tables I, II and III show that the one-loop coefficients in the vertex corrections diminishes by $`1020\%`$ for the tree-level Symanzik action compared to those for the plaquette action. Further reduction of the magnitude is observed for the renormalization group improved actions. These features are also found in the case of bilinear operators.
In the continuum, the vertex correction at one-loop level is expressed as
$$V_{R/L,Y}=C_{B/}6\mathrm{ln}\left|\frac{\mu ^2}{\lambda ^2}\right|+v_{\mathrm{diag}}.$$
(90)
The finite constant $`v_{\mathrm{diag}}`$ is given by
$`v_{\mathrm{diag}}^{\mathrm{NDR}}={\displaystyle \frac{8}{3}},`$ (91)
$`v_{\mathrm{diag}}^{\mathrm{DRED}}=4`$ (92)
for the NDR and DRED schemes with $`\overline{\mathrm{MS}}`$ subtraction.
The vertex corrections on the lattice give the expression of $`\stackrel{~}{𝒪}_{X,Y}^{\mathrm{latt}}`$ in eq.(37),
$$\stackrel{~}{𝒪}_{R/L,Y}^{\mathrm{latt}}(a)=ϵ^{abc}\left[(\overline{\psi }_1^c)^a\mathrm{\Gamma }_{L/R}(\psi _2)^b\right]\left[\mathrm{\Gamma }_Y(\psi _3)^c\right]\pm \frac{1}{4}\underset{\mu }{}ϵ^{abc}\left[(\overline{\psi }_1^c)^a\gamma _\mu \gamma _5(\psi _2)^b\right]\left[\mathrm{\Gamma }_Y\gamma _\mu (\psi _3)^c\right].$$
(93)
Comparing the results for the vertex corrections in the continuum and on the lattice, we obtain the vertex correction components in the renormalization factors of eqs.(38) and (39),
$`\mathrm{\Delta }_{V,\mathrm{diag}}^{\mathrm{NDR},\mathrm{DRED}}`$ $`=`$ $`v_{\mathrm{diag}}^{\mathrm{NDR},\mathrm{DRED}}{\displaystyle \underset{i=0,1,2}{}}c_{\mathrm{SW}}^{}{}_{}{}^{(i)}(v_{\mathrm{diag}}^{(i)}+zv_{\mathrm{diag}}^{(i)}z^2v_{\mathrm{diag}}^{\prime \prime (i)}),`$ (94)
$`\mathrm{\Delta }_{V,\mathrm{mix}}`$ $`=`$ $`{\displaystyle \underset{i=0,1,2}{}}c_{\mathrm{SW}}^{}{}_{}{}^{(i)}(v_{\mathrm{mix}}^{(i)}+zv_{\mathrm{mix}}^{(i)}z^2v_{\mathrm{mix}}^{\prime \prime (i)}),`$ (95)
where $`\mathrm{\Delta }_{V,\mathrm{mix}}`$ is independent of the renormalization scheme in the continuum. To obtain the diagonal part of the renormalization factor in eq.(38), we also need the wavefunction component $`\mathrm{\Delta }_\psi `$. This quantity is already evaluated in Ref. employing the NDR scheme in the continuum, where $`C_Fz_\psi `$ in their notation corresponds to our $`\mathrm{\Delta }_\psi ^{\mathrm{NDR}}`$. We note that $`\mathrm{\Delta }_\psi ^{\mathrm{DRED}}`$ can be obtained from $`\mathrm{\Delta }_\psi ^{\mathrm{NDR}}`$ by
$$\mathrm{\Delta }_\psi ^{\mathrm{DRED}}=\mathrm{\Delta }_\psi ^{\mathrm{NDR}}\frac{4}{3}$$
(96)
For the mixing part of the renormalization factor in eq.(39), $`\mathrm{\Delta }_\psi `$ has no contribution.
## IV Conclusions
In this paper we have calculated the one-loop contributions for the renormalization factors of the three-quark operators employing the improved quark and gauge actions. Detailed numerical values of the one-loop coefficients are presented for general values of the clover coefficients and for some specific values of $`c_1`$, $`c_2`$ and $`c_3`$ in the gauge action improved by the Symanzik and renormalization group approaches. We find that the magnitude of the one-loop coefficients are considerably reduced for the improved gauge actions compared to those for the plaquette action, which is a desirable feature in the practical implementation of lattice QCD.
###### Acknowledgements.
This work is supported in part by the Grants-in-Aid of the Ministry of Education (No. 10740125). One of us (Y.K.) is supported by the JSPS Fellowship. |
warning/0001/hep-ph0001097.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The inclusive pion double charge exchange (DCX) reaction on nuclei, e.g.
$$\pi ^{}+A(Z,N)\pi ^++X$$
(1)
as a two-step transition on two like nucleons (protons) is described by the diagrams of three types (Fig.1(a), (b), (c)) which correspond to the Glauber picture of elastic (a), quasielastic (b) and inelastic (c) rescatterings.
The conventional DCX mechanism of two sequential single charge exchanges (SSCX), i.e. the elastic rescattering, is found to describe , in general, experimental data on pion DCX at energies of the meson factories up to 0.5 $`GeV`$. It predicts a strong decrease of forward angle exclusive pion DCX at incident kinetic energies $`T_0\begin{array}{c}>\hfill \\ \hfill \end{array}0.6GeV`$ due to the rapid energy drop of $`\pi ^{}\pi ^0`$ transition rate. The SSCX mechanism leads to analogous behaviour of inclusive forward DCX rate too.
Inclusive reaction (1) is a unique process where, as it was shown in , the effect of the Glauber inelastic rescatterings (IR) becomes very important already at $`GeV`$ energies. This effect can be expected from a comparison of experimental total cross sections for two competing processes, $`\pi ^{}p\pi ^0n`$ and $`\pi ^{}p\pi ^+\pi ^{}n`$ (Fig.2, where data are given from the compilation ). The study of the forward angle reaction (1) by ITEP experiments on $`{}_{}{}^{16}O`$ and $`{}_{}{}^{6}Li`$ in a kinematics which forbids real production of additional pion has demonstrated a relatively weak energy dependence of its rate and considerable excess over SSCX mechanism prediction around $`T_0=1GeV`$. So the conventional elastic SSCX picture of inclusive pion DCX has to be modified by an important new mechanism of IR at $`GeV`$ energies.
Our goal is to elaborate recently suggested Gribov-Glauber approach to inclusive pion DCX reaction on $`{}_{}{}^{16}O`$. Below we present an updated analysis of some theoretical aspects of our model and clarify simplifying assumptions we have used. Then we estimate quantitatively upper and lower limits for inclusive pion DCX cross section in the energy region of $`T_0=1GeV`$ to 4$`GeV`$ to be tested in future experiments. Our working scheme for the calculation of $`\pi 2\pi `$ amplitude is the one pion exchange (OPE) model which we check by a comparison of theoretical cross section with experiment as well as distributions on two-pion effective mass for the process $`\pi ^{}p\pi ^+\pi ^{}n`$. We compare our results with predictions of other models and in particular of the meson exchange currents (MEC) model .
## 2 Gribov-Glauber picture of inclusive pion DCX
In this Section we shall discuss some aspects of the formalism for a calculation of the amplitude of the DCX process $`\pi ^{}pp\pi ^+nn`$ (Fig.3). Following Gribov it is convenient to introduce an integration over $`M_H^2`$, square of mass of the intermediate state in Fig.1. It is related to integration over Fermi motion of initial protons and momenta of final neutrons. These momenta are limited by experimental conditions. For example in the ITEP experiment on $`{}_{}{}^{6}Li`$ and $`{}_{}{}^{16}O`$ the inclusive DCX reaction $`A(\pi ^{},\pi ^+)X`$ was measured at kinetic energies $`T_0=0.6`$, $`0.75`$ and $`1.1GeV`$ ($`\theta 5^0`$) with kinematical condition $`\mathrm{\Delta }T=T_0Tm_\pi `$ ($`T`$ is the kinetic energy of the outgoing pion) in order to exclude additional pion production.
The kinematical region indicated above strongly limits momentum transfers $`t_{1(2)}`$ to neutrons in the transitions $`\pi ^{}pHn`$ and $`Hp\pi ^+n`$ (see Fig.3). $`M_H`$ can be very large at high incident energy, $`E_0`$, and increases as $`M_H^22m_NE_0`$. The value of $`M_H^2`$ varies also due to the Fermi motion. It is easy to obtain an estimate of this variation taking into account that at high energies, $`E_0m_N`$, $`|t_{1(2)}^{min}|`$ $`(M_H^2m_\pi ^2)^2/4E_0^2`$ $`\begin{array}{c}<\hfill \\ \hfill \end{array}p_{F}^{}{}_{}{}^{2}\begin{array}{c}<\hfill \\ \hfill \end{array}2m_N\mathrm{\Delta }T_{1(2)}`$ and $`(M_H^2)^{max}2`$$`E_0(2\mathrm{\Delta }T_im_N)^{1/2}`$. For energy $`E_01GeV`$, $`(M_H^2)^{max}`$ is equal to $`0.5GeV^2`$. As it was explained in Ref. , the integral over $`M_H^2`$ can be written either over the real axis or (due to analyticity in $`M_H^2`$) as a sum of the absorptive part and integral over semicircle at $`|M_H^2|`$ $``$ ($`M_H^2)^{max}`$. At high enough energies $`E_0`$ the last contribution can be neglected if the corresponding amplitude decreases faster than $`1/M_H^2`$ at large $`M_H^2`$. We gave arguments that the amplitude of Fig.3 should satisfy this property. In this case the DCX amplitude is proportional to the discontinuity in $`M_H^2`$ of the amplitude in Fig.3, and using unitarity we obtain the representation equivalent to the diagrams of Fig.1.
Contribution of a single particle intermediate state ($`\pi ^0,\eta ^0`$) or narrow resonance ($`\omega ,\mathrm{}`$) to the DCX cross section can be written in the form
$$\frac{d\sigma _{DCX}^{IR}}{d\mathrm{\Omega }}(\frac{d^2\sigma _{\pi H}}{dtdM^2}𝑑t𝑑M^2)^2,H=\pi ^0,\eta ^0,\omega ,\mathrm{}$$
(2)
Let us comment on contribution to DCX of different diagrams of Fig.1 that was taken in Ref. additively. The diagram ($`b`$) gives a small correction to the Glauber elastic rescattering of Fig.1($`a`$) due to the fact that $`\pi ^{}\eta `$ amplitude is relatively small (see Fig.2) while the production of two pions is a very important competing mechanism, especially for $`T_0\begin{array}{c}>\hfill \\ \hfill \end{array}0.6GeV`$. As for diagrams ($`a`$) and ($`c`$), if $`\pi 2\pi `$ amplitude is dominated by the pion exchange (see Fig.4), then their interference is absent. Really, the invariant form of the nucleon vertex in Fig.4 is $`\overline{u}\gamma _5u`$, while the amplitude of the transition $`\pi ^{}p\pi ^0n`$ is proportional to $`\overline{u}(A+B\widehat{q})u`$, where $`q=p_\pi ^{}+p_{\pi ^0}`$. So the interference term is Tr$`(\widehat{p}_1+m_N)\gamma _5(\widehat{p}_1^{}+m_N)(A+B\widehat{q})`$ = 0.
Note that for the $`s`$-wave $`\pi \pi `$ production which is a dominant inelastic process at energies $`E_0\begin{array}{c}<\hfill \\ \hfill \end{array}1GeV`$ the interference is absent by the same reason. For the $`s`$ wave (and in general for all even waves) the same expression (2) is valid. However for realistic case of $`\pi \pi `$ production, when both even and odd orbital momenta are important, the amplitude of $`\pi ^{}\pi ^+`$ transition (contrary to the diagonal $`\pi ^{}\pi ^{}`$ transition) can not be written in the form (2). This is especially clear in the $`\pi `$ exchange (OPE) model (Fig.4) where DCX amplitude is expressed in terms of the $`\pi ^{}\pi ^+\pi ^+\pi ^{}`$ amplitude (or backward elastic $`\pi ^{}\pi ^+`$ scattering amplitude).
In this paper we shall use the OPE model in order to study the problems outlined above: is it possible to neglect the contribution of the large semicircle at energies 1$`GeV`$ and what is the difference between integrals of forward and backward $`\pi ^{}\pi ^+`$ scattering amplitudes? In this way we shall find corrections $`\mathrm{\Gamma }_H`$ to simple expressions for DCX cross sections
$$\frac{d\sigma _{DCX}}{d\mathrm{\Omega }}=\frac{d\sigma _{DCX}^{\pi ^0}}{d\mathrm{\Omega }}\left(1+\underset{H}{}\mathrm{\Gamma }_H\right)$$
(3)
where
$$\mathrm{\Gamma }_{\pi ^+\pi ^{}}=\left(𝑑M\underset{t_{1\mathrm{min}}(M)}{\overset{t_{1\mathrm{max}}^{\mathrm{exp}}}{}}\frac{d^2\sigma _{\pi \pi ^+\pi ^{}}(M,t_1)}{dMdt_1}𝑑t_1/\underset{0}{\overset{t_{1\mathrm{max}}^{\mathrm{exp}}}{}}(d\sigma _{\pi ^0}/dt_1)𝑑t_1\right)^2,$$
(4)
$$\mathrm{\Gamma }_{\pi ^0\pi ^0}=\left(\sigma _{\pi \pi ^0\pi ^0}^{tot}/\sigma _{\pi \pi ^+\pi ^{}}^{tot}\right)^2\mathrm{\Gamma }_{\pi ^+\pi ^{}}$$
(5)
introduced in Ref. and shall obtain lower and upper bounds for these cross sections at higher energies.
## 3 Inelastic rescatterings and comparison with ITEP experiment
In this Section we shall calculate the contribution of IR to DCX using Eqs.(3)-(5) and experimental data on particle production in reactions
$$\pi ^{}pHn,H=\pi ^0,\pi ^+\pi ^{},\pi ^0\pi ^0,...$$
(6)
The quantities $`\mathrm{\Gamma }_H`$ in Eqs.(3)-(5) were determined using experimental data for the following hadronic states $`H`$: $`\pi ^0`$ , $`\pi ^+\pi ^{}`$ and $`\pi ^0\pi ^0`$ . While integrating over $`t`$ and $`M^2`$, experimental constraints on $`\mathrm{\Delta }T`$ are taken into account (see details in ): $`|t_{1(2)}|2m_N(\mathrm{\Delta }T^{max}/2)`$ \[<sup>1</sup><sup>1</sup>1In principle it is necessary to integrate over $`t_1`$ and $`t_2`$, but taking into consideration some uncertainty of the theoretical estimates as well as of the experimental information available we limited ourselves with this simple relation.\].
In the Table 1 we show the updated summary of energy dependence for forward inclusive DCX reaction cross sections of $`\pi ^{}`$ on $`{}_{}{}^{16}O`$ obtained in experiments at ITEP and calculated via Eqs.(3)-(5). The first line in the Table 1 contains the values of $`d\sigma _{DCX}^{\pi ^0}/d\mathrm{\Omega }`$ calculated in the framework of the elastic SSCX mechanism without taking into account effect of the Fermi motion (FM) on magnitude of the $`\pi N`$-interaction energy in the nucleus. In the values $`d\stackrel{~}{\sigma }_{DCX}^{\pi ^0}/d\mathrm{\Omega }`$ of the second line the FM effect <sup>2</sup><sup>2</sup>2We thank L.Alvarez-Ruso and M.J.Vicente Vacas for preparing the version of the code which takes into consideration the Fermi motion. is included. The resulting values are $`d\sigma _{DCX}/d\mathrm{\Omega }`$ = $`d\stackrel{~}{\sigma }_{DCX}^{\pi ^0}/d\mathrm{\Omega }`$ \+ ($`\mathrm{\Gamma }_{\pi ^+\pi ^{}}+\mathrm{\Gamma }_{\pi ^0\pi ^0})`$ $`d\sigma _{DCX}^{\pi ^0}/d\mathrm{\Omega }`$ are given at the fourth line. Due to a sharp change of DCX rate at energies of interest the FM effect gives smearing of the dip-bump structure, that in average leads in some increase of the values of DCX cross section. Note, that all values in the Table 1 correspond to the outgoing momentum spectra integrated over the region of $`\mathrm{\Delta }T`$ = 0 to 140 $`MeV`$.
## 4 Inclusive pion DCX in the energy range of 1 – 4 $`GeV`$
We see from the experimental data compilation of Fig.2 that the cross section values of the reaction (6) with $`H=\pi ^+\pi ^{}`$,
$`\pi ^{}p\pi ^+\pi ^{}n,`$ $`(6^{})`$
exceed ones with $`H\pi ^+\pi ^{}`$ not only at 1 $`GeV`$ but also at higher energies. So it is reasonable to assume that IR with two pion intermediate states will dominate inclusive pion DCX at least up to 4$`GeV`$. In this case the diagram of Fig.4 can be used to obtain a quantitative estimate of DCX rate within the OPE model. Note here that, besides of the well-known good description of experiment for (6) reaction, the OPE model is very convenient to our goals because of its factorization feature. This permits to separate $`M^2`$ and $`t`$ dependences of the (6) reaction amplitude, $`A_{\pi 2\pi }(M^2,t)`$, and to take into account limitations on their values in the experiment .
We start from the following formula
$$\frac{d\sigma _{DCX}^{\pi \pi }}{d\mathrm{\Omega }}\left(\underset{4m_\pi ^2}{\overset{M^2(t_{max})}{}}𝑑M^2\underset{t_{\mathrm{min}}(M^2)}{\overset{t_{\mathrm{max}}}{}}𝑑t\text{Im}A_{\pi 2\pi }^{forward}(M^2,t)\right)^2$$
(7)
equivalent to Eqs.(3)-(5). Then we use the OPE model to calculate a difference between imaginary parts of forward and backward amplitudes and to estimate the real part of the amplitude.
In OPE model it is possible to express the amplitude of the process (6) in terms of the $`\pi \pi `$ scattering amplitude,
$$A_{\pi 2\pi }(M^2,t)=\stackrel{~}{F}^2(t)A_{\pi \pi }(M^2)$$
(8)
where
$$\stackrel{~}{F}^2(t)=2G^2[|t|/(tm_\pi ^2)^2]\text{exp}[2R^2(tm_\pi ^2)],R^2=1.92GeV^2,$$
(9)
$`G`$ is the $`\pi NN`$-coupling constant.
$`(a)OPEmodelforDCX`$
As it was shown in the OPE model calculations are in a good agreement with the experimental data on the reaction (6) above 2 $`GeV`$. So for calculation of the cross section for this reaction we shall follow Ref., namely:
$$\frac{d^2\sigma _{\pi 2\pi }}{dM^2dt}=G^2\frac{Q(M^2,m_\pi ^2,m_\pi ^2)M}{2^4\pi ^2Q^2(s,m^2,m_\pi ^2)s}\stackrel{~}{F}^2(t)\sigma _{\pi \pi }^{tot}(M^2)$$
(10)
where $`G^2/4\pi =`$14.6, $`Q(s,m_1^2,m_2^2)`$ = $`\sqrt{s^22s(m_1^2+m_2^2)+(m_1^2m_2^2)}/2\sqrt{s}`$,
$`M`$ is the mass of two-pion state,
$$\sigma _{\pi ^+\pi ^{}}(M^2)=2\pi f_{\pi ^+\pi ^{}}^2𝑑z,z=cos\theta ,$$
(11)
and the $`\pi \pi `$ scattering amplitude is equal to
$$f_{\pi \pi }(M^2,z)=\frac{1}{Q(M^2)}\underset{l}{}(2l+1)[1+(1)^{I+l}]f_l^I(M^2)P_l(z).$$
(12)
By virtue of the unitarity condition we have:
$$\text{Im}A_{\pi \pi }(M^2,z=1)=2Q(M^2,m_\pi ^2,m_\pi ^2)M\sigma _{\pi \pi }^{tot}(M^2).$$
(13)
The partial-wave amplitudes $`f_l^I(M^2,z)`$ for orbital angular momentum $`l`$ and isospin $`I`$ are expressed via the phase shifts, $`\delta _l^I`$, and elasticities, $`\eta _l^I`$:
$$f_l^I(M^2)=\frac{1}{2i}[\eta _l^I(M^2)e^{2i\delta _l^I(M^2)}1],f_{\pi ^+\pi ^{}}=\frac{1}{6}f^2+\frac{1}{3}f^0+\frac{1}{2}f^1.$$
(14)
It should be noted that the equation (7) takes into account both $`\pi ^+\pi ^{}`$ and $`\pi ^0\pi ^0`$ intermediate states.
In the following we will use the complete set of amplitudes (phase shifts and elasticities) with $`l`$ = 0, 2 for $`I`$ = 0, 2 and with $`l`$ = 1, 3 for $`I`$ = 1 in the range of $`M`$ up to $``$1.8 $`GeV`$ from Ref. where fixed-$`t`$ and fixed-$`u`$ analyticity in conjunction with energy independent phase-shift analysis gave the single solution. The following features of the model confirm that its predictions are reasonable.
(1) The $`M^2`$ dependence of the total cross section of $`\pi ^+\pi ^{}`$scattering calculated according to Eq.(11), (13), (14) is in a good agreement with the data from Ref. (see Fig.5).
(2) The value of the parameter $`R^2`$ = 1.918 $`GeV^2`$ that we use was obtained on the base of our analysis of the $`t`$ dependence of the yield for the reaction (6) averaged over the energy region from 2 to 3 $`GeV`$ taken from . The total cross section of this reaction calculated via the integration of Eq.(10),
$$\sigma _{\pi ^{}p\pi ^+\pi ^{}n}=\underset{4m_\pi ^2}{\overset{M^2(t_{max})}{}}𝑑M^2\underset{t_{\mathrm{min}}(M^2)}{\overset{t_{\mathrm{max}}}{}}𝑑t\frac{d^2\sigma _{\pi ^{}p\pi ^+\pi ^{}n}}{dM^2dt}$$
(15)
(for $`t_{\mathrm{max}}`$ = $`t_{max}(s))`$ is compared to experimental data ($`fullstars`$) from Ref. in Fig.6. It is seen that the model used ($`fullcircles`$) reproduces experimental energy behaviour of $`\sigma _{\pi ^{}p\pi ^+\pi ^{}n}`$ only above $``$2.5 $`GeV`$. Some underestimate of the absolute value of the cross section decreases as $`T_0`$ increases.
In Fig.6 the cross section of the reaction (6), obtained with constraints of $`t_{max}=t_{max}^{exp}=0.135(GeV/c)^2`$ corresponding to the experimental limitation $`\mathrm{\Delta }T140MeV`$, is also presented. The experimental values ($`crosses`$ and $`emptystar`$) were calculated from the data of Ref.\[10a\] on $`M^2`$ and angular cross-section dependencies of the reaction (6) ($`crosses`$) and taken from Ref. for the interval $`1.5t/m_\pi ^28`$ ($`emptystar`$). The OPE model used ($`triangles`$ in Fig.6) agrees with experiment starting already from 1 $`GeV`$. Note that below $`T_0=1GeV`$ the OPE model predictions strongly deviate from experimental data. This is probably related to a substantial role of $`s`$-channel resonance production in this region.
(3) The experimental distributions on the two-pion mass ($`M`$ or $`M^2`$) for incident momentum values of 1.343 , 1.59 , 2.26 , and 4 $`GeV/c`$ are compared with the calculated ones (see Figs. 7-10). It is seen from Fig.7 that the model predicts too small cross section in the low $`M^2`$ region at the lowest energy. On the other hand, at higher energies the OPE model reproduces experiment rather well (including $`\rho `$ and $`f_2(1270)`$ resonances as it is seen in Figs. 9, 10). For our purpose it is essential to have a good description of experiment in $`t`$ interval close to $`t0.135(GeV/c)^2`$. It follows from Fig.8 that for the region of $`1.5t/m_\pi ^28`$ an agreement with experiment is better than for the total cross section of (6).
Thus we conclude that it is reasonable to use the OPE model for calculation of the high-energy DCX cross section.
$`b)Testsoftheassumptions`$
Now it is interesting to use the OPE model calculations for checking of the simplifying assumptions of our approach to DCX formulated in Ref. and discussed above.
The real and imaginary parts of forward ($`z`$ = 1) and backward ($`z=1`$) invariant amplitude, $`A_{\pi ^+\pi ^{}}(M^2,z)`$, were calculated using Eqs.(12) – (14). As can be seen from Fig.11, Re$`A_{\pi ^+\pi ^{}}(M^2,z=1)`$ oscillates but it is not small compared to Im$`A_{\pi ^+\pi ^{}}(M^2,z=1)`$ so its contribution should be taken into account at energies $`1GeV`$. This means that an integral over Re$`A_{\pi ^+\pi ^{}}(M^2,z=1)`$ can not be neglected in this energy range and thus the same is true for integral over the large semicircle. The backward parts of the amplitude appeared to be rather different from the forward ones within various intervals of $`M^2`$. The actual ranges of $`M^2`$ for a set of initial energies for experimental data avaliable on the reaction (6) are given in the Table 2.
In the Fig.12 we present the energy dependence of the cross section $`d\sigma _{DCX}^{\pi \pi }/d\mathrm{\Omega }`$ calculated according Eq.(7) at $`M_{max}^2|_{140}`$ for the following cases:
(1) Im$`A_{\pi 2\pi }^{forward}`$ is obtained through the Eq.(8) for the forward ($`z=1`$) amplitudes;
(2) Im$`A_{\pi 2\pi }^{forward}`$ is replaced by Im$`A_{\pi 2\pi }^{backward}`$ which is obtained according to Eq.(8) for the imaginary part of the backward ($`z=1`$) $`\pi \pi `$ amplitude;
(3) Same as in (2), but for Re$`A_{\pi 2\pi }^{backward}`$;
(4) both imaginary and real parts of the backward amplitude are taken into account.
The above analysis shows that for $`\pi ^+\pi ^{}`$-scattering amplitudes there are substantial differencies between integrals of forward and backward scattering amplitudes and they can be neglected only at rather low energies, where $`s`$-wave production dominates. We think however that OPE model overestimates this difference, especially for $`T_0\begin{array}{c}<\hfill \\ \hfill \end{array}1GeV`$, as it predicts too low $`s`$-wave production at these energies (see Fig.7). So the OPE-result for DCX can be considered as a lower bound for this cross section.
## 5 Results
Let us now calculate the cross section for the reaction
$$\pi ^{}+{}_{}{}^{16}O\pi ^++X$$
(16)
in the kinematical range $`0\mathrm{\Delta }T140MeV`$ of the experiment as a sum of the cross sections with $`\pi ^0`$ and 2$`\pi `$ in intermediate states (see Eq.(3) ),
$$\frac{d\sigma _{DCX}}{d\mathrm{\Omega }}_{140}=\frac{d\stackrel{~}{\sigma }_{DCX}^{\pi ^0}}{d\mathrm{\Omega }}_{140}+\frac{d\sigma _{DCX}^{\pi \pi }}{d\mathrm{\Omega }}_{140}.$$
(17)
Here $`d\stackrel{~}{\sigma }_{DCX}^{\pi ^0}/d\mathrm{\Omega }_{140}`$ is the cross section of the reaction (16) within the SSCX mechanism (FM effect is taken into account) which is shown in Fig.13 (see $`dashed`$ curve). $`d\sigma _{DCX}^{\pi \pi }/d\mathrm{\Omega }_{140}`$ is the contribution of the intermediate $`\pi \pi `$ state for the given kinematical range of $`\mathrm{\Delta }T`$. For Im$`A_{\pi 2\pi }^{forward}`$ according to Eqs.(3)-(5) and normalizing at $`T_o=0.6GeV`$ we obtain the $`dottedcurve`$ in Fig.13. The solid curve corresponds to an account of the full backward scattering amplitude. It can be considered as a lower bound for the DCX cross section, while the dotted curve can be considered as an upper bound for the cross section. Experimental point at 1.1 $`GeV`$ is close to this bound. This is related in our opinion to an underestimate of the $`s`$-wave contribution in the OPE model at these energies noted above. So at higher $`T_0`$ we expect future DCX measurements to be closer to the solid curve. Note, that both curves at $`T_o<1GeV`$, where OPE model strongly deviates from experimental data (see Fig.6), are calculated using not theoretical but experimental values (see the first two $`crosses`$ in Fig.6).
The contribution from three-pion intermediate state to Eq.(17) estimated using the cross section of $`\pi ^{}p\omega n`$ appeared to be not more than 10% even at the highest energy considered.
## 6 Comparison with other approaches
Let us discuss a relation between inelastic rescatterings considered in this paper, and other approaches which have been proposed for the DCX process.
Relation of the present approach to a model of meson exchanged currents (MEC) has been discussed in Ref.. The main difference is that in the OPE model we treat the $`\pi ^{}\pi ^+`$-scattering amplitude as a function of $`M^2`$ and integrate over this variable while in the MEC model the $`\pi ^{}\pi ^+`$-amplitude is approximated as a point-like interaction taken at the threshold ($`M=2m_\pi `$). Due to soft pion theorems this amplitude is small and thus leads to small modifications of SSCX predictions . We take into account in the OPE model both real and imaginary parts of $`\pi ^{}\pi ^+`$ amplitude in the regions of $`M^2`$ where they are not small, so it is not surprising that our results for DCX cross sections are substantially higher than SSCX predictions (even for a solid curve in Fig. 13).
Another possible mechanism of DCX is a production of an intermediate $`\mathrm{\Delta }`$ isobar on one proton of a nucleus, $`\pi ^{}p\pi ^+\mathrm{\Delta }^{}`$, with its absorption on a second proton, $`\mathrm{\Delta }^{}pnn`$. A characteristic feature of this mechanism is a fast decrease of DCX cross section at $`T_0>1GeV`$ due to energy dependence of the $`\sigma (\pi ^{}p\pi ^+\mathrm{\Delta }^{})`$ as is shown in Fig.2 by the curve.
Our study shows that at least up to energies $`4GeV`$ the region of $`M_H^2M_{max}^2`$ is very essential in integrals over $`M_H^2`$, and integrals of both imaginary and real parts of amplitudes are important. This means that average distances, $`d`$, between participating nucleons are relatively small: $`d(2m_N\mathrm{\Delta }T)^{1/2}`$. These nucleons are closely correlated and can be even considered as a 6-quark system (as, for example, in Ref.).
## 7 Experimental outlook and conclusions
We have extended the conventional Glauber mechanism for DCX process with an account of inelastic rescatterings. In this paper we investigated the problems which arise due to a non-diagonal nature of the process using the OPE model. This model gives a reasonable description of the reaction $`\pi ^{}p\pi ^{}\pi ^+n`$ at energies above $`2GeV`$. So we think that our predictions for DCX process based on this model, which are substantially higher than SSCX results, are reliable in this energy region. DCX is not yet studied experimentally at energies above $`T_0=1.1GeV`$. If strong deviations from our predictions will be found it will have important implications: either other mechanisms (e.g. those mentioned in the previous Section) are important or the phase-shift analysis of $`\pi \pi `$ scattering , which we have used, is not reliable in the large-mass region. An interesting information on these problems can be obtained by varying experimental interval of $`\mathrm{\Delta }T`$. Note that for small $`\mathrm{\Delta }T`$ (and in particular in the exclusive limit) an integration region in $`M^2`$ is small, and results should be closer to SSCX predictions.
Thus experimental study of the energy dependence of DCX in the region of $`T_0`$ =$`1÷5GeV`$ can give an important information on dynamics of this process. Cross section in this region according to our estimate is not too small, and experiments seems feasible. However, at higher energies the possibility to measure the cross section in the region of $`\mathrm{\Delta }T=0÷140MeV`$ will be limited by the energy resolution of the detector.
In conclusion, it is important to stress once more that the investigation of the pion DCX reactions in the intermediate energy region can provide a test on a validity of the Glauber–Gribov approach in the case when effects of inelastic rescatterings are large. Another class of processes where such inelastic contributions are very important is heavy ions interactions at very high energies. It is known (see, for example, ) that these contributions strongly modify a usual Glauber picture and an equilibration in the system can appear only as a result of such rescatterings effects.
## Acknowledgments
A.P.K. thanks I.S. and I.I.Tsukerman for permanent support. This work was supported in part by RFBR Grant No. 98-02-17179.
| $`T_0,GeV`$ | 0.6 | 0.75 | 1.12 |
| --- | --- | --- | --- |
| $`d\sigma _{DCX}^{\pi ^0}/d\mathrm{\Omega },\mu b/sr[8]`$ | 125.0 | 10.4 | 3.1 |
| $`d\stackrel{~}{\sigma }_{DCX}^{\pi ^0}/d\mathrm{\Omega },\mu b/sr`$ | 139.3 | 25.4 | 4.7 |
| $`\mathrm{\Gamma }_{\pi ^+\pi ^{}}+\mathrm{\Gamma }_{\pi ^0\pi ^0}`$ | 0.11 | 1.75 | 11.3 |
| $`d\sigma _{DCX}/d\mathrm{\Omega },\mu b/sr`$ | 153.1 | 43.6 | 39.7 |
| $`d\sigma _{DCX}^{\mathrm{exp}}/d\mathrm{\Omega },\mu b/sr[6]`$ | $`59.6\pm 7.4`$ | $`43.3\pm 5.5`$ | $`26.6\pm 8.9`$ |
Table 1.
| $`T_0,GeV`$ | $`M_{max}^2(s)`$ | $`M_{max}^2(s)|_{140}`$ |
| --- | --- | --- |
| 0.61 | 0.36 | 0.36 |
| 0.88 | 0.57 | 0.51 |
| 1.02 | 0.69 | 0.60 |
| 1.20 | 0.84 | 0.69 |
| 1.45 | 1.08 | 0.84 |
| 1.76 | 1.74 | 1.02 |
| 2.12 | 2.02 | 1.26 |
| 2.86 | 2.5 | 1.68 |
| 3.86 | 3.2 | 2.2 |
Table 2. |
warning/0001/hep-ph0001244.html | ar5iv | text | # Test of the QCD vacuum with the sources in higher representations
## Abstract
Recent accurate measurement of static potentials between sources in various $`SU(3)`$ representations provides a crucial test of the QCD vacuum and of different theoretical approaches to the confinement. In particular, the Casimir scaling of static potentials found for all measured distances implies a strong suppression of higher cumulants and a high accuracy of the Gaussian stochastic vacuum. Most popular models are in conflict with these measurements.
1. An accurate measurement of static potentials between sources in the eight different representations of $`SU(3)`$ group made recently in reveals a new quantitative picture of the QCD vacuum and provides a crucial test of existing theoretical models. Other measurements of static interaction are in general agreement with .
The most useful way to represent static potentials $`V_D(r)`$ in representations $`D=3,8,6,15a,10,27,24,15s`$ is through the complete set of field correlators in the framework of the Field Correlator Method (FCM) :
$$V_D(r)=\underset{T\mathrm{}}{lim}\frac{1}{T}\mathrm{ln}W(C),$$
(1)
where Wilson loop $`W(C)`$ for the rectangular contour $`C=r\times T`$ in the (34) plane has the cumulant expansion,
$$W(C)=Tr_Dexp_S\underset{n=2,4\mathrm{}}{}(ig)^nF(1)F(2)..F(n)d\sigma (1)\mathrm{}d\sigma (n)$$
(2)
Here $`F(k)d\sigma (k)=F_{34}(u^{(k)},x_0)d\sigma _{34}(^{(k)})`$ and the component $`F_{34}(u,x_0)E_3(u,x_0)=\varphi (x_0,u)E_3(u)\varphi (u,x_0)`$, where $`\varphi `$ is a parallel transporter and $`x_0`$ is an arbitrary point on the surface $`S`$ inside contour $`C`$; $`Tr_D\widehat{1}=1`$.
Dependence on $`D`$ enters in (2) through the generators $`T^a`$, since $`F(k)=F^a(k)T^a(a=1,\mathrm{}N_c^21),`$ and the main characteristics of $`D`$ is the quadratic Casimir operator $`C_D`$: $`T^aT^a=\widehat{1}C_D`$, so that the invariant square of the color charge in the representation $`D`$ is $`g^2C_D`$.
One can now express the connected correlators (cumulants) in (2) via $`C_D`$ and $`D`$–independent averages as follows (for more details see last reference in and),
$$Tr_DF(1)F(2)=C_D\frac{F^a(1)F^a(2)}{N_c^21},$$
(3)
$$Tr_DF(1)F(2)F(3)F(4)=\frac{C_D^2}{(N_c^21)^2}\{F^a(1)F^a(2)F^b(3)F^b(4)$$
$$+F^a(1)F^b(2)F^b(3)F^a(4)F^a(1)F^a(2)F^b(3)F^a(4)$$
$$+(1\frac{N_c}{2C_D})F^a(1)F^b(2)F^a(3)F^b(4)\}+O(\frac{1}{N_c^2}).$$
(4)
Note that the arguments of $`F(k)`$ in (4) and in (2) are ordered (e.g. clockwise, $`u^{(1)}<u^{(2)}<u^{(3)}<u^{(4)})`$ and therefore the only vacuum insertion is possible in the first term on the r.h.s. of (4) leading to the cancellation with the third term: hence the correlator (4) is a connected one vanishing at large distances, $`|u^{(1)}+u^{(2)}u^{(3)}u^{(4)}|\mathrm{}`$.
One can show in a similar way that the $`n`$-th cumulant in (2) contributes proportionally to $`C_D^n`$. As a result the static potential $`V_D(r)`$ has the expansion
$$V_D(r)=d_DV^{(2)}(r)+d_D^2V^{(4)}(r)+\mathrm{},$$
(5)
where in notations of ref. $`d_D=C_D/C_F`$ and $`C_F`$ is the fundamental Casimir operator, $`C_F=\frac{N_c^21}{2N_c}`$. The fundamental static potential contains perturbative Coulomb part $`V_{Coul}`$, confining linear and constant terms.
The Coulomb part ,which is also obtainable from the perturbative component of the FC in (3) is now known up to two loops and is proportional to $`C_D`$. Therefore one may expect quartic contributions proportional to $`C_D^2d_D^2`$, to the constant and linear terms, writing (5) as
$$V_D(r)=d_DV^{(2)}(r)+d_D^2(\overline{v}_0^{(4)}+\overline{\sigma }_4r).$$
(6)
Here $`\overline{v}_0,\overline{\sigma }_4`$ measure the contribution of the quartic cumulants to the constant term and string tension respectively.
Now the measurements of ref. allow to find $`\overline{v}_0,\overline{\sigma }_4`$ from all 8 sets of data. To this end one forms 7 combinations $`\zeta _DV_D(r)d_DV_F(r)=d_D(d_D1)(\overline{v}_0^{(4)}+\overline{\sigma }r)`$. As a typical example one can take fundamental and adjoint potentials,at distances between 0.05fm and 1.1fm from the data the $`\chi ^2`$ fit yields for $`\overline{v}_0,\overline{\sigma }_4`$
$$\overline{v}_0^{(4)}=(0.6\pm 0.67)10^3GeV$$
(7)
$$\overline{\sigma }^{(4)}=(1.136\pm 0.69)10^3GeV^2$$
(8)
The quality of the fit is reasonable, $`\chi ^2/N=0.45,N=43`$ One obtains similar figures also for D=6,15a,10(while 3 higher representations do not yield additional information)suggesting that $`\overline{\sigma }^{(4)}`$is negative while $`\overline{v}_0^{(4)}`$ is compatible
with zero,confirming in this way the parametrization (6). This analysis demonstrates the phenomenon of the Casimir scaling, i.e. proportionality of static potential $`V_D(r)`$ to the Casimir operator $`C_D`$ with accuracy better than one percent.
Physical consequences of the Casimir scaling are numerous and important.
First of all, the sign and magnitude of quartic correction (7),(8) can be understood in the FCM. Indeed the quartic term enters the potential $`V_D`$ with the factor $`(g^4)`$, as compared to $`+g^2`$ for the quadratic (Gaussian) term. Secondly, one can estimate $`E_3^2`$ term from the standard gluonic condensate as follows:
$$g^2E_3^aE_3^a\frac{4\pi ^2}{12}(0.04\pm 0.02)GeV^4(0.10\pm 0.06)GeV^4$$
(9)
and take into account that the cumulant expansion in (2) is actually in powers of the parameter
$$\xi g^2E_3^aE_3^aT_g^4$$
(10)
Here $`T_g`$ is the correlation length of the QCD vacuum; for bilocal correlator it was measured on the lattice , $`T_g^{(2)}1GeV^1`$. With the use of (10) one could expect that $`\overline{\sigma }_4`$ would be from 4 to 10% of the standard string tension, $`\sigma =0.2GeV^2`$ provided $`T_g=T_g^{(2)}`$. The value of $`\overline{\sigma }^{(4)}`$ calculated in (8) is at least 6 times smaller and suggests that quartic correlation length $`T_g^{(4)}`$ may be smaller than the Gaussian one, $`T_g^{(2)}0.2fm`$.
This result means that the Gaussian Stochastic Model (GSM), suggested in and successfully used heretofore in many applications , can be more accurate than it was even expected, at least in processes where string tension plays the most important role. On the other hand, the smallness of quartic and higher contributions implies a very specific picture of vacuum correlations.
Indeed the smallness of $`T_g^{(4)}`$ implies that color fields tend to form compact white bilocal combinations $`F^a(1)F^a(2)`$ which are almost noninteracting between themselves and therefore not contributing to the higher connected correlators. This looks like the picture of small white dipoles made of fields $`FF`$ (or of vector potentials $`A_\mu A_\mu `$ connected to $`FF`$ in the Fock–Schwinger or contour gauge).
One can also understand qualitatively the difference between $`T_g^{(2)}`$ and $`T_g^{(4)}`$, since $`T_g^{(2)}`$ measures correlation length between adjoint fields $`E_3^a(x)`$ and $`E_3^a(y)`$ in the Gaussian correlator $`E_3^a(x)E_3^a(y)`$, while $`T_g^{(4)}`$ refers to the correlation of two white complexes, and should be connected to the lowest glueball mass, $`M_G2GeV`$; hence $`T_g^{(4)}1/M_G0.1fm<T_g^{(2)}.`$
Finally, the Casimir scaling imposes severe restrictions on existing models of the QCD vacuum. For example, the center–symmetry flux model was tested and ruled out in , since in the original formulation it predicts vanishing adjoint string tension, while in the later modification – the fat vertex model it is still far from the accurate data .
Next one should mention models of the abelian projected vacuum which fail to provide Casimir scaling , at least in the simplest version .
Consider now the dilute instanton gas model and the $`SU(2)`$ group. The instantons may be present in the confining vacuum as an important source of chiral symmetry breaking. Then the Casimir scaling imposes a strict bound on the admixture of instantons in the QCD vacuum. Indeed, insertion in (2) and (4) of the instanton field strength $`gF_{34}^a(x,z)=\frac{4\delta _{a3}\rho ^2}{[(xz)^2+\rho ^2]^2}`$, $`(\rho `$ – instanton size, $`z`$ – its position, contribution of parallel transporters is neglected for simplicity since it gives for bilocal correlators a reduction of 20-30%, see for details) yields the following expression for the quartic contribution to the static potential ($`r<\rho `$) (in $`SU(2)`$ case)
$$V_D^{(4)}(r0)=\frac{N}{V}\frac{r^4}{\rho }\frac{3C_D^2C_D}{15}\frac{\pi ^6}{320}$$
(11)
where $`\frac{N}{V}`$ is the density of instanton in the vacuum. Inserting here the values of $`v_0^{(4)}`$ above(Eq(7)) one gets the following bounds on the density of instantons
$$\frac{N}{V}0.17fm^4,$$
(12)
which is much smaller than normal instanton density of $`1fm^4`$.
With such density the role of instantons in chiral symmetry breaking and other effects would be negligible.
A more stringent bound can be obtained from the quartic string tension generated by instantons . However the nonzero value of $`\overline{\sigma }_4`$ for instantons does not imply confinement. One should take into account that at large distances the sum of all partial string tensions $`_{n=2,4,\mathrm{}}\overline{\sigma }^{(n)}`$ for the dilute instanton gas vanishes ,.
In previous discussion we have ignored the fact that at large enough distances the adjoint charges are screened by the vacuum gluons, and the limiting value of the adjoint potential is equal to the doubled gluelump mass, $`2M_{gl}`$. This leads to an estimate of the screening distance $`r_0`$ from the relation; $`V_{adj}(r_0)=2M_{gl}`$ where $`M_{gl}`$ from is around 1.4 GeV and therefore $`r_01.4fm`$, which is beyond the distance where Casimir scaling was measured in .
Thus the Casimir scaling is a stringent test for all models considered and displays a strong suppression of quartic and higher connected correlators,hence supporting a good accuracy of the GSM. At this point one may wonder how the negligible small higher correlators are compatible with the screening of the adjoint potential at $`rr_0`$,which is not seen in the cumulant expansion (2).The solution to this problem was suggested in , where the screening terms have been identified as an addition to the Wilson loop,with the small coefficient proportional to $`N_c^2`$. The corresponding term actually comes from the two- and more Wilson loop averages,and therefore has a perimeter rather than area-law behaviour.Hence sreening cannot be seen in the one-Wilson-loop expansion (2), and the transition due to the properties of the definition (1) of the static potential occurs rather sharply in large $`T`$ limit.
This work was partially supported by the joint grant RFFI-DFG, 96-02-00088G. The author is grateful to Dr. G.Bali for sending his lattice data and for correspondence, to A.M.Badalian, D.I.Diakonov M.I.Polykarpov and G. ’t Hooft for useful discussions, and V.I.Shevchenko for discussions and help in calculations. |
warning/0001/cond-mat0001121.html | ar5iv | text | # Quantum Projector Method on Curved Manifolds
## I Introduction
The importance and difficulty of solving models of interacting quantum particles is hard to overstate. It is well known that the correlated motion of those particles gives rise to a wide variety of physical phenomena at different length and time scales, spanning disciplines like chemistry, condensed matter, nuclear, and high energy physics. Novel complex structures can emerge as a consequence of the competing multiple-length scales in the problem. Nonetheless, only a reduced set of interacting problems admits exact closed form solutions and the use of numerical techniques becomes essential if one is looking for accurate solutions not subjected to uncontrolled approximations. Among those techniques, the statistical methods offer the potential to study systems with large number of degrees of freedom, reducing the computational complexity from exponential to polynomial growth. This scaling behavior is particularly relevant when one recognizes that most of the interesting phenomena in many-body physics occurs in the thermodynamic limit . Unfortunately, for fermions (i.e. quantum particles obeying Fermi statistics) the sign problem plagues all useful stochastic algorithms and causes the variance of computed results to increase exponentially with increasing number of fermions .
On the other hand, the growing interest in physical systems whose state functions are defined on a general metric space makes the quantum mechanics of interacting particles in curved manifolds no longer a mere intellectual exercise, but one with very practical consequences. Perhaps the most well-known examples can be found in cosmology (e.g., matter in strong gravitational fields, atomic spectroscopy as probe of space-time curvature ), but the subject is certainly not exclusive to this field. In condensed matter a very elementary case is provided by a deformed crystal. Less well-known ones are mesoscopic graphitic microtubules and fullerenes. All these physical systems are ubiquitous in nature and the crucial role the curvature of the manifold plays has been confirmed by experimental observations (e.g. spectrum of collective excitations ). Therefore, the development of stable quantum methods with polynomial complexity in Riemannian manifolds represents a real challenge for many-body theorists.
The present manuscript deals with the (non-relativistic) many-particle Schrödinger equation in a general metric space and its solution using stochastic techniques. In particular, we will show how to construct approximate solutions (wave functions) for systems with broken time-reversal symmetry (e.g. electrons in the presence of external electromagnetic sources) avoiding the infamous “phase problem” . The main difficulty is to define a probability measure (semi-positive definite) which allows one to get the complex-valued state with no asymptotic signal-to-noise ratio decay in Euclidean time. This translates into a problem of geometric complexity, which is solved approximately using constraints in the manifold where the wave function has its support. In this way, we get stable but approximate solutions which can be systematically improved.
Among the large variety of problems one can attack, we decided to choose the general problem of fermions in the presence of external gauge fields to illustrate the main ideas. The effects of an external magnetic field on a system of electrons can be profound . The field couples to the electron’s charge and spin, modifying its spatial motion and lifting its spin degeneracies. The field can also create spatial anisotropy, effectively reducing the dimensionality of the system from three to two. The combination of the reduced dimension and the field itself is known to have novel consequences. For example, in a system of non-interacting electrons hopping on a square lattice, the field transforms the energy spectrum from the simplicity of trigonometric functions to the complexity of a field-dependent self-similar structure (Hofstadter’s butterfly) whose depth mathematicians are still fathoming . The combination of the reduced dimensionality, strong particle interactions and the field itself is known to have novel consequences, like the formation of isotropic fractional quantum Hall fluids , which are incompressible states of the two-dimensional homogeneous Coulomb gas.
The projector (zero temperature) method we will introduce uses random-walks to solve a general multidimensional partial differential equation second order in space coordinates and first order in time. Whenever mention is made of a random-walk we mean a Markov chain that is defined as a sequence $`_1,_2,\mathrm{},_K`$ of $`K`$ random variables that take values in configuration space, i.e. the space of particle positions. As usual, what characterizes a random-walk is its initial probability distribution and a conditional probability that dictates the transition from $`_i`$ to $`_{i+1}`$. This transition probability is non-unique and discretization dependent . Among all the possible choices we will require a prepoint discretization of the transition probability (short-time propagator) because we will use Monte Carlo methods to generate the walkers.
The paper is organized as follows. In Section II we present the formulation of the general problem of fermions in curved manifolds. In particular, for illustration purposes and to fix notation, we develop the formalism for spin-$`\frac{1}{2}`$ particles in the presence of an external electromagnetic potential. Then, we show how to project out the lowest energy state of a given symmetry in a manifold with curvature, and discuss the resulting Fokker-Planck equations for various distribution functions. Once the problem is precisely defined we develop, in Section III, path-integral solutions to those multidimensional differential equations, and give an interpretation of the emergent “quantum corrections” in the Euclidean action. The path-integral solutions are evaluated using Monte Carlo techniques in Section IV. There, we provide an stable step by step practical algorithm which emphasizes the subtle changes (with respect to the standard Diffusion Monte Carlo (DMC) technique) due to the metric of the manifold. In Section V we apply such computational implementation to the problem of an electron moving on the surface of a sphere in the presence of a Dirac monopole. Finally, Section VI summarizes the main findings and discusses the relevance of the stochastic method as applied to the physics of quantum Hall fluids.
## II Fermions on Riemannian manifolds
Notation. Consider a differentiable manifold $``$ of dimension $`d`$ (e.g., for the two-sphere S<sup>2</sup>, $`d`$=2) with coordinates $`𝐫_i=(x_i^1,\mathrm{},x_i^d)`$ defined on it. If $``$ is a Riemannian manifold, then it is a metric space, with metric tensor $`g^{\mu \nu }(𝐫_i)=g^{\mu \nu }(i)`$, such that the distance $`ds`$ between two points in $``$ is $`ds^2=g_{\mu \nu }(i)dx_i^\mu dx_i^\nu `$ in the usual way . The metric tensor is positive definite and symmetric $`g^{\mu \nu }=g^{\nu \mu }`$ (as we will see, this condition is important to define a probability density distribution), and is a function of the coordinates $`𝐫_i`$ with the property $`g_{\mu \gamma }g^{\gamma \nu }=\delta _\mu ^\nu `$. Let us consider the coordinate transformation $`h`$: $`x_i^\mu =h^\mu (x_i^1,\mathrm{},x_i^d)`$. Then, a generic second order contravariant ($`T^{\mu \nu }`$) and covariant tensor ($`T_{\mu \nu }`$) transform as
$$T^{\mu \nu }=\frac{x_i^\mu }{x_i^\alpha }\frac{x_i^\nu }{x_i^\beta }T^{\alpha \beta },T_{\mu \nu }=\frac{x_i^\alpha }{x_i^\mu }\frac{x_i^\beta }{x_i^\nu }T_{\alpha \beta }^{},$$
(1)
respectively. Throughout the paper Einstein’s summation convention for repeated indices is assumed ($`\mu ,\nu =1,\mathrm{},d`$).
Formulation of the problem. In this article we will be concerned with finite interacting fermion systems in the presence of an external electromagnetic potential $`a_\mu (𝐫_i)=a_\mu (i)=(𝐀(i),\varphi (i)=0)`$ ($`𝐁=𝐀`$ represents a uniform field, $`𝐀`$ and $`\varphi `$ are the vector and scalar potentials, respectively) whose Hamiltonian for motion on the manifold, in the coordinate representation, is given by
$$\widehat{\mathrm{I}\mathrm{H}}=\widehat{\mathrm{I}\mathrm{H}}_0+\widehat{V}(\{𝐫_i\},\{s_i\})$$
(2)
with
$$\widehat{\mathrm{I}\mathrm{H}}_0=D𝚫+i\frac{e\mathrm{}}{2m^{}c}\underset{i=1}{\overset{N}{}}\left[2a^\mu (i)_\mu +g^{1/2}(i)_\mu \left(g^{1/2}(i)a^\mu (i)\right)\right]+\frac{e^2}{2m^{}c^2}\underset{i=1}{\overset{N}{}}a^\mu (i)a_\mu (i),$$
(3)
$`_\mu =/x_i^\mu `$, and $`𝚫=_{i=1}^N\mathrm{\Delta }(i)`$, where
$$\mathrm{\Delta }=g^{1/2}_\mu \left(g^{\mu \nu }g^{1/2}_\nu \right)$$
(4)
is the covariant Laplace-Beltrami operator and $`\widehat{V}`$ is a potential energy operator. Notice that we use the conventional notation where the transformation between different forms of a given tensor is achieved by using the metric tensor (e.g., $`a^\mu =g^{\mu \nu }a_\nu `$, $`a_\mu =g_{\mu \nu }a^\nu `$), and $`g^{1/2}=\sqrt{\mathrm{det}g_{\mu \nu }}`$. This Hamiltonian characterizes the dynamics of $`N`$ non-relativistic indistinguishable particles of mass $`m^{}`$, charge $`e`$ and spin $`s_i=\frac{1}{2}`$ in a curved space with metric tensor $`g^{\mu \nu }`$, and $`D=\mathrm{}^2/2m^{}`$. We have assumed that the quantum Hamiltonian $`\widehat{\mathrm{I}\mathrm{H}}`$ in curved space has the same form as in flat space (this amounts to a particular operator ordering prescription.)
Given the previous ordering, one can rewrite the Hamiltonian above in terms of the generalized (hermitian) canonical momentum $`𝐩_\mu =i\mathrm{}(_\mu +\frac{1}{2}_\mu (\mathrm{ln}g^{1/2}))`$
$$\widehat{\mathrm{I}\mathrm{H}}=\frac{1}{2m^{}}\underset{i=1}{\overset{N}{}}g^{1/4}(i)𝚷_\mu g^{1/2}(i)g^{\mu \nu }(i)𝚷_\nu g^{1/4}(i)+\widehat{V}(\{𝐫_i\},\{s_i\}),$$
(5)
where the kinetic momentum $`𝚷_\mu =𝐩_\mu \frac{e}{c}a_\mu `$. The first term in Eq. 5 represents the kinetic energy of the system and is the non-relativistic approximation to the Dirac operator. $`\widehat{V}`$ includes the sum of one and two-body local interaction terms (and background potential in the case of a charge neutral system) and Zeeman contribution. The potentials are assumed to be finite almost everywhere and can only be singular at coincident points ($`𝐫_i=𝐫_j`$, $`ij`$)
We are interested in the stationary solutions of the resulting multidimensional Schrödinger equation
$$i\mathrm{}_t|\mathrm{\Psi }=\widehat{\mathrm{I}\mathrm{H}}|\mathrm{\Psi },$$
(6)
and will restrict ourselves to Hamiltonians $`\widehat{\mathrm{I}\mathrm{H}}`$ which are time-translation invariant. In the usual space-spin formalism the $`N`$-fermion states characterizing the system, $`X|\mathrm{\Psi }=\mathrm{\Psi }(X)`$, and all its first derivatives belong to the Hilbert space of antisymmetric (with respect to identical particle $`(𝐫_i,s_i)`$-exchanges) square-integrable functions $`_N=^2(^N)\text{ }\mathrm{C}^{2N}`$, defined as
$$_N=\{\mathrm{\Psi }\widehat{P}_{ij}\mathrm{\Psi }=\mathrm{\Psi },\mathrm{and}\mathrm{\Psi }=\sqrt{\mathrm{\Psi }|\mathrm{\Psi }}<\mathrm{}\},$$
(7)
where $`X=(,\Sigma )`$ ($`=(𝐫_1,\mathrm{},𝐫_N)`$ and $`\Sigma =(\sigma _1,\mathrm{},\sigma _N)`$ are discrete spin variables) and $`\widehat{P}_{ij}`$ represents the permutation of the pairs $`(𝐫_i,\sigma _i)`$ and $`(𝐫_j,\sigma _j)`$. $`^N`$ is the Cartesian product manifold of dimension $`dN`$.
Since the system Hamiltonian can be written as $`\widehat{\mathrm{I}\mathrm{H}}=\widehat{\mathrm{I}\mathrm{H}}_{}()+\widehat{\mathrm{I}\mathrm{H}}_S(\Sigma )`$, the last term representing the Zeeman coupling, the many-body wave function $`\mathrm{\Psi }(,\Sigma )`$ can be expressed as a tensor product of a coordinate and a spin function (or a linear combination of such products),
$$\mathrm{\Psi }(,\Sigma )=\mathrm{\Phi }()\mathrm{\Xi }(\Sigma ).$$
(8)
We want to construct $`N`$-fermion eigenstates of $`\widehat{\mathrm{I}\mathrm{H}}`$, $`\mathrm{\Psi }`$ , that are also eigenfunctions of the total spin $`S^2`$ ($`S=_{i=1}^Ns_i`$),
$$S^2\mathrm{\Psi }(X)=\mathrm{}^2s(s+1)\mathrm{\Psi }(X),$$
(9)
and this is always possible since $`[\widehat{\mathrm{I}\mathrm{H}},S^2]`$ = 0. Thus, the configuration part $`\mathrm{\Phi }()`$ must have the right symmetry in order to account for the Pauli principle. It turns out that a coordinate state $`\mathrm{\Phi }(𝐫_1,\mathrm{},𝐫_k,𝐫_{k+1},\mathrm{},𝐫_N)`$ which is symmetrized according to the Young scheme and has total spin $`s=\frac{N}{2}k`$ will be antisymmetric in the variables $`𝐫_1,\mathrm{},𝐫_k`$, and antisymmetric in the variables $`𝐫_{k+1},\mathrm{},𝐫_N`$. Moreover, $`\mathrm{\Phi }`$ possesses the property of Fock’s cyclic symmetry,
$$\left(1\mathrm{l}\underset{j=k+1}{\overset{N}{}}\widehat{P}_{kj}\right)\mathrm{\Phi }=\mathrm{\hspace{0.25em}0},$$
(10)
where, in this case, $`\widehat{P}_{kj}`$ refers to the transposition of particle coordinates $`𝐫_k`$ and $`𝐫_j`$. This last condition is a very useful one for testing the symmetry of a given coordinate function.
Quantum projection on curved manifolds. For a given total spin $`s`$ we are thus left with the task of solving the stationary many-body Schrödinger equation $`\widehat{\mathrm{I}\mathrm{H}}_{}\mathrm{\Phi }()=E\mathrm{\Phi }()`$, where $`\mathrm{\Phi }()=|\mathrm{\Phi }`$ satisfies the symmetry constraint discussed above. In particular, we are interested in the zero temperature properties of this quantum system, i.e. its ground state properties. To this end, we study the Euclidean time evolution of the state $`\mathrm{\Phi }`$, i.e. we analytically continue Eq. 6 to imaginary time (Wick rotation, $`tit\mathrm{}`$)
$$_t\mathrm{\Phi }=\left[\widehat{\mathrm{I}\mathrm{H}}_{}E_T\right]\mathrm{\Phi },$$
(11)
whose formal solution $`\mathrm{\Phi }(t)=\widehat{𝒰}(t)\mathrm{\Phi }_T=\mathrm{exp}[t(\widehat{\mathrm{I}\mathrm{H}}_{}E_T)]\mathrm{\Phi }_T`$ is used to determine the limiting distribution
$$\mathrm{\Phi }_0\underset{t\mathrm{}}{lim}\mathrm{\Phi }(t),$$
(12)
which is the largest eigenvalue solution of the evolution operator $`\widehat{𝒰}(t)`$ compatible with the condition $`\mathrm{\Phi }_0|\mathrm{\Phi }_T0`$, where $`\mathrm{\Phi }_T`$ is a parent state and $`E_T`$ is a suitable (constant) energy that shifts the zero of the spectrum of $`\widehat{\mathrm{I}\mathrm{H}}_{}`$.
We would like to solve the multidimensional differential equation Eq. 11 using initial value random walks. In this way, starting with an initial population of walkers (whose state space is $`^N`$) distributed according to $`p(,t=0)=\mathrm{\Phi }_T`$ ($`\mathrm{\Phi }_T`$ must be positive semi-definite), the ensemble is evolved by successive applications of the short (imaginary) time propagator $`\widehat{𝒰}(\tau )`$ ($`\tau =t/M`$, and $`M`$ is the number of time slices) to obtain the limiting distribution $`\mathrm{\Phi }_0`$. Then, we can introduce a “pseudo partition function”
$$𝒵=\mathrm{\Phi }_T|\widehat{𝒰}(t)\mathrm{\Phi }_T$$
(13)
in terms of which we can determine the ground state energy $`E_0`$ as
$$E_0E_T=\underset{t\mathrm{}}{lim}\frac{1}{t}\mathrm{ln}𝒵.$$
(14)
Similarly, other ground state expectation values, e.g., $`\mathrm{\Phi }_0|\widehat{𝒪}\mathrm{\Phi }_0`$, can be obtained as derivatives (with respect to a coupling constant $`J`$) of a modified pseudo partition function $`𝒵_J`$ whose evolution operator has a modified Hamiltonian, $`\widehat{\mathrm{I}\mathrm{H}}_{}+J\widehat{𝒪}`$.
In order to reduce statistical fluctuations in the measured quantities (i.e., observables) one can guide the random walk with an approximate wave function, $`\mathrm{\Phi }_G`$, which contains as much of the essential physics as possible (including cusp conditions at possible singularities of the potential $`\widehat{V}`$). Then, instead of sampling the wave function $`\mathrm{\Phi }(t)`$ one samples the distribution $`\stackrel{~}{f}(,t)=\mathrm{\Phi }(t)\mathrm{\Phi }_G`$ (properly normalized) with the initial time condition $`\stackrel{~}{f}(,t=0)=\mathrm{\Phi }_T\mathrm{\Phi }_G`$. Expectation values of operators $`\widehat{𝒪}`$ (observables) that commute with the Hamiltonian have a particularly simple form for guided walkers. For instance,
$$\underset{t\mathrm{}}{lim}\frac{\mathrm{\Phi }_T|\widehat{𝒪}\widehat{𝒰}(t)\mathrm{\Phi }_T}{\mathrm{\Phi }_T|\widehat{𝒰}(t)\mathrm{\Phi }_T}=\mathrm{\Phi }_G^1\widehat{𝒪}\mathrm{\Phi }_T_{\stackrel{~}{f}(t\mathrm{})},$$
(15)
where the average $`𝒜_{\stackrel{~}{f}}`$ stands for
$$𝒜_{\stackrel{~}{f}}=\frac{_^N𝝎\stackrel{~}{f}(,t\mathrm{})𝒜()}{_^N𝝎\stackrel{~}{f}(,t\mathrm{})},$$
(16)
$`\stackrel{~}{f}(,t\mathrm{})`$ is the long-time stationary probability of the system, and the (invariant) volume element $`𝝎`$ is given by the $`dN`$-form
$$𝝎=\left[\underset{i=1}{\overset{N}{}}g^{1/2}(i)\right]\mathrm{d}x_1^1\mathrm{}\mathrm{d}x_1^d\mathrm{}\mathrm{d}x_N^d.$$
(17)
Remember that in a general metric space the resolution of the identity operator with respect to the spectral family of the position operator is
$$1\mathrm{l}=_^N𝝎||.$$
(18)
It is important to stress that $`\mathrm{\Phi }_T`$ and the guiding function $`\mathrm{\Phi }_G`$ can, in principle, be different functions, although most of the practical calculations use the same function. It turns out that this importance sampling procedure is decisive to get sensible results when the potential $`\widehat{V}`$ presents some singularities.
Notice, however, that the quantum Hamiltonian $`\widehat{\mathrm{I}\mathrm{H}}`$ breaks explicitly time-reversal symmetry, meaning that in general $`\mathrm{\Phi }`$ will be a complex-valued function. Even if $`\mathrm{\Phi }`$ were real-valued, because it represents a fermion wave function it can acquire positive and negative values (the case where $`\mathrm{\Xi }(\Sigma )`$ is totally antisymmetric being the exception). Then, it is clear that we cannot in principle interpret $`\mathrm{\Phi }`$ or $`\stackrel{~}{f}`$ as a probability density.
For reasons that will become clear later we will be interested in sampling the probability density $`\overline{f}(,t)=|\stackrel{~}{f}(,t)|`$. The generalized diffusion equation in curved space for the importance-sampled function $`\overline{f}`$ can be derived directly from Eq. 11 with the result
$$_t\overline{f}=D\underset{i=1}{\overset{N}{}}\left[g^{1/2}(i)_\mu \left(g^{\mu \nu }(i)g^{1/2}(i)(_\nu \overline{f}\overline{f}F_\nu )\right)\right](E_LE_T)\overline{f},$$
(19)
where the drift velocity $`F_\nu ()=_\nu \mathrm{ln}\mathrm{\Phi }_G^2`$, and the “local energy” of the effective (“Fixed-Phase”) Hamiltonian $`\widehat{H}_{FP}`$ (see Eq. 50 and its derivation) is $`E_L()=\mathrm{\Phi }_G^1\widehat{H}_{FP}\mathrm{\Phi }_G`$ with
$$\widehat{H}_{FP}=D𝚫+D\underset{i=1}{\overset{N}{}}\left[(^\mu \chi ()\frac{e}{\mathrm{}c}a^\mu (i))(_\mu \chi ()\frac{e}{\mathrm{}c}a_\mu (i))\right]+\widehat{V}(),$$
(20)
where $`\chi ()`$ is the phase of the many-body state $`\mathrm{\Phi }`$, i.e. $`\mathrm{\Phi }=|\mathrm{\Phi }|\mathrm{exp}[i\chi ]`$. The differential equation satisfied by the distribution function $`\overline{f}`$ is formally equivalent to the one describing Brownian motion on a general manifold (including generation and recombination processes), and corresponds to a Kramers-Moyal expansion with exactly two terms. In fact, we can rewrite the equation above as a Fokker-Planck equation for $`dN`$ continuous stochastic variables $`\{𝐫_i\}_{i=1,\mathrm{},N}`$
$$_t\overline{f}=\left\{\overline{}_{\mathrm{FP}}(\overline{E}_LE_T)\right\}\overline{f},$$
(21)
where the (time-independent) Fokker-Planck operator $`\overline{}_{\mathrm{FP}}`$ is given by
$$\overline{}_{\mathrm{FP}}=\underset{i=1}{\overset{N}{}}[_\mu _\nu (\overline{D}^{\mu \nu }(i))_\mu (\overline{D}^\mu ())].$$
(22)
The diffusion matrix (contravariant tensor) $`\overline{D}^{\mu \nu }`$ and drift $`\overline{D}^\mu `$ (which does not transform as a contravariant vector) are given by
$`\overline{D}^{\mu \nu }`$ $`=`$ $`Dg^{\mu \nu }`$ (23)
$`\overline{D}^\mu `$ $`=`$ $`\overline{D}^{\mu \nu }F_\nu +_\nu \overline{D}^{\mu \nu }\overline{D}^{\mu \nu }\mathrm{\Gamma }_{\nu \sigma }^\sigma ,`$ (24)
where $`\mathrm{\Gamma }_{\mu \nu }^\sigma `$ is the Christoffel symbol of the second kind
$`\mathrm{\Gamma }_{\mu \nu }^\sigma `$ $`=`$ $`{\displaystyle \frac{1}{2}}g^{\sigma \rho }\left(_\mu g_{\nu \rho }+_\nu g_{\mu \rho }_\rho g_{\mu \nu }\right)`$ (25)
$`\mathrm{\Gamma }_{\nu \sigma }^\sigma `$ $`=`$ $`{\displaystyle \frac{1}{2}}_\nu \mathrm{ln}g,`$ (26)
and the modified local energy
$$\overline{E}_L=E_L+\overline{D}^{\mu \nu }\mathrm{\Gamma }_{\mu \sigma }^\sigma F_\nu +_\mu \left(\overline{D}^{\mu \nu }\mathrm{\Gamma }_{\nu \sigma }^\sigma \right).$$
(27)
Notice, however, that singularities in the “quantum corrections” to the local energy $`E_L`$ due to the metric, can induce very large fluctuations in $`\overline{E}_L`$. Moreover, the probability density $`\overline{f}`$ does not transform as a scalar function ($`\overline{f}(,t)\overline{𝝎}=\overline{f}(^{},t)\overline{𝝎}^{}`$, where the primes represent the transformed coordinates and $`\overline{𝝎}=\mathrm{d}x_1^1\mathrm{}\mathrm{d}x_1^d\mathrm{}\mathrm{d}x_N^d`$ is a volume element in $`^N`$). Therefore, it is more convenient to work with a probability density that is a scalar and which is defined as
$$f(,t)=\left[\underset{i=1}{\overset{N}{}}g^{1/2}(i)\right]\overline{f}(,t).$$
(28)
The differential equation $`f`$ satisfies is of the form Eq. 21 with bar quantities replaced by unbar ones (e.g. $`\overline{}_{\mathrm{FP}}_{\mathrm{FP}}`$). It turns out that $`D^{\mu \nu }=\overline{D}^{\mu \nu }`$ and the drift (which is not a tensor)
$$D^\mu =D^{\mu \nu }F_\nu +_\nu D^{\mu \nu }+D^{\mu \nu }\mathrm{\Gamma }_{\nu \sigma }^\sigma .$$
(29)
Note that in this case the quantum correction to the local energy vanishes. Furthermore, if the metric is diagonal, i.e. $`g_{\mu \nu }=g^{1/2}\delta _{\mu \nu }`$, then the correction to the flat space drift also vanishes, i.e. $`_\nu D^{\mu \nu }+D^{\mu \nu }\mathrm{\Gamma }_{\nu \sigma }^\sigma =0`$, and $`D^\mu =Dg^{1/2}F^\mu `$. This last remark is quite important, specially for $`d=2`$ where it is always possible to choose a coordinate system ($`𝐫_i=(\xi _i^1,\xi _i^2)`$) where the metric tensor is diagonal (conformal gauge ), and use the conformal parameterization ($`z_i=\xi _i^1+i\xi _i^2,\overline{z}_i=\xi _i^1i\xi _i^2`$) which greatly simplifies the resulting expressions (see Section V).
## III Path Integral Solutions
The generalized Fokker-Planck Eq. 21 describes the time evolution of a distribution function $`f`$ which is completely determined by the distribution function at $`t=t_0=0`$. In this sense it describes a continuous stochastic process that is Markovian. Because it represents a Markov process, the conditional probability that if the system is in $``$ at time $`t=0`$ it will jump to $`^{}`$ in time $`t`$ (importance-sampled Green’s function) $`G(^{};t)`$ contains the complete information about the process, and it follows that the probability densities $`f(,t+\tau )`$ and $`f(,t)`$ are connected by
$$f(^{},t+\tau )=_^N𝝎G(^{};\tau )f(,t),$$
(30)
where the Green’s function $`G(^{};\tau )`$ is a transition probability for moving particles from $``$ to $`^{}`$ in time $`\tau `$ with the initial value $`G(^{};0)=\left[_{i=1}^Ng^{1/2}(i)\right]\delta (^{})`$, and is formally given by
$$G(^{};\tau )=\stackrel{~}{\mathrm{\Phi }}_G(^{})^{}|\mathrm{exp}[\tau (\widehat{H}_{FP}E_T)]|\stackrel{~}{\mathrm{\Phi }}_G^1(),$$
(31)
with
$$\stackrel{~}{\mathrm{\Phi }}_G()=\left[\underset{i=1}{\overset{N}{}}g^{1/2}(𝐫_i)\right]\mathrm{\Phi }_G().$$
(32)
Path integral solutions for $`f`$ may be derived from the transition probability density for small $`\tau `$. Iteration of the Chapman-Kolmogorov equation for $`G`$ allows one to express the evolution of $`f(^{},t)`$ from the initial distribution $`f(,t=0)`$ in terms of the short-time Green’s function as
$$f(^{},t)=_^N𝝎_{M1}\mathrm{}_^N𝝎_0G(_{M1}_M;\tau )\mathrm{}G(_0_1;\tau )f(_0,0),$$
(33)
where $`t=M\tau `$; we identify $`_0=`$ and $`_M=^{}`$. By simple inspection the solution of the generalized Fokker-Planck equation $`f`$ stays positive if it was initially positive (i.e., if $`f(,0)>0`$.)
It is clear then that the functional integral representation of $`f`$ requires knowledge of the infinitesimal evolution operator. It is well-known in a similar context that the integrand of the functional integral is not unique, it is discretization dependent (compatible with the Markovian property of the paths since $`(t)`$ is only sampled at $`t\tau `$ and $`t`$.) On the other hand, it is crucial for numerically simulating those paths to use a discretization where the drift velocity and diffusion are evaluated at the prepoint in the integral equation. Following Feynman we have determined the functional form of $`G(^{};\tau )`$ to $`𝒪(\tau ^2)`$. We are going to present the final result and omit the details of the calculation which are just simple (although lengthy) manipulations of Taylor series expansions and gaussian integration. Thus, to $`𝒪(\tau ^2)`$ the short-time conditional probability is given by
$`G(^{};\tau )=G_b(^{};\tau ){\displaystyle \underset{i=1}{\overset{N}{}}}G_i^0(^{};\tau )`$ (34)
where
$`G_b(^{};\tau )=\mathrm{exp}\left[\tau \left({\displaystyle \frac{[E_L()+E_L(^{})]}{2}}E_T\right)\right],`$ (35)
and
$`G_i^0(^{};\tau )=\left({\displaystyle \frac{1}{4\pi D\tau }}\right)^{d/2}\mathrm{exp}\left[{\displaystyle \frac{\left(x_i^\mu x_i^\mu \tau D^\mu ()\right)g_{\mu \nu }(𝐫_i)\left(x_i^\nu x_i^\nu \tau D^\nu ()\right)}{4D\tau }}\right],`$ (36)
that is, a gaussian distribution with variance matrix $`2D^{\mu \nu }`$ and mean $`x_i^\mu +\tau D^\mu ()`$. Sample trajectories (continuous but nowhere differentiable) are generated by using the Langevin equation associated with the process, i.e. $`x_i^\mu =x_i^\mu +\tau D^\mu ()+\sqrt{\tau }\eta `$ , where $`\eta `$ is a gaussian random variable with zero mean. Note that $`D^\mu `$ and $`D^{\mu \nu }`$ are evaluated at the prepoint in the integral equation.
Therefore, the Wick-rotated path integral for $`G(^{};t)`$ is
$$G(^{};t)=_{(0)=}^{(t)=^{}}𝒟[𝝎(t)]\mathrm{exp}[S\left[(t)\right]],$$
(37)
where the measure $`𝒟[𝝎(t)]=lim_M\mathrm{}(4\pi D\tau )^{Md/2}𝝎_1\mathrm{}𝝎_{M1}`$, and the Euclidean action
$$S\left[(t)\right]=_0^t𝑑t^{}\left\{\frac{1}{4D}\underset{i=1}{\overset{N}{}}(\dot{x}_i^\mu D^\mu [(t^{})])g_{\mu \nu }(i)(\dot{x}_i^\nu D^\nu [(t^{})])+E_L[(t^{})]E_T\right\}.$$
(38)
The integrand above represents a generalized Onsager-Machlup function , and the dot is a short-hand for time derivatives. In closing this Section, we would like to mention that functional integral solutions for $`\overline{f}`$ can be obtained from previous expressions after making the replacement $`(D^\mu ,E_L)(\overline{D}^\mu ,\overline{E}_L)`$.
Interpretation of the quantum corrections. The general methodology we have developed so far, can be equally applied to other situations which do not necessarily involve a curved manifold such as, for instance, particles moving in a medium with position dependent diffusion constant. In order to adapt our previous formalism, we need to understand qualitatively the origin of the quantum corrections to the short-time propagator obtained above. To this end, we will illustrate the general idea with the following 1$`d`$ equation ($`=\mathrm{I}\mathrm{R}`$, $`N=1`$)
$$_x^2(D(x)f(x,t))=_tf(x,t).$$
(39)
The “standard” approach to finding the Green’s function for this problem is simply to solve the equation $`_x^2(D(x)G(x,t))=_tG(x,t)`$ subject to the boundary condition $`G(x,0)=\delta (x)`$. We can do this by taking the Fourier transform in $`x`$:
$$\frac{k^2}{2\pi }_{}𝑑k^{}D(kk^{})G(k^{},t)=_tG(k,t),$$
(40)
with the boundary condition $`G(k,0)=1`$. This is more complex than the usual diffusion equation because we get a convolution of $`D(k)`$ and $`G(k,t)`$. However, we can find an approximate solution for $`G(k,t)`$, valid for small time, by noting that the boundary condition implies that, for small times $`t`$, $`G(k,t)G(k^{},t)𝒪(t)`$ and so we have
$$\frac{k^2}{2\pi }_{}𝑑k^{}D(kk^{})G(k,t)+𝒪(t)=_tG(k,t)$$
(41)
or, simply, $`k^2D(x=0)G(k,t)+𝒪(t)=_tG(k,t)`$, where $`D(x=0)`$ is $`D(x)`$ evaluated at the prepoint. For small times we can ignore the order $`t`$ term and solve for $`G`$ with the result $`G(k,t)=\mathrm{exp}[tk^2D(0)]`$. Taking the inverse Fourier transform we finally have
$$G(x,t)=\frac{1}{\sqrt{4\pi D(0)t}}\mathrm{exp}[x^2/4D(0)t]+𝒪(t^2),$$
(42)
which is just the plain Green’s function for a 1$`d`$ random-walk with $`D(x)`$ evaluated at the prepoint and with no quantum corrections. This is, in fact, the result that the Green’s function for the Jacobian times $`\overline{f}`$ has no quantum corrections.
To make things clear let us look at a different equation with the same boundary conditions
$$D(x)_x^2G(x,t)=_tG(x,t),$$
(43)
which characterizes a free Brownian particle in 1$`d`$. Again, taking the Fourier transform we get
$$\frac{1}{2\pi }_{}𝑑k^{}k_{}^{}{}_{}{}^{2}D(kk^{})G(k^{},t)=_tG(k,t).$$
(44)
Making the same approximation as above, replacing $`G(k^{},t)`$ with $`G(k,t)`$ in the integrand, making an error of $`𝒪(t)`$, and noticing that
$$\frac{1}{2\pi }_{}𝑑k^{}(kk^{})^2D(k^{})=k^2D(0)+2ki_xD(0)_x^2D(0)$$
(45)
we get $`G(k,t)=\mathrm{exp}[t(k^2D(0)+2ki_xD(0)_x^2D(0))]`$. When we take the inverse Fourier transform the terms with derivatives of $`D(x)`$ at $`x=0`$ give precisely the quantum corrections for this simple case
$$G(x,t)=\frac{1}{\sqrt{4\pi D(0)t}}\mathrm{exp}[(x2t_xD(0))^2/4D(0)t+t_x^2D(0)]+𝒪(t^2).$$
(46)
This is the result that the Green’s function for $`\overline{f}`$ has quantum corrections.
The general case is just as simple, and the following rule emerges. Given any second order differential equation (first order in $`t`$), no matter how many dimensions, with or without curvature, with or without a position dependent diffusion constant, the rule for obtaining the short-time Green’s function with everything evaluated at the prepoint is as follows: Bring all derivatives in each term all the way to the left of that term. Once this is done one can simply write down the Green’s function for a generalized diffusion process assuming the $`D(x)`$, $`D^\mu `$, $`D^{\mu \nu }`$, whatever position dependent terms they may be, are constant and evaluated at the prepoint. The quantum corrections are then seen to be simply those extra terms we get when commuting the derivatives to the left.
Fermion-phase problem: Fixed-Phase method. It is evident that one cannot make a probability density out of a complex and/or antisymmetric wave function. This is the reason why we decided to write down Fokker-Planck equations for the distribution $`f`$ (or $`\overline{f}`$) and not $`\stackrel{~}{f}`$. Nevertheless, the phase factor associated with the original complex distribution must show up in the evaluation of the expectation values. It is well-known that this causes the variance of the computed results to increase exponentially with increasing number of degrees of freedom. This problem is known as the fermion-phase catastrophe, and in this Section we will review a method to obtain stable, albeit approximate, path-integral solutions whose stochastic determination has a polynomial, instead of exponential, complexity. The generalization of the ideas presented below to bosons (with complex-valued states) or anyons in general is straightforward .
In order to avoid cumbersome notation which would obscure the main ideas, here we will consider a simplified version of our original Hamiltonian
$$\widehat{\mathrm{I}\mathrm{H}}_{}=\frac{𝚷𝚷}{2m^{}}+\widehat{V}(),$$
(47)
where, for simplicity, we introduced the vector notation, $`𝚷=(𝚷(1),\mathrm{},𝚷(N))`$, and in this subsection the same convention will be used for other bold quantities. The microscopic equations governing the imaginary-time evolution of our interacting system can be found from a variational principle of the form $`\delta S[(t)]=0`$, where the Euclidean action is given by
$$S[(t)]=_0^t𝑑t^{}_^N𝝎\left\{\frac{1}{2m^{}}\left[𝚷\mathrm{\Phi }\right]^{}\left[𝚷\mathrm{\Phi }\right]+\widehat{V}\mathrm{\Phi }^{}\mathrm{\Phi }+\mathrm{\Phi }^{}_t^{}\mathrm{\Phi }\right\}.$$
(48)
Finding the states $`\mathrm{\Phi }`$ and $`\mathrm{\Phi }^{}`$ which minimize $`S[(t)]`$ is equivalent to solving the Schrödinger equation and its complex conjugate. Equivalently, we can consider as independent real fields the phase and modulus of the wave function, i.e. $`|\mathrm{\Phi }|`$ and $`\chi `$ such that $`\mathrm{\Phi }=|\mathrm{\Phi }|\mathrm{exp}[i\chi ]`$, and perform independent functional variations on $`S[(t)]`$. The resulting Euler-Lagrange equations are:
$$\{\begin{array}{ccccc}\hfill _t|\mathrm{\Phi }|& =& \text{ }\mathrm{}e\left\{\mathrm{exp}\left[i\chi \right]\widehat{\mathrm{I}\mathrm{H}}_{}\mathrm{\Phi }\right\}\hfill & =& \widehat{H}_{FP}|\mathrm{\Phi }|\hfill \\ & & & & \\ \hfill (_t\chi )|\mathrm{\Phi }|^2& =& \text{ }\mathrm{}m\left\{\mathrm{exp}\left[i\chi \right]\widehat{\mathrm{I}\mathrm{H}}_{}\mathrm{\Phi }\right\}|\mathrm{\Phi }|\hfill & =& \frac{\mathrm{}}{2}\mathbf{}_𝝁𝐉^\mu \hfill \end{array},$$
(49)
where $`\mathrm{}e`$ and $`\mathrm{}m`$ stand for the real and imaginary parts of the expressions in brackets, respectively,
$$\widehat{H}_{FP}=\frac{𝐩𝐩}{2m^{}}+\stackrel{~}{V}(),\stackrel{~}{V}=\widehat{V}+D(\mathbf{}^\mu \chi ()\frac{e}{\mathrm{}c}𝒂^\mu )(\mathbf{}_\mu \chi ()\frac{e}{\mathrm{}c}𝒂_\mu ),$$
(50)
and
$$𝐉^\mu =\frac{1}{2m^{}}\left(\mathrm{\Phi }^{}\left[𝚷^\mu \mathrm{\Phi }\right]+\left[𝚷^\mu \mathrm{\Phi }\right]^{}\mathrm{\Phi }\right)=\frac{|\mathrm{\Phi }|^2}{m^{}}\left(\mathrm{}\mathbf{}_\mu \chi ()\frac{e}{c}𝒂_\mu \right)$$
(51)
is the probability current. The singularities of $`\chi ()`$ occur at the zeros of $`\mathrm{\Phi }`$, which generically have codimension two.
So far we have simply mapped the original fermion problem into a bosonic one for $`|\mathrm{\Phi }|`$ but still coupled to its phase fluctuations. Alternatively, one can regard this as a gauge transformation of the original fermion problem, whose effect is to add a non-local gauge field potential, $`\mathbf{}^\mu \chi `$, giving rise to a fictitious magnetic field. Notice that it is this gauge field that contains information on particle statistics. Moreover, although the geometry of $`\chi ()`$ can be altered by a gauge transformation, the singularities remain invariant.
The Fixed-Phase (FP) method consists in making a choice for the phase, $`\chi _T`$, and solving the bosonic problem for $`|\mathrm{\Phi }|`$ exactly using stochastic methods. The method is stable and has the property of providing a variational upper bound to the exact ground energy $`E_0`$, $`E_{FP}E_0`$ (the equality holds when $`\chi _T`$ is the exact ground state phase), and for a given $`\chi _T`$ the lowest energy consistent with this phase. The trial phases $`\chi _T`$ should conserve the symmetries of $`\widehat{\mathrm{I}\mathrm{H}}_{}`$ and particle statistics (for time-reversal invariant systems there is a way for systematically improving a given mean-field phase using projection techniques ). Notice that the FP method projects out the lowest energy state of a given symmetry. Therefore, the method allows one to compute also excitations which are “ground states” of a particular symmetry. For ground state properties of real symmetric Hamiltonian operators the FP approach reduces to the Fixed-Node method .
## IV Computational Implementation
In this Section we present an algorithm for computing the ground state properties of quantum many-body systems defined on a curved manifold with general metric $`g^{\mu \nu }`$. As mentioned in the Introduction, this can be accomplished by performing all multidimensional integrals using Monte Carlo techniques. In this way, ground state expectation values are obtained by averaging over a large number of particle configurations generated according to a certain limiting probability distribution $`p(,t\mathrm{})`$. There is some freedom in the choice of this distribution $`p(,t)`$, however, to reduce statistical fluctuations in the observables to be computed it is more efficient to use the so-called importance-sampled distribution $`\overline{f}(,t)`$, which is the product of the absolute value of the solution of the time dependent Schrödinger equation $`\mathrm{\Phi }(,t)`$ and some positive function $`\mathrm{\Phi }_G()`$ that is the best available approximation to the modulus of the ground state eigenfunction. In a curved manifold, on the other hand, it is more convenient to work with the modified importance-sampled distribution $`f(,t)`$, which is defined as a product of the conventional importance-sampled distribution $`\overline{f}(,t)`$ and the metric (see Eq. 28).
The propagation of particle configurations in time $`\tau `$ is determined by the conditional probability (Green’s function) $`G(^{};\tau )`$, whose separation into a diffusion (plus drift) and branching parts (see Eq. 34) makes it very simple to simulate numerically. The gaussian term represents propagation according to the equation $`x_i^\mu =x_i^\mu +\tau D^\mu ()+\sqrt{\tau }\eta `$ , where $`\eta `$ is a gaussian random variable with zero mean. The effect of the term $`\tau D^\mu ()`$ is to superimpose a drift velocity on the random diffusion process so that particle configurations are directed towards regions of configuration space where $`\mathrm{\Phi }_G()`$ is large. The branching term $`G_b(^{};\tau )`$ in Eq. 35, determines the creation and annihilation of configurations (walkers) at the point $`^{}`$ after a move. If the size of the ensemble of walkers at any time $`t`$ is defined as
$$𝒫(t)=_^N\overline{𝝎}f(,t)$$
(52)
then, its rate of change is given by
$$_t𝒫(t)=_^N\overline{𝝎}\left[E_L()E_T\right]f(,t).$$
(53)
Therefore, if the local energy $`E_L()`$ is a smooth function of $``$, and the trial energy $`E_T`$ is suitably adjusted, the size of the ensemble of walkers will remain approximately constant as the configurations propagate. In particular, if the local energy is constant and equal to $`E_T`$ then the fluctuations in the ensemble size will vanish. To ease notation, in the rest of the paper we will only consider the standard situation $`\mathrm{\Phi }_T=\mathrm{\Phi }_G`$. In such a case, ground state expectation values of a generic observable $`\widehat{𝒪}`$ will be computed as
$$\underset{t\mathrm{}}{lim}\frac{\mathrm{\Phi }_G|\widehat{𝒪}\widehat{𝒰}(t)\mathrm{\Phi }_G}{\mathrm{\Phi }_G|\widehat{𝒰}(t)\mathrm{\Phi }_G}=\mathrm{\Phi }_G^1\widehat{𝒪}\mathrm{\Phi }_G_{f(t\mathrm{})}=_^N\overline{𝝎}\frac{f(,t\mathrm{})}{𝒫(t\mathrm{})}[\mathrm{\Phi }_G^1\widehat{𝒪}\mathrm{\Phi }_G]().$$
(54)
In the following we present a step by step computational algorithm for implementing the stochastic approach discussed above. Most parts of the algorithm follow closely the standard DMC method, described for instance in , but we believe that it is useful to present these steps in detail here because a number of straightforward but subtle modifications due to the space curvature are involved. To be more specific (without loosing the general features), let us present the algorithm as it is applied to fermionic systems within the FP approach.
Algorithm.
$`\mathrm{\S }`$ 1. Construct a guiding wave function $`\mathrm{\Phi }_G`$, which is the best available approximation to the exact ground state (or lowest energy state of a given symmetry). In principle any choice of $`\mathrm{\Phi }_G`$ which has finite overlap with the exact state is acceptable, but the more accurate $`\mathrm{\Phi }_G`$ is the faster the convergence to the stationary solution will be, and the lower the statistical fluctuations will be as well. Recall that zero variance is obtained when the guiding wave function is equal to the desired ground state (and the ground state is bosonic).
$`\mathrm{\S }`$ 2. Given the guiding function compute the quantum drift velocity $`F_\nu =_\nu \mathrm{ln}\mathrm{\Phi }_G^2`$ and the local energy $`E_L=\mathrm{\Phi }_G^1\widehat{H}_{FP}\mathrm{\Phi }_G`$, where $`\widehat{H}_{FP}`$ is the FP Hamiltonian defined in Eq. 20. $`F_\nu `$ is used to construct the drift vector $`D^\nu `$ according to Eq. 29, and the drift and local energy are used in evaluating the short-time Green’s function according to Eqs. 34, 35, and 36. When a simple guiding function can be constructed, the expressions for the drift and the local energy can be evaluated analytically, but in most cases this must be done numerically.
$`\mathrm{\S }`$ 3. A set of $`N_w`$ initial configurations or walkers $`\{_j(t=0)\}`$ ($`j=1,\mathrm{},N_w`$) is created, such that particles in each walker are distributed according to the modified importance-sampled distribution $`f(,t=0)=\left[_{i=1}^Ng^{1/2}(i)\right]|\mathrm{\Phi }_G|^2`$ (which is equivalent to using the Variational Monte Carlo (VMC) technique).
$`\mathrm{\S }`$ 4. Each walker $`_j`$ is diffused for a time $`\tau `$ according to the gaussian part of the propagator $`_{i=1}^NG_i^0(_j_j^{};\tau )`$. This can be accomplished by moving each particle coordinate $`x_i^\mu `$ according to
$$x_i^\mu =x_i^\mu +\tau D^\mu (_j)+\sqrt{\tau }\eta ,$$
(55)
where $`\eta `$ is a Gaussian random variable with a mean of zero and a variance of $`2D^{\mu \nu }=2Dg^{\mu \nu }`$.
$`\mathrm{\S }`$ 5. The move from $`_j`$ to $`_j^{^{}}`$ is then accepted with a probability
$$A(_j_j^{};\tau )\mathrm{min}(1,W(_j,_j^{};\tau )),$$
(56)
where
$$W(,^{};\tau )\left[\underset{i=1}{\overset{N}{}}\frac{g(𝐫_i^{^{}})}{g(𝐫_i)}\right]\left|\frac{\mathrm{\Phi }_G(^{^{}})}{\mathrm{\Phi }_G()}\right|^2\frac{G(^{};\tau )}{G(^{};\tau )}.$$
(57)
This step ensures detailed balance in the Monte Carlo procedure. A typical acceptance ratio is in excess of 99%. Notice that if $`G(^{};\tau )`$ is the exact Green’s function, and not its short-time approximation, then $`W`$ is unity and this step is not necessary. This is because the Green’s function of an Hermitian operator is symmetric, i.e. $`G(^{};\tau )=G(^{};\tau )`$, but this is not the case for any importance-sampled distribution function equation.
$`\mathrm{\S }`$ 6. After all the particles of a given walker have been diffused from the initial position $`_j`$ to the position $`_j^{}`$ and the move is accepted, then the values of the drift $`D^\mu `$, the local energy $`E_L`$, and $`\mathrm{\Phi }_G`$ are updated. We could have equally well move all particles at once in step 4 before step 5, however, the acceptance probability for a given time step is reduced considerably. Depending upon the particular problem one can adopt one of the two strategies: single or multiparticle moves.
$`\mathrm{\S }`$ 7. The multiplicity (weight) $`𝖬`$ of a given walker, is computed from the branching part of the Green’s function Eq. 35
$$𝖬=G_b(^{};\tau _a),$$
(58)
where $`\tau _a\tau `$ because some of the moves can be rejected. If the mean-squared distance the particles would diffuse in the absence of the rejection step is $`r^2_{tot}`$, and the actual mean-squared distance is $`r^2_a`$ then the actual time used in a branching step $`\tau _a=\tau r^2_a/r^2_{tot}`$. In the case of multiparticle moves, the effective time step must be calculated separately from a computation of the accepted to attempted moves. Since $`𝖬`$ is, in general, not an integer one can use instead the integer
$`\widehat{𝖬}=\mathrm{int}(𝖬+\xi ),`$where $`\xi `$ is a random number uniform in the range $`[0,1]`$. In this case the average density of walkers is conserved and $`\widehat{𝖬}=𝖬`$. If $`\widehat{𝖬}=0`$ then the walker is deleted from the ensemble; otherwise $`\widehat{𝖬}1`$ copies of the configuration are made and added to the ensemble. Note that fixing $`\widehat{𝖬}=1`$ is equivalent to eliminate branching and, consequently, to perform a VMC calculation with limiting distribution $`f(,t\mathrm{})=\left[_{i=1}^Ng^{1/2}(i)\right]|\mathrm{\Phi }_G|^2`$.
$`\mathrm{\S }`$ 8. After diffusing and branching all walkers in time $`\tau `$ the mean value of the local energy over all walkers is computed from the obtained distribution. One has to start averaging over these values, after some target (equilibration) time, when the configurations are sampled according to the limiting stationary distribution $`f(,t\mathrm{})`$. The target time depends upon the particular problem and how close $`\mathrm{\Phi }_G`$ is to the exact state.
$`\mathrm{\S }`$ 9. Using the average value of the local energy over the whole ensemble of walkers $`E_L()_{f(t\mathrm{})}`$, the trial energy $`E_T`$ is updated according to $`E_T=(E_T+E_L)/2`$. Mixing this estimate with an old value of the trial energy, allows one to improve the convergence.
$`\mathrm{\S }`$ 10. After computing the average energy over a sufficiently long time ($`\tau N_m`$, where $`N_m`$ is the number of moves per configuration), its value is stored and the first block is completed. $`N_m`$ should be large enough for there to be little statistical correlation between energy subaverages obtained in different blocks. A new ensemble of $`N_w`$ walkers is generated by randomly copying or deleting configurations, and steps 4 to 9 are repeated completing another block. One has to do as many blocks as it takes in order to reach the desired statistical accuracy. Notice, however, that because the propagator is only accurate to $`𝒪(\tau ^2)`$, the distribution $`f(,t\mathrm{})`$ and the resulting estimates will have a time step bias. The way to eliminate this bias is by extrapolating all computed expectation values to $`\tau 0`$.
Figure 1 shows a schematic flow diagram of the algorithm presented in this Section.
## V Example: Electron-Monopole in S<sup>2</sup>
As an example application of the method we have developed in the previous Sections we consider the problem of a single particle of charge $`e`$, mass $`m^{}`$ and vector position $`𝐫=(x^1,x^2,x^3)`$ in $`\mathrm{I}\mathrm{R}^3`$ confined to the surface of a sphere of radius $`R`$ centered at the origin ($``$ = S<sup>2</sup>, $`N=1`$) moving in the presence of the vector potential of a Dirac monopole at the origin. This problem can be solved in closed form and so constitutes an ideal model system for testing the accuracy of the stochastic solutions we can derive using the formalism developed in previous Sections. Moreover, the case of $`N`$ interacting electrons confined to S<sup>2</sup> in the presence of a monopole field serves as the basic model which captures the essential physics of the quantum Hall effect .
The Pauli Hamiltonian for a spinless particle in S<sup>2</sup> is
$$\widehat{\mathrm{I}\mathrm{H}}_{}=\frac{\left|\widehat{𝐫}(i\mathrm{}(e/c)𝐀)\right|^2}{2m^{}},$$
(59)
where $`\widehat{𝐫}=𝐫/R`$, and $`𝐀`$ is the monopole vector potential ($`𝐀=B\widehat{𝐫}`$, $`B`$ being the strength of the radial field.) Therefore, the total number of flux quanta $`2𝒮`$ piercing the surface of the sphere is given by $`2𝒮=4\pi R^2B/\varphi _0`$, where $`\varphi _0=hc/|e|`$ is the elementary flux quantum. Following Wu and Yang we can construct angular momentum operators $`𝐋=𝐫(i\mathrm{}(e/c)𝐀)+\mathrm{}𝒮\widehat{𝐫}`$ in terms of which the Hamiltonian reads
$$\widehat{\mathrm{I}\mathrm{H}}_{}=\frac{|𝐋|^2\mathrm{}^2𝒮^2}{2m^{}R^2}.$$
(60)
If we choose a gauge where the vector potential is $`𝐀=BR\mathrm{cot}\theta \widehat{\phi }`$, then the Hamiltonian, Eq. 60, can be written as
$$\widehat{\mathrm{I}\mathrm{H}}_{}=\frac{D}{R^2}\left[_\theta ^2\frac{1}{\mathrm{sin}^2\theta }_\phi ^2\mathrm{cot}\theta _\theta +2i𝒮\frac{\mathrm{cot}\theta }{\mathrm{sin}\theta }_\phi +𝒮^2\mathrm{cot}^2\theta \right],$$
(61)
in terms of the usual spherical angles $`\theta `$ and $`\phi `$ ( $`0\theta \pi `$, $`0\phi <2\pi `$, see Fig. 2.) The eigenstates of this Hamiltonian are monopole harmonics (normalized to 1)
$`Y_{𝒮,n,m}`$ $`=`$ $`𝒩_{𝒮nm}(1)^{𝒮+nm}\mathrm{exp}[i𝒮\phi ]u^{𝒮+m}v^{𝒮m}(|u|,|v|),`$ (62)
$`(|u|,|v|)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{n}{}}}(1)^k\left({\displaystyle \genfrac{}{}{0pt}{}{n}{k}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{2𝒮+n}{𝒮+nmk}}\right)(v\overline{v})^{nk}(u\overline{u})^k,`$ (63)
$`𝒩_{𝒮nm}`$ $`=`$ $`\left({\displaystyle \frac{2𝒮+2n+1}{4\pi }}{\displaystyle \frac{(𝒮+nm)!(𝒮+n+m)!}{n!(2𝒮+n)!}}\right)^{1/2},`$ (64)
where $`u=\mathrm{cos}(\theta /2)\mathrm{exp}[i\phi /2]`$ and $`v=\mathrm{sin}(\theta /2)\mathrm{exp}[i\phi /2]`$ are spinor coordinates, $`n`$ is the Landau level quantum number, and $`m=𝒮n,𝒮n+1,\mathrm{},𝒮+n`$ is the ($`L_{x^3}`$) angular momentum quantum number which labels degenerate states within the $`n^{\mathrm{th}}`$ level. In the sum above the binomial coefficient $`\left(\genfrac{}{}{0pt}{}{\alpha }{\beta }\right)`$ vanish when $`\beta >\alpha `$ or $`\beta <0`$. For a given $`𝒮`$ and $`m`$, the ground state ($`n=0`$) and first excited state ($`n=1`$) are
$`\psi _{gs}`$ $``$ $`u^{S+m}v^{Sm},`$ (65)
$`\psi _{es}`$ $``$ $`\left[2(𝒮+1)v\overline{v}(𝒮m+1)\right]\psi _{gs},`$ (66)
respectively. The energy of a state with Landau level quantum number $`n`$ is given by
$$E_n=\left(2n+1+\frac{n(n+1)}{𝒮}\right)\frac{\mathrm{}\omega _c}{2},$$
(67)
where $`\omega _c`$ is the cyclotron frequency ($`\omega _c=|e|B/m^{}c`$).
Let us now reformulate the electron-monopole problem in a way consistent with the notations introduced in the previous Sections. In this way one can compare the exact result to the numerical one obtained with the algorithm developed in this paper, thus testing the numerical technique. First, instead of the spherical angles $`\theta `$ and $`\varphi `$ we introduce new coordinates $`z`$ and $`\overline{z}`$, where $`z=\mathrm{tan}(\theta /2)\mathrm{exp}[i\phi ]`$, and $`\overline{z}`$ is its complex conjugate. Geometrically, this transformation can be viewed as a stereographic projection of the sphere onto the plane, as illustrated in Fig. 2. The Hamiltonian can be rewritten as
$$\widehat{\mathrm{I}\mathrm{H}}_{}=\frac{i}{m^{}}g^{1/4}(p_z\overline{𝒜}(z))\left(p_{\overline{z}}𝒜(z)\right)g^{1/4}+\frac{D𝒮}{R^2},$$
(68)
in terms of the (non-hermitian) canonical momenta $`p_z`$ and $`p_{\overline{z}}`$, and
$$𝒜(z)=i\frac{\mathrm{}𝒮}{2}z\left(\frac{1|z|^2}{|z|^2(1+|z|^2)}\right),$$
(69)
with metric tensor
$$g^{\mu \nu }(z,\overline{z})=\left(\begin{array}{cc}0& \frac{(1+z\overline{z})^2}{2R^2}\\ \frac{(1+z\overline{z})^2}{2R^2}& 0\end{array}\right).$$
(70)
Naturally, the particles moving in the projected plane are in a space with curvature, corresponding to that of the sphere. Notice that the metric tensor is diagonal when written in terms of ($`\xi ^1,\xi ^2`$), such that $`z=\xi ^1+i\xi ^2`$, i.e. $`g^{\mu \nu }(\xi ^1,\xi ^2)=\frac{(1+|z|^2)^2}{4R^2}\delta ^{\mu \nu }`$ (i.e, it corresponds to the conformal gauge). Then, the drift is simply $`D_\mu =DF_\mu `$.
The stochastic method developed above allows one to obtain the exact energy eigenvalues of the electron-monopole problem in S<sup>2</sup> iff we know the exact phase of the eigenfunctions. In other words, if the trial state is chosen such that it has the exact ground state phase, then independently of its modulus our stochastic approach will lead to the exact ground state energy. Similarly, if the trial state has a phase corresponding to an excited state eigenfunction then we will obtain the exact excited state energy eigenvalue. In the next two subsections we construct simple trial states for the ground and first excited states of the one particle problem. Their modulus are then used as guiding functions $`\mathrm{\Phi }_G`$. Using these trial states we will apply our technique and illustrate the main ideas of our method.
Ground State $`(n=0)`$: The $`2𝒮+1`$ degenerate ground states of the electron-monopole system are labeled by their $`L_{x^3}=\mathrm{}[\overline{z}_{\overline{z}}z_z]`$ angular momentum quantum numbers $`m=𝒮,\mathrm{},𝒮`$. Here we consider the $`m=𝒮`$ ground state for which the exact (unnormalized) wave function can be written
$`\psi _{gs}=\left({\displaystyle \frac{|z|}{z(1+|z|^2)}}\right)^𝒮|\psi _{gs}|e^{i\phi _{gs}}.`$ (71)
To illustrate our method we imagine that we do not know this exact ground state $`\psi _{gs}`$ but instead we have constructed the following two trial states
$$\psi _{T1}=\left(\frac{|z|}{z(1+|z|^2)}\right)^𝒮\frac{1}{1+\lambda |z|^2}|\psi _{T1}|e^{i\phi _{T1}},$$
(72)
and
$$\psi _{T2}=\left(\frac{|z|}{z(1+|z|^2)}\right)^𝒮\left(\frac{|z|}{z}\right)^\alpha |\psi _{T2}|e^{i\phi _{T2}},$$
(73)
where $`\lambda `$ and $`\alpha `$ are real valued constants.
The trial states $`\psi _{T1}`$ and $`\psi _{T2}`$ have been constructed so that for $`\lambda =0`$ and $`\alpha =0`$ they are both equal to the exact ground state, $`\psi _{gs}`$. For $`\lambda 0`$, the modulus of $`\psi _{T1}`$ is no longer equal to that of the exact ground state, but the phase of the wave function is exact, i.e.
$`|\psi _{T1}||\psi _{gs}|,\phi _{T1}=\phi _{gs}.`$ (74)
In contrast, for $`\alpha 0`$, the modulus of $`\psi _{T2}`$ is exact, but the phase is approximate,
$`|\psi _{T2}|=|\psi _{gs}|,\phi _{T2}\phi _{gs}.`$ (75)
It follows that if $`\psi _{T1}`$ is used as a trial state in a FP DMC simulation the resulting energy will be the exact ground state energy $`E_0=\mathrm{}\omega _c/2`$, while if $`\psi _{T2}`$ is used as a trial state the simulation will not lead to the exact ground state energy, but instead will provide a variational upper bound.
As has already been emphasized the trial state used in a FP DMC simulation should be constructed to be the best available approximation to the exact eigenstate, since the quality of the trial state can greatly influence the speed of convergence and the statistical accuracy of the result of the simulation. This can be clearly illustrated by considering the trial state $`\psi _{T2}`$. Since the modulus of $`\psi _{T2}`$ is exact the drift velocity $`F`$ which depends only on the modulus of the trial state will correspond to the exact drift velocity and in the absence of the branching term will lead to the exact density distribution. It is straightforward to show that
$`F_1={\displaystyle \frac{4𝒮\xi ^1}{1+|z|^2}},F_2={\displaystyle \frac{4𝒮\xi ^2}{1+|z|^2}},`$ (76)
where $`F_\mu =_{\xi ^\mu }\mathrm{ln}|\psi _{T2}|^2`$, indicating that walkers are guided away from the regions where the wave function is small and, in this way, the particle tends to spend most of the time near the top of the sphere ($`\theta =0`$). A potential problem appears when we consider the local energy,
$`E_L=|\psi _{T2}|^1\widehat{H}_{FP}|\psi _{T2}|={\displaystyle \frac{\mathrm{}\omega _c}{2}}\left[1+{\displaystyle \frac{(1+|z|^2)^2}{𝒮}}\left({\displaystyle \frac{\alpha 𝒮}{1+|z|^2}}+{\displaystyle \frac{\alpha ^2}{4|z|^2}}\right)\right]`$ (77)
which, of course, is not exact due to the approximate phase of the trial state. In particular, $`E_L`$ diverges as $`|z|0`$. As we have just shown, the drift will tend to push the particle towards $`z=0`$, leading to large fluctuations in the local energy. This in turn can lead to huge fluctuations in the population size (number of walkers), since the size of the population depends exponentially on the local energy. Thus, in this particular example, one has to take small values for $`\alpha `$ ($`\alpha <<1`$) in order to assure fast convergence and good statistical accuracy.
Figure 3 shows the results of FP DMC simulations, using the algorithm developed in this paper, for the difference between computed and exact ground state energies for trial state $`\psi _{T1}`$ (circles) and trial state $`\psi _{T2}`$ (squares) using different values of the time step $`\tau `$. The $`\tau 0`$ extrapolated values are also shown. For trial state $`\psi _{T1}`$ we used $`\lambda =1`$, the number of walkers was chosen to be $`N_w=300`$, the number of Monte Carlo steps per walker was $`2\times 10^7`$, and the acceptance rate was between $`97\%`$ and $`99.5\%`$. For trial state $`\psi _{T2}`$ we used $`\alpha =0.001`$, the number of walkers was $`N_w=100`$, the total number of Monte Carlo steps per walker was $`10^6`$, and the acceptance rate was the same as for $`\psi _{T1}`$. Since the parameter $`\alpha `$ in $`\psi _{T2}`$ was chosen so that $`\alpha 1`$, the trial state and the corresponding Green’s function were nearly exact and we were able to reach reasonable statistical accuracy with a relatively small number of Monte Carlo steps. As expected we find that when trial state $`\psi _{T1}`$ is used the extrapolated energy agrees within statistical accuracy with the exact result, but when we use trial state $`\psi _{T2}`$, for which the phase is not exact, we obtain a variational upper bound for the exact ground state energy.
In Fig. 4 the density profiles for the exact ground state $`\psi _{gs}=\psi _{T1}(\lambda =0)`$ (Exact), the trial state $`\psi _{T1}(\lambda =1)`$ (VMC), the density obtained in FP DMC with trial state $`\psi _{T1}(\lambda =1)`$ at time step $`\tau =0.001`$ (FP mixed estimator), and the extrapolated density defined as ratio of the square of FP density to the variational density corresponding to $`\psi _{T1}(\lambda =1)`$, are shown. Note that since the density in our DMC calculation is determined as a mixed estimate (see Eq. 15), and the density operator does not commute with the Hamiltonian between the DMC solution and the trial state, the corresponding density profile (FP mixed estimator) improves on the variational result but still differs from the exact density. The extrapolated estimator for the density constructed by combining both, the FP mixed estimator and the variational density makes it possible to improve on the FP density and is seen to be very close to the exact result.
First Excited State $`(n=1)`$: As a further demonstration of the validity of our method we turn to the first excited state of the electron-monopole system. Again, we specify the $`L_{x^3}`$ angular momentum quantum number and take $`m=𝒮+1`$ for which the exact excited state wave function is
$`\psi _{es}=\left({\displaystyle \frac{|z|}{z(1+|z|^2)}}\right)^{𝒮+1}|z|=|\psi _{es}|e^{i\phi _{es}}.`$ (78)
We then introduce two new trial states
$$\psi _{T1}=\left(\frac{|z|}{z(1+|z|^2)}\right)^{𝒮+1}\frac{|z|}{1+\lambda |z|^2}=|\psi _{T1}|e^{i\phi _{T1}}$$
(79)
and
$$\psi _{T2}=\left(\frac{|z|}{z(1+|z|^2)}\right)^{𝒮+1}|z|\left(\frac{|z|}{z}\right)^\alpha =|\psi _{T2}|e^{i\phi _{T2}}$$
(80)
with the property that for $`\lambda =0`$ and $`\alpha =0`$ they each reduce to the exact excited state, $`\psi _{es}`$. As before, for $`\lambda 0`$ and $`\alpha 0`$ the modulus of $`\psi _{T1}`$ is approximate and the phase is exact ($`|\psi _{T1}||\psi _{es}|`$, $`\phi _{T1}=\phi _{es}`$) and the modulus of $`\psi _{T2}`$ is exact and the phase is approximate ($`|\psi _{T2}|=|\psi _{es}|`$, $`\phi _{T2}\phi _{es}`$).
Figure 5 shows the difference between FP DMC energies computed using trial states $`\psi _{T1}`$ and $`\psi _{T2}`$ and the exact excited state energy for different values of time step $`\tau `$ as well as the extrapolated $`\tau =0`$ result. The parameters (number of walkers, number of Monte Carlo steps, etc.) used for these simulations were the same as those used for the ground state simulations except that we took $`\alpha =0.0015`$ in $`\psi _{T2}`$. Again, our simulations gave the expected results — when the phase of the trial state is exact we obtain the exact energy (circles), $`E_1=(3/2+1/𝒮)\mathrm{}\omega _c`$, and when the phase is approximate we obtain a variational upper bound on that energy (squares).
In Figure 6 the density profiles corresponding to the trial state $`\psi _{T1}(\lambda =1)`$ (VMC), the FP density using the same state at time step $`\tau =0.001`$ (FP mixed estimator), the extrapolated density, computed as above by taking the ratio of square of the FP density and the VMC density (Extrap. estimator), and the exact density (Exact). The results are qualitatively similar to those for the ground state – the FP estimator improves on the VMC result, and the extrapolated density is nearly equal to the exact excited state density.
The simulation results presented in this Section provide a simple test of both the FP DMC method and the method developed in this paper for dealing with quantum corrections due to curvature. The results clearly show that these methods can be used to study quantum systems on curved manifolds.
## VI Discussion and conclusions
In this paper we have introduced a stochastic method to solve the many-body Schrödinger equation on curved manifolds. This method is essentially a generalized Diffusion Monte Carlo (DMC) technique allowing one to deal with the effects arising from the space curvature. We have shown that due to the curvature the diffusion matrix and drift vector, which appear in the Green’s function used as a conditional probability in DMC simulations, acquire additional terms, the so-called quantum corrections. The explicit expressions for a general metric tensor are worked out in detail. Since the presence of the curvature leads to a number of other nontrivial modifications, we have presented a step by step algorithm which can be used to implement a code dealing with DMC simulations in curved space.
It is worth emphasizing that our method can be applied to a wide variety of inhomogeneous systems (e.g., inhomogeneous semiconductors with a position dependent effective mass), not just systems on curved manifolds. The reason for this, as discussed in Section III, is that the quantum corrections to the Green’s function can be interpreted as being due to those terms which appear in the generalized diffusion equation describing the system once each of the derivatives in that equation have been commuted all the way to the left of each term. This definition of quantum corrections is quite general and can be applied to any differential equation which is second order in space and first order in time, regardless of the number of dimensions or any spatial inhomogeneity in the system.
To illustrate the general methodology we have concentrated on the problem of interacting fermions in external electromagnetic potentials. In this case a variational upper bound to the exact ground state energy can be found by applying the Fixed-Phase approximation, where the fermionic problem is treated as a bosonic one by fixing the phase of the many-body wave function (which is complex-valued in general) by some trial phase. As an example, we have considered the problem of a single electron confined to the surface of a two-sphere, which has a magnetic monopole at its center. The electron thus moves in a space with curvature in the presence of a magnetic field which breaks time-reversal symmetry. This simple problem can be solved in closed form and, therefore, we have used it as a playground for testing our technique. In the paper we have presented two calculations, where the ground and first excited state energies are computed using the exact phases, but approximate modulus for the corresponding guiding functions. We have shown that the exact energies are reproduced within statistical accuracy thus proving that the approach for dealing with the quantum corrections is valid.
As emphasized in the Introduction, the method presented in this paper for performing DMC simulations on curved manifolds can be used to study many interesting physical systems. An important example is the quantum Hall effect, a phenomenon which occurs when a two-dimensional electron system is placed in a strong magnetic field. As first pointed out by Haldane , the electron-monopole system described in Section V provides a convenient geometry for performing finite size numerical studies of quantum Hall systems when many interacting electrons are placed on the sphere. This is in part because the spherical geometry has no boundary so that finite size effects are suppressed. In addition, the spherical geometry is conceptually simpler than the (flat metric) torus geometry, which also has no boundary, because the topological order exhibited by quantum Hall states leads to certain nontrivial degeneracies on the torus .
Recently we have used the method developed in this paper to study some of the exotic excitations which occur in quantum Hall systems, specifically the fractionally charged quasiparticle excitations of the fractional quantum Hall effect , and the charged spin texture excitations (skyrmions) of the integer quantum Hall effect . Previous numerical studies of these excitations have been based on either VMC or exact diagonalization calculations which, for the most part, have assumed that the wave functions describing the excitations are confined to the lowest ($`n=0`$) Landau level. In fact, this is a rather poor approximation for real experimental systems which can exhibit significant mixing of higher Landau levels due to the electron-electron interaction. Because the FP DMC method allows one to go beyond the lowest Landau level approximation it can be used to study the effect of Landau level mixing on quantum Hall states and, by employing the method described in this paper, we were able to perform such studies using Haldane’s spherical geometry. These simulations have provided useful quantitative results for various properties of quantum Hall excitations for realistic experimental parameters . Along with the more rigorous test case of the electron-monopole system described in this paper, these calculations of Landau level mixing effects in quantum Hall systems using the Haldane sphere have shown that the method we have developed for performing FP DMC calculations on curved manifolds works well and can be used to study many other interesting physical systems.
###### Acknowledgements.
We acknowledge stimulating discussions with David Ceperley. VM is supported by NSF grant number DMR-9725080. NEB is supported by US DOE grant number DE-FG02-97ER45639. Work at Los Alamos is sponsored by the US DOE under contract W-7405-ENG-36. |
warning/0001/nlin0001024.html | ar5iv | text | # Periodic orbit action correlations in the Baker map
## I Introduction
Numerical evidence suggests that eigenvalue spectra of individual low dimensional quantum systems show correlations which are solely determined by the underlying classical dynamics, symmetries and Planck’s constant $`\mathrm{}`$ (see e.g. Bohigas et al 1984, Berry 1987, Bohigas 1991). Especially quantum systems exhibiting fully chaotic or completely regular dynamics in the classical limit show spectral fluctuations coinciding with random matrix theory (RMT). Semiclassical trace formulae provide a link between the set of eigenenergies and the set of all periodic orbits in the classical system and have proven to be important in understanding universality in spectral statistics and deviations thereof (Hannay and Ozorio de Almeida 1984, Berry 1985, Berry 1988, Aurich and Steiner 1995, Bogomolny and Keating 1996a). A complete description of spectral properties in terms of the underlying classical dynamics of the system is, however, still missing.
Universal spectral statistics is intimately connected to correlations in periodic orbit length or actions distributions. The validity of the random matrix conjecture implies generalised universal periodic action correlation for systems for which the trace formula is exact (Argaman et al 1993). Keating (1993) applied this idea to the Riemann zeta function whose non-trivial zeros are conjectured to follow random matrix statistics. Here, the Hardy - Littlewood conjecture provides an explicit expression for the pair - correlations of prime numbers, (the analogue of periodic orbits in the Riemann case). Expressions for the two-point correlation function of Riemann zeros can be obtained making use of the conjectured prime number correlations which coincide asymptotically with RMT results for an ensemble of Gaussian random matrices invariant under unitary transformation (GUE). A generalisation of these arguments to n-point correlation functions for Riemann-zeros has been given by Bogomolny and Keating (1995, 1996).
Far less is known about action correlations of periodic orbits for generic chaotic systems and even numerical evidence for the existence of these universal periodic orbit correlations is weak. Argaman et al (1993) present numerical results which agree qualitatively with the RMT-prediction presented in the same paper; Cohen et al (1998) find correlations in the length spectrum of periodic orbits in the Stadium billiard which indicate RMT-like correlations but a quantitative analysis can not compete with the accuracy reached for spectral eigenvalue statistics. Other groups report complete failure in their search for any kind of action correlations (Aurich 1998) or find action correlations which do not coincide with the RMT - prediction by Argaman et al (1993), see Sano (1999); in addition there are yet no arguments other than from the duality between quantum eigenvalues and classical periodic orbits which indicate why universal classical correlations should exist. Cohen et al (1998) present a variety of ideas based on identifying relationships within various subgroups of orbits in the Sinai billiard. The results are, however, still of speculative nature.
In this paper I will study quantum and semiclassical spectra and discuss necessary conditions for the existence of universal periodic orbit action correlations. I will focus in particular on the classical and quantum Baker map. It will become evident that the periodic orbit correlations deviate from the universal behavior for long orbits due to the violation of quantum unitarity in the semiclassical approximations. Correlations following the RMT-prediction do, however, exist for short periodic orbits which can not be explained by classical sum rules alone. Properties of periodic orbits can be studied in detail for the Baker map by writing periodic orbit sums in terms of a suitable classical Perron-Frobenius operator. Exponentially increasing terms in the large action limit can be regularised by imposing unitarity onto a semiclassical quantization as proposed by Bogomolny and Keating (1996a).
## II Periodic Orbit Correlation Functions for Quantum Maps
In the following I will limit the discussion to quantum maps. A quantum map acts on a finite dimensional Hilbert space and the quantum dynamics is governed by the equation
$$\psi _{n+1}=𝐔\psi _n$$
(1)
with $`𝐔`$ being a unitary matrix of dimension $`N`$. We assume that the map has a well defined classical limit for $`N\mathrm{}`$ described by a discrete dynamical map. $`\psi _n`$ is the discretised $`N`$ \- dimensional wave vector and $`N`$ is equivalent to the inverse of Planck’s constant $`h`$, i.e. $`N1/h`$. For possible quantization procedures of classical maps see e.g. Hannay and Berry (1980) for the cat map, Balazs and Voros (1989) and Saraceno and Voros (1994) for the Baker map or Bogomolny (1992), Doron and Smilansky (1991), and Prosen (1994, 1995) for semiclassical and quantum Poincaré maps. The spectral density of eigenphases $`\theta _i`$ of $`𝐔`$ can be written in the form
$$d(\theta ,N)=\underset{i=1}{\overset{N}{}}\delta _{2\pi }(\theta \theta _i)=\frac{N}{2\pi }+\frac{1}{\pi }Re\underset{n=1}{\overset{\mathrm{}}{}}\text{Tr}𝐔^ne^{in\theta }$$
(2)
with $`\overline{d}=N/2\pi `$ being the mean density of eigenphases in the interval $`[0,2\pi ]`$ and $`\delta _{2\pi }`$ denotes the periodically continued delta function with period $`2\pi `$. The spectral measure of interest here is the two-point correlation function defined as
$$R_2(x,N)=\frac{1}{\overline{d}^2}<d(\theta )d(\theta +x/\overline{d})>.$$
(3)
The average is taken over the unit circle, i.e. $`<.>=\frac{1}{2\pi }_0^{2\pi }.d\theta `$ which leads to
$$R_2(x,N)=\frac{1}{N}\underset{i,j=1}{\overset{N}{}}\delta _{2\pi }(x_ix_jx)=1+\frac{2}{N^2}\underset{n=1}{\overset{\mathrm{}}{}}|\text{Tr}𝐔^n(N)|^2\mathrm{cos}(2\pi \frac{n}{N}x);$$
(4)
the spectrum $`\{x_i=\overline{d}\theta _i;i=1,2\mathrm{}N\}`$ is unfolded to mean level density one. Note that $`R_2`$, as defined in (3), is still a distribution and further averaging has to be applied. Fourier transformation of the two-point correlation function yields the so-called form factor which can be written as
$$K(\tau ,N)=\{\begin{array}{ccc}\frac{1}{N}|\text{Tr}𝐔^{N\tau }(N)|^2\hfill & \text{for}\hfill & \tau 0\hfill \\ N\hfill & \text{for}\hfill & \tau =0;\hfill \end{array}$$
(5)
the variable $`\tau `$ takes the discrete values $`\tau =n/N`$, $`n`$, $`N\text{I}\text{N}`$ which reflects the periodicity of the two-point correlation function $`R_2`$. Note that the form factor does not converge in the limit $`N\mathrm{}`$ but fluctuates around a mean value. The limit distribution of $`K(\tau ,N)`$ for $`N\mathrm{}`$ is, however, expected to exist and is conjectured to be Gaussian (Prange 1997). I will refer to the form factor as the mean value of the distribution (5) near a given $`\tau `$ value, i.e. further averaging over small $`\tau `$-intervals has to be performed.
The traces of $`𝐔^n`$ can in semiclassical approximation be written as sum over periodic orbits of topological length $`n`$ of the underlying classical map (Gutzwiller 1990, Bogomolny 1992), i.e.
$$\text{Tr}𝐔^n(N)\text{Tr}\stackrel{~}{𝐔}^n(N):=\underset{p}{\overset{(n)}{}}A_pe^{2\pi iNS_p}.$$
(6)
The complex prefactors $`A_p`$ contain information about the stability of the periodic orbit and $`S_p`$ denotes the action or generating function of the map corresponding to the periodic orbit $`p`$. The relation between quantum traces and semiclassical periodic orbit sums is in general obtained by stationary phase approximation and is exact only in the limit $`\tau ^1=N/n\mathrm{}`$ (Andersson and Melrose 1997). Exceptions thereof are e.g. the cat map (Hannay and Berry 1980, Keating 1991) or quantum graphs (Kottos and Smilansky 1997) for which Eq. (6) is exact for all $`\tau `$.
The relation (6) makes it possible to connect purely classical action correlation functions with the statistics of quantum eigenvalues. A weighted action correlation function can be obtained by considering
$$P(s,n)=\underset{N=\mathrm{}}{\overset{\mathrm{}}{}}|\text{Tr}\stackrel{~}{𝐔}^n(N)|^2e^{2\pi isN}=\underset{p,p^{}}{\overset{(n)}{}}A_pA_p^{}^{}\delta _{2\pi }(S_pS_p^{}s)$$
(7)
which is a Fourier sum in $`N`$ depending only on classical quantities such as the topological length of the orbit, $`n`$, and the actions $`S_p`$ and weights $`A_p`$. Eqs. (4) and (7) together with (5) indicate that energy level statistics and periodic orbit action correlations are connected provided that both the quantum traces $`\text{Tr}𝐔^n`$ and the semiclassical ’traces’ $`\text{Tr}\stackrel{~}{𝐔}^n`$ follow the same probability distributions, i.e. random matrix statistics. Little is known about the statistical properties of semiclassical traces and eigenvalues and I will study the existence of universal RMT-behavior for semiclassical expressions in more detail in section IV.
The RMT-conjecture implies for quantum traces $`\text{Tr}𝐔^n(N)`$ of systems whose classical limit is chaotic and non-time reversal symmetric
$$\underset{N\mathrm{}}{lim}\frac{1}{N}<|\text{Tr}𝐔^n(N)|^2>=\{\begin{array}{ccc}n/N\hfill & \text{for}\hfill & nN\hfill \\ 1\hfill & \text{for}\hfill & n>N\hfill \end{array};$$
(8)
the brackets $`<.>`$ indicate the average over an $`n`$-interval small compared to the dimension $`N`$. Assuming that the above relation holds also for the semiclassical traces $`<|\text{Tr}\stackrel{~}{𝐔}^n(N)|^2>`$ we can insert (8) in (7) which leads to the asymptotic result (Argaman et al 1993)
$$P(s,n)=\overline{P}(n)\left(\frac{\mathrm{sin}n\pi s}{\pi s}\right)^2+n\delta (s)\text{for}|s|1.$$
(9)
Here, $`\overline{P}(n)=_{p,p^{}}^{(n)}A_pA_p^{}^{}`$ is the mean part of the weighted periodic orbit action pair density. Note that the number of periodic orbits increase exponentially with $`n`$ for chaotic maps which in turn implies an exponential increase for the density of periodic orbit actions (modulus 1) and in general also for the weighted density $`\overline{P}(n)`$.
After rescaling the $`s`$ \- variable according to $`s=\sigma /n`$, Eq. (9) can be written as
$$P_{scal}^{GUE}(\sigma )=\frac{1}{n^2}P\left(\frac{\sigma }{n}\right)=\frac{\overline{P}(n)}{n^2}\left(\frac{\mathrm{sin}\pi \sigma }{\pi \sigma }\right)^2+\delta (\sigma )\text{for}|\sigma |n,$$
(10)
which puts the non-trivial, $`\sigma `$-dependent correlations into $`n`$ \- independent form.
The corresponding equation for systems with time-reversal symmetry is (Argaman et al 1993)
$`P_{scal}^{GOE}(\sigma )={\displaystyle \frac{\overline{P}(n)}{n^2}}`$ $``$ $`2\left({\displaystyle \frac{\mathrm{sin}\pi \sigma }{\pi \sigma }}\right)^2`$ (11)
$`+`$ $`{\displaystyle \frac{2}{\pi \sigma }}[\mathrm{cos}2\pi \sigma (\text{si}(2\pi \sigma )\mathrm{cos}2\pi \sigma \text{Ci}(2\pi \sigma )\mathrm{sin}2\pi \sigma )`$ (12)
$`+`$ $`\text{Ci}(4\pi \sigma )\mathrm{sin}4\pi \sigma \text{si}(4\pi \sigma )\mathrm{cos}4\pi \sigma ]+2\delta (\sigma )\text{for}|\sigma |n,`$ (13)
with si($`x`$), Ci($`x`$) being the Sine - and Cosine - integrals.
For integrable systems conjectured to have Poissonian spectral statistics one expects
$`{\displaystyle \frac{1}{N}}<|\text{Tr}𝐔^n(N)|^2>=<{\displaystyle \underset{p,p^{}}{\overset{(n)}{}}}A_pA_p^{}^{}e^{2\pi iN(S_pS_p^{})}>=1\text{for}|N|0`$
and thus
$$P^{Pois}(s,n)=\underset{N=\mathrm{}}{\overset{\mathrm{}}{}}\frac{1}{N}|\text{Tr}\stackrel{~}{𝐔}^n(N)|^2e^{2\pi isN}=\overline{P}(n)1+\delta (s),$$
(14)
i.e. periodic orbit actions are uncorrelated.
The $`\delta `$-functions in (9) – (14) can be identified with the diagonal terms $`p=p^{}`$ in the sum (7) after making use of the classical sum rule (Hannay and Ozorio de Almeida 1984)
$$\underset{p}{\overset{(n)}{}}|A_p|^21\text{for large }n.$$
(15)
Before investigating the existence of correlations of the form (10) or (11) in more detail, I would like to make a few remarks: the weighted periodic orbit correlation function (7) contains a background term $`\overline{P}(n)`$ which increases in general exponentially with $`n`$. This is in contrast to the linear behaviour of the mean quantum level density $`\overline{d}(N)`$ in (2) and (3). Note also, that the complex weights $`A_p`$ carry phases which may result in cancellations in $`\overline{P}`$ and may occasionally lead to $`\overline{P}`$ = 0. The correlation functions (10) and (11) are in addition not rescaled with respect to the mean action density as e.g the two-point spectral correlation function $`R_2`$ in (3). The differences in action between adjacent orbits of the same length $`n`$ is exponentially small on the $`s`$ or $`\sigma `$ scale and the term $`(\mathrm{sin}\pi \sigma /\pi \sigma )^2`$ in (10), (11) indicates long range correlations on the scale of the mean periodic action separation. The correlations (10) and (11) imposed by random matrix theory are thus a small modulation of order $`𝒪(1)`$ on top of an exponentially large background when considering all periodic orbit pairs (Dahlqvist 1995). The two-point correlation function for periodic orbit actions approaches the Poissonian limit for $`n\mathrm{}`$ when rescaling $`P(\sigma ,n)`$ with respect to the mean action density in accordance to e.g. the definition of the spectral correlation function (3); periodic orbits are thus uncorrelated on scales of the mean action separation in agreement with the numerical results by Harayama and Shudo (1992).
The search for an $`𝒪(1)`$ – effect on top of an exponentially large background is one of the main obstacles when investigating RMT-induced classical action correlations numerically. The second major problem in studying periodic orbit properties especially in the large $`n`$ \- limit is the exponential increase in the number of orbits with increasing $`n`$. In the following I will focus on a specific example, the classical and quantum Baker map, in which the problems mentioned earlier can be overcome by constructing a quasiclassical Perron–Frobenius type operator (Dittes et al 1994).
## III The Baker Map
The Baker map has become a standard example when studying classical and quantum chaos (see e.g. Balazs and Voros 1989, Saraceno and Voros 1994, O’Connor et al 1992, Hannay et al 1994). The classical dynamics is given by a two-dimensional, piecewise linear, area-preserving map on the unit square defined as
$`q^{}`$ $`=`$ $`2qϵ`$ (16)
$`p^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(p+ϵ)\text{with}ϵ=[2q]`$ (17)
and $`(q,p)`$, $`(q^{},p^{})`$ being the initial and final points. The notation $`[x]`$ stands for the integer part of $`x`$. The Baker map is hyperbolic with a well defined Markov - partition and a complete binary symbolic dynamics given in terms of the $`ϵ=\{0,1\}`$. The map is invariant under the anti-unitary symmetry $`T:(q,p)(p,q)`$ (being equivalent to time reversal symmetry) and a parity transformation $`R:(q,p)(1q,1p)`$. A proper desymmetrisation of quantum spectra and semiclassical periodic orbit sums with respect to parity will be crucial when studying spectral statistics and periodic orbit correlations in section IV.
Each periodic orbit of the map can be associated with a finite symbol string $`\mathit{ϵ}=(ϵ_1,ϵ_2,\mathrm{},ϵ_n)`$ and the phase space coordinates of the orbit are given by the expression
$$q_\mathit{ϵ}=\frac{1}{2^n1}\underset{i=0}{\overset{n1}{}}ϵ_i2^{ni1};p_\mathit{ϵ}=\frac{1}{2^n1}\underset{i=0}{\overset{n1}{}}ϵ_i2^i.$$
(18)
A suitable generating function of the map is
$$W_ϵ(p^{},q)=2p^{}qϵp^{}ϵq\text{with}ϵ=[2q],$$
(19)
which in turn defines the action; the action of a periodic orbit $`\mathit{ϵ}`$ of the map can up to an additive integer be written as (Dittes et al 1994)
$$S_\mathit{ϵ}=\underset{i=1}{\overset{n}{}}(q_i1)ϵ_i=\underset{i=1}{\overset{n}{}}(q_i1)[2q_i]$$
(20)
with $`q_i`$ being the $`q`$-coordinates of the periodic orbit. The sum in (20) is invariant under time reversal symmetry and parity-transformation $`R(q_i)=1q_i`$; the later can be shown using the relation $`_{i=1}^nq_i=_{i=1}^nϵ_i`$.
A quantized version of the Baker map is obtained by making use of a discretised version of the generating function (19) and the classical structure of the map in the mixed representation. A suitable formulation preserving all the classical symmetries is provided by the choice (Balazs and Voros 1989, Saraceno and Voros 1994)
$$𝐔(N)=𝐅_N^1\times \left(\begin{array}{cc}𝐅_{N/2}& 0\\ 0& 𝐅_{N/2}\end{array}\right)$$
(21)
with
$$(𝐅_N)_{\mathrm{ij}}=\frac{1}{\sqrt{N}}e^{2\pi i(\mathrm{i}\frac{1}{2})(\mathrm{j}\frac{1}{2})/N}\mathrm{i},\mathrm{j}=1,\mathrm{}N.$$
(22)
The Fourier-matrix $`𝐅`$ is nothing but the transformation from position to momentum representation in the finite dimensional Hilbert space. The dimension $`N`$ of the map $`𝐔(N)`$ has to be even in this construction and $`N`$ is equivalent to the inverse of Planck’s constant. The unitary map $`𝐔(N)`$ commutes with the parity operator $`(𝐑_N)_{\mathrm{i},\mathrm{j}}=\delta _{\mathrm{i},N\mathrm{j}+1}`$ due to the particular choice of half-integer phases (Saraceno and Voros 1994). Symmetry reduction is obtained by considering the matrices $`𝐔_\pm (N)`$ defined as
$$𝐔_\pm (N)=\frac{1}{2}\left(𝐔(N)\pm 𝐔(N)𝐑_N\right)$$
(23)
acting on the symmetric/anti-symmetric wave-vectors only. This leads effectively to a reduction in dimension by a factor of 2. The traces of $`𝐔_\pm ^n`$ can in the large $`N`$-limit be written as (Saraceno and Voros 1994, Toscano et al 1997)
$$\text{Tr}𝐔_\pm ^n(N)=\frac{1}{2}\underset{\mathit{ϵ}}{\overset{(n)}{}}\left(\frac{2^{n/2}}{2^n1}e^{2\pi iNS_\mathit{ϵ}}\pm \frac{2^{n/2}}{2^n+1}e^{2\pi iNS_\mathit{ϵ}^{}+i\pi }+a_\mathit{ϵ}\mathrm{log}N+b_\mathit{ϵ}\right)+𝒪(N^{1/2})+\mathrm{}.$$
(24)
The first two terms in the sum over the possible finite symbol strings of length $`n`$ are the usual semiclassical Gutzwiller-like periodic orbit contributions. They arise as stationary phase points in a continuum approximation of $`\text{Tr}𝐔^n`$, $`\text{Tr}(𝐔^n𝐑)`$, respectively, i.e. sums over matrix elements are replaced by integrals (Tabor 1983). The action $`S_\mathit{ϵ}`$ of a periodic orbit of length $`n`$ is given by Eq. (20). $`S_\mathit{ϵ}^{}`$ corresponds to half the action of an orbit of length $`2n`$ with symbol code ($`ϵ_1,\mathrm{}ϵ_n,\overline{ϵ_1},\mathrm{}\overline{ϵ_n}`$) and $`\overline{ϵ_i}=1ϵ_i`$. These are the orbits being invariant under the classical parity transformation $`R`$ and are hence the stationary phase contributions of $`\text{Tr}(𝐔^n𝐑)`$. The anomalous $`a_\mathit{ϵ}\mathrm{log}N+b_\mathit{ϵ}`$ corrections are due to the discretisation of the phase-space and arise in ’sub leading’ integrals when performing Poisson-summation on $`\text{Tr}𝐔^n`$ i.e. replacing sums by sums over integrals (Saraceno and Voros 1994, Toscano et al 1997)). Next leading terms are diffraction corrections arising from the discontinuous nature of the classical and quantum map (16) and (21). Note the extra phase $`\pi `$ in the $`\text{Tr}(𝐔^n𝐑)`$ \- contributions which originates from the discretised stationary phase approximation.
In the following I will mainly concentrate on the Gutzwiller - like periodic orbit contributions, i.e. I will consider the approximation
$`\text{Tr}𝐔^n(N)`$ $``$ $`\text{Tr}\stackrel{~}{𝐔}^n(N)={\displaystyle \frac{2^{n/2}}{2^n1}}{\displaystyle \underset{\mathit{ϵ}}{\overset{(n)}{}}}e^{2\pi iNS_\mathit{ϵ}}`$ (25)
$`\text{Tr}𝐔_\pm ^n(N)`$ $``$ $`\text{Tr}\stackrel{~}{𝐔}_\pm ^n(N)={\displaystyle \frac{1}{2}}\left[{\displaystyle \frac{2^{n/2}}{2^n1}}{\displaystyle \underset{\mathit{ϵ}}{\overset{(n)}{}}}e^{2\pi iNS_\mathit{ϵ}}{\displaystyle \frac{2^{n/2}}{2^n+1}}{\displaystyle \underset{\mathit{ϵ}}{\overset{(n)}{}}}e^{2\pi iNS_\mathit{ϵ}^{}}\right].`$ (26)
I will henceforth refer to the periodic orbit sums (25) as the semiclassical approximation of the quantum Baker map. The influence of the remaining contributions in (24), including the leading order ($`\mathrm{log}N`$) - corrections, will be discussed at the end of this section.
The main problem in calculating periodic orbit sums like (25) is the exponential increase of periodic orbit contributions. Even for the Baker map where periodic orbit actions are given by the simple analytic formula (20) a direct summation of $`\text{Tr}\stackrel{~}{𝐔}^n`$ becomes a computational challenge for $`n30`$. For the Baker map it is, however, possible, to construct a quasiclassical operator $`\stackrel{~}{U}`$ whose traces coincide with the semiclassical periodic orbit sum (25) (Dittes et al 1994). The spectrum of this operator can be computed explicitly and simplifies the task of calculating traces, i.e. periodic orbit sums, considerably. Other quasiclassical operators for the Baker map have been proposed by Kaplan and Heller (1996), Sano (1999) and Hannay (1999). Note, that an identification of semiclassical periodic orbit sums with classical or quasiclassical operator is not possible for generic systems and the notation $`\text{Tr}\stackrel{~}{𝐔}^n`$ is then a mere substitution for the periodic orbit sum itself.
The Baker map is special for the following reasons: firstly, the periodic orbit amplitudes $`A_p`$, (cf. Eq. (6)), do not depend on the specific orbit, but only on the topological length $`n`$, which is a consequence of the piecewise linearity of the classical map. Secondly, the dynamics of the $`q`$ \- coordinate as well as the periodic orbit - action (20) are independent of the momentum $`p`$ which allow for a classical separation of momentum and configuration space variables. As a consequence, one can define a one dimensional Perron-Frobenius type integral kernel acting on the $`q`$-coordinate only (Dittes et al 1994). In a form preserving the parity-symmetry $`R(q)=1q`$ it can be written as
$$\stackrel{~}{U}(q,q^{};N)=\sqrt{2}\delta \left(q(2q^{}[2q^{}])\right)e^{2\pi iNS(q^{})}$$
(27)
with
$$S(q^{})=q^{}[2q^{}]\frac{1}{2}([2q^{}]+q^{}).$$
(28)
One easily deduces $`\stackrel{~}{U}(q,q^{})=\stackrel{~}{U}(Rq,Rq^{})`$. The traces of $`\stackrel{~}{U}(q,q^{})`$ coincide with the periodic orbit sum (25) if the operator is defined on the space of analytic functions on the unit interval with periodic boundary conditions (Dittes et al 1994, Rugh 1992). The operator is thus infinite dimensional (having infinitely many eigenvalues) in contrast to the finite dimensional quantum matrix $`𝐔(N)`$ and $`N`$ is a mere parameter in (27). A symmetry-reduced version of this quasiclassical operator is obtained by considering
$$\stackrel{~}{U}_\pm (q,q^{};N)=\frac{1}{2}\left(\stackrel{~}{U}(q,q^{};N)\stackrel{~}{U}(q,Rq^{};N)\right).$$
(29)
Note, that the quasiclassical operator acting on antisymmetric functions (denoted $`\stackrel{~}{U}_+`$ here) corresponds to the symmetric quantum subspace $`𝐔_+`$ and vice versa, reflecting the ‘quantum’ origin of the extra phase for the periodic orbit contributions to $`\text{Tr}𝐔^n𝐑`$, see Eq. (24). The operator (27) or (29) can now be represented as an infinite dimensional matrix after choosing a suitable basis set and the calculation of traces can be reduces to matrix calculus (provided the operator is trace class). An obvious choice for the basis functions is the Fourier-basis, i.e.
$`\stackrel{~}{U}_{k,m}`$ $`=`$ $`{\displaystyle _0^1}𝑑q{\displaystyle _0^1}𝑑q^{}\stackrel{~}{U}(q,q^{};N)e^{2\pi i(mq^{}kq)}`$ (30)
$`=`$ $`(1)^m\sqrt{2}{\displaystyle \frac{e^{\pi iN/2}}{2\pi i}}\left[{\displaystyle \frac{e^{\pi i(N/2+m2k)}1}{N/2+m2k}}+{\displaystyle \frac{e^{\pi i(N/2m+2k)}1}{N/2m+2k}}\right].`$ (31)
Similarly one obtains
$`(\stackrel{~}{U}R)_{k,m}`$ $`=`$ $`{\displaystyle _0^1}𝑑q{\displaystyle _0^1}𝑑q^{}\stackrel{~}{U}(q,Rq^{};N)e^{2\pi i(mq^{}kq)}`$ (32)
$`=`$ $`(1)^m\sqrt{2}{\displaystyle \frac{e^{\pi iN/2}}{2\pi i}}\left[{\displaystyle \frac{e^{\pi i(N/2+m+2k)}1}{N/2+m+2k}}+{\displaystyle \frac{e^{\pi i(N/2m2k)}1}{N/2m2k}}\right]`$ (33)
with $`k,m`$ integers. The various terms $`(e^{\pi i(N/2\pm m\pm 2k)}1)/(N/2\pm m\pm 2k)`$ are equal to $`i\pi `$ for $`N/2\pm m\pm 2k=0`$.
A typical eigenvalue spectrum $`\{\mathrm{\Lambda }_i\},i=0,1,2\mathrm{}`$ of the quasiclassical operator $`\stackrel{~}{U}_\pm `$ is shown in Fig. 1 together with the quantum eigenvalues. In the following, the eigenvalues will be ordered with decreasing modulus, i.e. $`|\mathrm{\Lambda }_i|>|\mathrm{\Lambda }_j|`$ for $`i<j`$. The operator $`\stackrel{~}{U}`$ approximates unitarity of the corresponding quantum map $`𝐔`$ quite well; approximately $`N/2`$ eigenvalues in each parity subspace lie near the unit circle and these eigenvalues agree well with quantum results. The modulus of semiclassical eigenvalues becomes exponentially small for $`i>N/2`$ in each subspace and infinitely many eigenvalues ‘disappear’ by spiraling into the origin. This behavior reflects the trace-class property of the operator and makes it possible to truncate the size of the matrix $`\stackrel{~}{𝐔}_\pm `$ in actual calculations. (Choosing dim$`(\stackrel{~}{𝐔}_\pm )=3\times \frac{N}{2}`$ ensures convergence of the first $`N/2`$ eigenvalues in each parity subspace to 4 significant digits).
A measure for the error of the semiclassical approximation compared to quantum calculations is the deviation of the semiclassical eigenvalue with largest modulus, $`\mathrm{\Lambda }_0(N)`$, from the unit circle. The distribution of $`\gamma ^\pm (N):=\mathrm{log}|\mathrm{\Lambda }_0^\pm (N)|`$ for integer values of $`N`$ and both subspaces is shown in Fig. 2. The deviation of semiclassical eigenvalues from the unit - circle decreases with increasing $`N`$ following a $`1/\sqrt{N}`$ behavior. The $`1/\sqrt{N}`$ \- scaling suggest that corrections to individual semiclassical eigenvalues are dominated by diffraction contributions. The $`\mathrm{log}(N)`$ corrections in the traces, cf. (24), become dominant when summing over the periodic orbit contributions or equivalently over the semiclassical eigenvalues and are thus a collective effect of the various contributions.
The deviations from the unit circle for semiclassical eigenvalues with modulus greater than one is systematically larger for eigenvalues corresponding to the $`\stackrel{~}{U}_{}`$ \- operator compared to those from the $`\stackrel{~}{U}_+`$ \- operator. This can qualitatively be understood for small $`N`$ when considering the limit $`N=0`$; the operator $`\stackrel{~}{U}(q,q^{};N=0)=\sqrt{2}\delta \left(q(2q^{}[2q^{}])\right)`$ is the classical Perron-Frobenius operator for the Sawtooth-map up to a factor of $`\sqrt{2}`$. The spectrum $`\stackrel{~}{U}(q,q^{};0)`$ has therefor a single non-zero eigenvalue $`\sqrt{2}`$ in the symmetric subspace which corresponds to $`U_{}`$ in our notation all other eigenvalues are zero. The eigenvalues are smooth functions of the continuous variable $`N`$ which leads to $`|\mathrm{\Lambda }_0^{}|>|\mathrm{\Lambda }_0^+|`$ for $`N`$ smaller one. The fact that this behavior persists for large $`N`$ values , cf. Fig. 1, can certainly not be explained by the argument above but is worth noting as a purely numeric result.
## IV Semiclassical spectral statistics and action correlation functions: numerical results
The existence of a quasiclassical operator for the Baker map makes it possible to write traces as convergent sums over the quasiclassical eigenvalue spectrum, i.e.
$$\text{Tr}\stackrel{~}{U}_\pm ^n(N)=\underset{i=0}{\overset{\mathrm{}}{}}\left(\mathrm{\Lambda }_i^\pm (N)\right)^n.$$
(34)
The violation of unitarity and the existence of semiclassical eigenvalues with modulus greater than one leads to the asymptotic result
$`|\text{Tr}\stackrel{~}{U}_\pm ^n(N)|\mathrm{exp}(n\gamma ^\pm (N))\text{for}n\mathrm{}`$
with $`\gamma ^\pm (N)=\mathrm{log}|\mathrm{\Lambda }_0^\pm (N)|`$, i.e. the traces grow exponentially in the limit $`n\mathrm{}`$ (Keating 1994, Aurich and Sieber 1994). This kind of behavior is connected to the notorious convergence problems for periodic orbit sums of the form
$`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{p}{\overset{(n)}{}}}A_pe^{2\pi iNS_p}`$
which are absolutely convergent only for complex $`N`$ with $`Im(N)>h_t/2`$ and $`h_t`$ is the topological entropy (Eckhard and Aurell 1989). For the Baker map the absolute sum over periodic orbit contributions (6) corresponds to $`_{n=1}^{\mathrm{}}\text{Tr}\stackrel{~}{U}^n(0)`$; this sum diverges like $`2^{n/2}`$ for large $`n`$, i.e. $`h_t=\mathrm{log}2`$ here. A summation over the semiclassical eigenvalues (34) gives a detailed account of the regime of conditional convergence for periodic orbit sums as a function of the real part of $`N`$.
The exponential increase in the traces (34) leads to exponentially increasing terms in a semiclassical approximation of the form factor
$$K_{sc}(\tau ,N)=\frac{1}{N}|\text{Tr}\stackrel{~}{U}_\pm ^{N\tau }|^2=\frac{1}{N}\underset{i,j=0}{\overset{\mathrm{}}{}}\left(\mathrm{\Lambda }_i^\pm \mathrm{\Lambda }_j^\pm ^{}\right)^{N\tau }\mathrm{exp}(2N\tau \gamma ^\pm (N))\text{for}\tau \mathrm{}.$$
(35)
A stationary distribution of $`K_{sc}`$ for fixed $`\tau `$ and $`N\mathrm{}`$ is obtained if $`<\gamma (N)>`$ falls off faster than $`N^1`$. Such a behavior is expected for smooth hyperbolic maps where the semiclassical error is dominated by $`1/N`$ \- corrections due to higher order terms in a stationary phase expansion. This is, however, not the case for the Baker map where diffraction effects of the order $`N^{1/2}`$ dominate, see Fig. 2.
Using the semiclassical traces $`\text{Tr}\stackrel{~}{U}^n`$ to evaluate the two-point correlation function (4) will, on the other hand, always be affected by the exponentially increasing terms in the form factor; a semiclassical calculation of $`R_2(x,N)`$ diverges thus independently of $`N`$ and the behavior of $`<\gamma (N)>`$.
The weighted classical correlation function (7) is a finite sum and thus well defined for fixed topological length $`n`$. From (25) one obtains in the large $`n`$ \- limit
$`\overline{P}(n)`$ $`=`$ $`2^n\text{in the ‘}+\text{’ - subspace}`$
$`\overline{P}(n)`$ $`=`$ $`2^n\text{in the ‘}\text{’ - subspace}`$
and we will discard this ’trivial’ mean part from now on. The large $`n`$ \- asymptotics of the non-trivial correlations in (7) is dominated by the eigenvalue being the largest among all $`\mathrm{\Lambda }_0(N)`$; after defining
$`\gamma _0^\pm :=\mathrm{log}|\mathrm{\Lambda }_0^\pm (N_0)|:=\underset{N\text{I}\text{N}}{sup}(\mathrm{log}|\mathrm{\Lambda }_0^\pm (N)|)`$
where $`N_0`$ is the $`N`$ \- value corresponding to the largest eigenvalue $`\mathrm{\Lambda }_0`$, one obtains
$$P(s,n)e^{2n\gamma _0^\pm }\mathrm{cos}(2\pi N_0s)\mathrm{}\text{for}n\mathrm{}.$$
(36)
It follows from the considerations above that a full description of quantum spectral statistics in terms of classical actions and amplitudes needs to incorporate unitarity into a semiclassical approximation (Keating 1994). I will come back to this point at the end of the section. Apart from that there seems to be little information in the asymptotic behavior of $`K_{sc}(\tau ,N)`$ or $`P(s,n)`$. The rates $`|\gamma ^\pm (N)|>0`$ correspond to ‘semiclassical escape rates’ and contain little information about the classical dynamics or the quantum map. The distribution of these rates is instead of statistical nature due to accumulation of errors in the semiclassical approximation.
The more interesting question in this context is whether or not semiclassical periodic orbit formulae for classical chaotic systems are able to reproduce RMT-statistics quantitatively in the non-asymptotic regime. The form factor obtained from the quasiclassical operator is displayed in Fig. 3 together with the result obtained from the quantum spectrum and the GOE-prediction (applicable for systems with time reversal symmetry). The semiclassical result is indeed capable of reproducing the RMT-result in both subspaces for small $`\tau `$ values; especially the deviation of the RMT-behavior from the classical result $`K(\tau )=2\tau `$ obtained from the sum rule (15) in the limit $`\tau 0`$ is here well reproduced by periodic orbit formulae. This clearly indicates that there are non-trivial correlations between periodic orbit actions for systems with time-reversal symmetry.
Representing periodic orbit formulae by sums over quasiclassical eigenvalues allows one to study the large $`N`$ and $`\tau `$ limit of a semiclassical approximation of the form factor. Exponentially increasing components start to dominate $`K_{sc}(\tau ,N)`$ for $`\tau `$ larger than a semiclassical break-time of the order
$$\tau _b^\pm \frac{1}{\gamma ^\pm (N)N}$$
(37)
The break-time $`\tau _b^\pm `$ is typically a factor 5 larger in the ‘$`+`$’ subspace than in the ‘$``$’ subspace, see Fig. 2.
The action correlation function (7), on the other hand, can be obtained either by sampling action-differences directly (which is limited to $`n<15`$ – 20 due to the exponential increase in the number of orbits) or with the help of the traces $`\text{Tr}\stackrel{~}{U}^n(N)`$. The latter method demands calculation of the spectra of $`\stackrel{~}{U}(N)`$ for integer $`N`$-values and summation of the Fourier-sum in (7) directly. In practice, the sum is truncated at a finite $`N`$-value, $`N_{\text{max}}`$, which corresponds to a smoothing of the original correlation function $`P(s,n)`$ on scales $`1/N_{\text{max}}`$; after rescaling according to Eq. (10) and subtracting $`\overline{P}(n)`$ one obtains
$$P_{scal}(\sigma ,n,N_{\text{max}})=\frac{2}{n^2}\underset{N=1}{\overset{N_{\text{max}}}{}}|\text{Tr}\stackrel{~}{𝐔}^n(N)|^2\mathrm{cos}(2\pi \frac{N}{n}\sigma )_{\mathrm{}}^{\mathrm{}}𝑑\sigma ^{}P_{scal}(\sigma ^{},n)g_\alpha (\sigma \sigma ^{})$$
(38)
with
$$g_\alpha (x)=\frac{\mathrm{sin}(2\pi \alpha x)}{\pi x}\text{and}\alpha =\frac{N_{\text{max}}}{n}.$$
(39)
The smoothed periodic orbit correlation function $`P_{scal}(\sigma ,n,N_{\text{max}})`$ is shown in Fig. 4 for various $`n`$-values and $`N_{\text{max}}=n`$. (Contributions from diagonal terms leading to the $`\delta `$ \- functions in (9) and (10) are subtracted here). Universal periodic orbit correlations in the ‘$`+`$’ subspace is observed up to $`n70`$; these are $`2^{70}`$ distinct periodic orbits, a number inaccessible to pure periodic orbit calculations. One can study even longer orbits and finds exponentially increasing terms dominating the correlation function for $`n`$ values above the transition point $`n70`$. The modulations in the periodic orbit pair density function $`P(\sigma ,n,n)`$ are indeed orders of magnitudes larger than the RMT-correlations for the $`2^{500}`$ periodic orbits of length $`n=500`$, see Fig. 4c
Things look similar in the ‘$``$’ subspace, the deviation from universality occur, however, for much smaller $`n`$-values, i.e. at $`n6`$. The correlations which occur for larger $`n`$ have the simple form $`P(s,n)=e^{\gamma _0^{}n}\mathrm{cos}(2\pi s)`$, see Fig. 4d for $`n=65`$. This oscillatory behavior has already been observed by Sano (1999) when studying periodic orbit correlations in the Baker map without separating the two symmetry subspaces. Sano’s result can now be interpreted in terms of the spectrum of the quasiclassical operator (27); the exponents $`\gamma _0^\pm `$ for the two different subspaces, which determine the asymptotic behavior of $`P(s,n)`$ for large $`n`$, see Eq. (36), are
$`\gamma _0^+=0.022885`$ with $`N_0=14`$
$`\gamma _0^{}=0.148508`$ with $`N_0=1,`$
which can also be deduced from Fig. 2. The exponentially growing terms in the periodic orbit correlation function have thus a seven times larger leading exponent in the ‘$``$’ subspace compared to the ‘$`+`$’ subspace. The correlation function $`P(s,n)`$ in the ‘$``$’ subspace is dominated by the Fourier coefficient $`|\text{Tr}\stackrel{~}{U}^n(1)|^2`$ for $`n6`$ and RMT-type correlation can not develop. The same is obviously true when considering the full Baker map without symmetry-reduction.
Universal correlations coinciding with the RMT-prediction do, however, exist. This is clearly demonstrated in the ‘$`+`$’ subspace for $`n70`$. It is the violation of unitarity in a semiclassical approximation which leads (in general) to exponentially diverging terms which in turn overwhelm the pair-distribution function $`P(s,n)`$ for large $`n`$. The existence of universal periodic orbit correlations is thus strongly linked to the preservation of unitarity in a semiclassical formulation. The absence of RMT - like periodic orbit correlation, as observed in the ‘$``$’ subspace, on the other hand, contains little information about the phenomena of universal spectral statistics but merely reflects the limitations of the approach due to the underlying semiclassical approximations.
A possibility to enforce unitarity onto a semiclassical description has been proposed by Bogomolny and Keating (1996a). The starting point is the quantum staircase function $`𝒩(\theta )=_{i=1}^N\mathrm{\Theta }(\theta \theta _i)`$ with $`\theta _i`$ being the quantum eigenphases and $`\mathrm{\Theta }(x)`$ denotes the Heaviside step-function. A smoothed, semiclassical version of the staircase function may be written as the truncated Fourier-series of $`𝒩(\theta )`$
$$\stackrel{~}{𝒩}(\theta ,N)=\frac{N}{2\pi }\theta \frac{1}{\pi }Im\underset{n=1}{\overset{N}{}}\frac{1}{n}\text{Tr}\stackrel{~}{𝐔}^n(N)e^{in\theta },$$
(40)
and the traces $`\text{Tr}\stackrel{~}{𝐔}^n`$ are expressed as sums over periodic orbit contributions, see Eq. (6). The unitarity condition is implemented by choosing a quantization condition
$$\stackrel{~}{𝒩}(\stackrel{~}{\theta }_n,N)\stackrel{!}{=}n+\frac{1}{2}$$
(41)
with $`\stackrel{~}{\theta }_n`$ being the solutions of (41); the $`\stackrel{~}{\theta }_n`$ represent a semiclassical approximation of the quantum eigenphases $`\theta _n`$. A density of states is defined according to
$$D(\theta ,N)=\underset{i=1}{\overset{N}{}}\delta _{2\pi }(\theta \stackrel{~}{\theta }_i)=\stackrel{~}{d}(\theta ,N)\underset{i=1}{\overset{N}{}}\delta _{2\pi }\left(\stackrel{~}{𝒩}(\theta ,N)n\frac{1}{2}\right)$$
(42)
with $`\stackrel{~}{d}(\theta ,N)=\frac{}{\theta }\stackrel{~}{𝒩}(\theta ,N)`$. The density (42) is again written in terms of periodic orbit contributions only but ‘preserves’ unitarity by construction. Bogomolny and Keating (1996a) started from the expression (42) to derive periodic orbit corrections to the quantum two-point correlation function $`R_2`$ beyond the results obtained from the classical sum-rule (15).
The Perron-Frobenius representation of the semiclassical traces for the Baker map allows one to test the quantization condition (41) for large $`N`$ up to $`N500`$ by calculating the semiclassical traces in (40) directly; the ‘regularised’ spectrum {$`\stackrel{~}{\theta }_i`$} thus obtained is used to construct new semiclassical traces $`\text{Tr}\stackrel{~}{\stackrel{~}{U}}^n=_i\mathrm{exp}(in\stackrel{~}{\theta }_i)`$ which in turn are used to calculate the periodic orbit correlation function (7). Note, however, that the simple relation between periodic orbits of length $`n`$ and the $`n`$-th semiclassical trace gets lost here and the new periodic orbit correlation function $`P(s,n)`$ contains also correlation between orbits and pseudo - orbits (being composites of shorter orbits).
The correlation function obtained from the ‘regularised’ semiclassical data is shown in Fig. 5 together with correlation functions where regularisation has not been applied, cf. also Fig. 4. The functions $`P(\sigma ,n,N_{\text{max}})`$ are shown here with enhanced resolution compared to Fig. 4, i.e. with a cut-off $`N_{\text{max}}=\frac{5}{2}n`$. The periodic orbit correlation function for $`n=70`$ follows the random matrix prediction also with enhanced resolution, but again deviates from the universal behavior for $`n=200`$. The exponentially growing terms are eliminated in the corresponding regularised correlation function using the spectrum obtained from (41) for $`n=200`$. The regularisation procedure unveils the underlying universal correlations in the classical actions, see Fig. 5.
The regularisation process (41) has the disadvantage that the Fourier-components of the new density of states $`D(\theta ,N`$) can no longer be written as closed expressions of a finite number of periodic orbits; they become complicated infinite periodic orbit sums after writing the $`\delta _{2\pi }`$ – function in its Fourier-components (Bogomolny and Keating 1996a). The resulting periodic orbit correlation functions, as e.g. shown in Fig. 5 for $`n=200`$, consist of periodic orbit contributions of orbits and pseudo - orbits of topological length up to and including $`n=200`$ which does not simplify the task of understanding the origin of universal periodic orbit correlations.
## V Conclusions and Outlook
Correlations between actions of periodic orbits of the Baker map up to orbits of topological length $`n=500`$ (corresponding to $`2^{500}`$ different periodic orbits) have been studied with the help of a quasiclassical Perron-Frobenius operator. The spectral form factor can be calculated by purely semiclassical expressions which coincide for small $`\tau `$ with random matrix results beyond the validity of the classical sum-rule, but diverge for $`\tau \mathrm{}`$. Action correlations of periodic orbits have been investigated which show universal non-trivial correlations linked to random matrix theory for short periodic orbits, but depart from the universal behavior for long orbits due to the violation of unitarity in a semiclassical approximation. The transition point from universal to non-universal statistics is distinctively different for the two parity subspaces in the Baker map. It is linked to the magnitude of the semiclassical error made when approximating the quantum density of states by semiclassical periodic orbit formulae and is controlled by the largest deviation of a semiclassical eigenvalue from the unit circle. This behaviour is probably generic for quantum maps. Statistical properties of semiclassical expressions are universal in the non-asymptotic regime, the transition point at which universality breaks down is, however, system dependent and controlled by the semiclassical error. Imposing unitarity on a semiclassical approximation makes it possible to discard exponentially growing terms in the periodic orbit correlation function which in turn uncovers universal correlation even above the transition point, as demonstrated here for the Baker map.
The study presented here establishes for the first time unambiguously the existence of universal periodic orbit correlations in a classical chaotic system whose quantum counterpart shows RMT-eigenvalue statistics; the limitations due to the semiclassical approximations are discussed in detail. This sets a proper framework for studying the origin of these classical correlations. A better understanding of the interplay between classical periodic orbit action correlations and unitarity might shed light on the existence of universality in quantum spectra in general.
Acknowledgments
I would like to thank Stephen Creagh, John Hannay and Jon Keating for stimulating discussions and the Isaac Newton Institute, Cambridge, for the hospitality during the workshop ’Disordered Systems and Quantum Chaos’ where parts of this work have been carried out. |
warning/0001/math0001119.html | ar5iv | text | # Finite groups with the same character tables, Drinfel’d algebras and Galois algebras
## Abstract
We prove that finite groups have the same complex character tables iff the group algebras are twisted forms of each other as Drinfel’d quasi-bialgebras or iff there is non-associative bi-Galois algebra over these groups. The interpretations of class-preserving automorphisms and permutation representations with the same character in terms of Drinfel’d algebras are also given.
1. Introduction. The theory of quasi-Hopf algebras was developed by V.G.Drinfel’d for the description of quantizations of Lie groups and algebras or so-called quantum groups.
Althought the deformational quantization approach which is so useful in the theory of quantum groups can’t be applied for the the case of finite groups, the idea of twisting seems to be very suitible for reformulating of various problems from representation theory of finite groups.
The key observations of this article is that any bijection between character tables of finite groups corresponds to the quasi-isomorphism of the group algebras considered as quasi-Hopf algebras and any two homomorphisms of the group algebras define the same map of character tables iff they are twisted forms.
In particular, we can give the definitions in terms of (quasi-)Hopf algebras of such objects as class-preserving automorphisms, permutation representations with the same character, groups with the same character tables. Namely, any class-preserving automorphism is twisted form of identity maps as homomorphisms of Hopf algebras. Two permutation representations have the same complex character iff the corresponding homomorphisms into symmetric group are twisted forms as homomorphisms of Hopf algebras. Two groups have the same character tables iff their group algebras are twisted forms as quasi-Hopf algebras.
This point of view allows to select the subclass of pairs of groups with the same character tables. This subclass consists of pairs of groups whose group algebas are twisted forms as Hopf algebras.
The notion of quasi-homomorphism of group algebras can be formulated in terms of Galois algebras. Using the calculation of automorphisms of associative Galois algebras we can describe quasi-isomorphisms of group algebras as Hopf algebras. These quasi-isomorphisms correspond to normal abelean $`2`$-subgroups equipped with some non-degenerated bimultiplicative forms.
2. Semirings and character tables. A semiring is a set $`S`$ with a collection of non-negative integers $`\{m_{x_1,x_2}^x,x,x_1,x_2S\}`$ (structural constants) which satisfy the (associativity) condititon
$$m_{x_1,x_2,x_3}^x=\underset{tS}{}m_{x_1,t}^xm_{x_2,x_3}^t=\underset{sS}{}m_{x_1,x_2}^sm_{s,x_3}^x,x,x_1,x_2,x_3S.$$
An element $`e`$ of the semiring $`S`$ is an identity if $`m_{t,e}^s=m_{e,t}^s=\delta _{s,t}`$ for all $`s,tS`$.
A morphism from the semiring $`S`$ to the semiring $`S^{}`$ is a collection $`\{n_t^s,sS,tS^{}\}`$ of non-negative integers which satisfy the following condition:
$$\underset{sS}{}m_{s_1,s_2}^sn_s^t=\underset{t_1,t_2S^{}}{}m_{}^{}{}_{t_1,t_2}{}^{t}n_{s_1}^{t_1}n_{s_2}^{t_2}$$
$`(1)`$
for any $`s_1,s_2S`$ and $`tS^{}`$. A degree map $`d`$ for the semiring $`S`$ is a morphism from $`S`$ to the one-element semiring with identity, e.g. a colection $`\{d(s),sS\}`$ of non-negative integers such that $`d(s_1)d(s_2)=_{sS}m_{s_1,s_2}^sd(s)`$.
The enveloping ring $`A(S)`$ of the semiring $`S`$ is the free $`𝐙`$-module with the basis $`\{[s],sS\}`$ labeled by the elements of $`S`$ and with the multiplications $`[i][j]=_{sS}m_{i,j}^s[s]`$. We will denote by $`A_0(S)`$ the cone of non-negative combinations of basic elements (the cone of non-negative elements).
A morphism of semirings defines a homomorphism of their enveloping rings $`f:A(S)A(S^{})`$ where $`f([s])=_{tS^{}}n_s^t[t]`$.
Denote by $`(,)`$ canonical bilinear form on $`A(S)`$
$$(x,y)=\delta _{x,y},x,yS.$$
The semiring $`S`$ is rigid if there defined an antihomomorphism $`()^{}:SS`$ (conjugation) such that
$$(xy,z)=(y,x^{}z),x,y,zS.$$
Note that conjugation is an anti-endomorphism of the enveloping ring $`A(S)`$:
$$(z,(xy)^{}w)=((xy)z,w)=(x(yz),w)=(yz,x^{}w)=(z,y^{}x^{}w),$$
or $`(xy)^{}=y^{}x^{}`$.
It follows from the definition that the kernel of the conjugation $`()^{}`$ lies in the kernel of the bilinear form $`(,)`$
$$(x,y)=(1,x^{}y)=0,\text{for}xker()^{},yS.$$
Since the bilinear form $`(,)`$ is non-degenerated the conjugation is injective. So it is bijective for the finite semiring $`S`$. In that case the conjugation has a finite order as an automorphism of the finite set $`S`$.
Lemma 1. Let $`d`$ be a degree map for the rigid seniring $`S`$. Then $`\rho =_{sS}d(s^{})sA(S)`$ satisfies to the conditions
$$x\rho =d(x)\rho ,xA(S).$$
Proof. Since $`m_{x,s}^t=(t,xs)=(x^{}t,s)=(t^{}x,s^{})=m_{t^{},x}^s^{}`$ we have
$$x\rho =\underset{sS}{}d(s^{})xs=\underset{s,tS}{}d(s^{})m_{x,s}^tt=\underset{s,tS}{}d(s^{})m_{t^{},x}^s^{}t=\underset{tS}{}d(t^{}x)t=\rho \underset{}{\text{(}}x).$$
$`\mathrm{}`$
Proposition 1(Uniqueness of degree map). Any two degree maps for commutative rigid semiring $`S`$ coincides.
Proof. Let $`d,d^{}`$ be degree maps for $`S`$. Define $`\rho =_{sS}d(s^{})s,\rho ^{}=_{sS}d^{}(s^{})s`$. Then $`d(\rho ^{})\rho =\rho ^{}\rho =d^{}(\rho )\rho ^{}`$, which means $`d=d^{}`$. $`\mathrm{}`$
Since the enveloping ring $`A(S)`$ of rigid semiring $`S`$ is equipped with non-degenerated semi-invariant bilinear form, the enveloping algebra $`A_𝐐(S)=A(S)𝐐`$ over rational numbers $`𝐐`$ is semisimple.
For commutative rigid semiring $`S`$ the enveloping algebra $`A_{\overline{𝐐}}(S)=A(S)\overline{𝐐}`$ over algebraic closure $`\overline{𝐐}`$ of $`𝐐`$ is isomorphic to the algebra of functions on some finite set $`Cl(S)`$ (”conjugacy classes” of $`S`$). The set $`Cl(S)`$ can be identified with the set of maximal ideals of $`A_{\overline{𝐐}}(S)`$, so that the value $`x(c)`$ of $`xA_{\overline{𝐐}}(S)`$ on $`cCl(S)`$ is unique element of $`\overline{𝐐}`$ such that $`xx(c)+m_c`$. Here $`m_c`$ is maximal ideal of $`A_{\overline{𝐐}}(S)`$ corresponding to $`c`$.
For any commutative rigid semiring $`S`$ we can associate a character table $`(s(c))_{sS,cCl(S)}`$, which is $`|S|\times |S|`$-matrix with entries in $`\overline{𝐐}`$.
Proposition 2. The map $`f:A(S)A(S^{})`$ given by the collection $`\{n_t^s,sS,tS^{}\}`$ is a homomorphism of (semi)rings (the collection satisfies to the condition (1)) if and only if there is a map $`f^{}:Cl(S^{})Cl(S)`$ such that $`f(s)(c)=s(f^{}(c))`$ for any $`sS,cCl(S^{})`$.
Proof. The map $`f:A(S)A(S^{})`$ is a ring homomorphism iff $`f_{\overline{𝐐}}:A_{\overline{𝐐}}(S)A_{\overline{𝐐}}(S^{})`$ is a homomorphism of $`\overline{𝐐}`$-algebras. Any homomorphism of algebras of functions corresponds to the map of sets $`f^{}:Cl(S^{})Cl(S)`$. $`\mathrm{}`$
Example 1. The set $`Irr(G)`$ of irreducible characters of the finite group $`G`$ has a natural semiring structure:
$$\chi \psi =\eta Irr(G)m_{\chi ,\psi }^\eta \eta ,\chi ,\psi Irr(G).$$
The map of character tables of the finite groups $`G_1,G_2`$ is a pair consisting of
i) the map of the sets of conjugacy classes $`cl(G_1)cl(G_2),CC^{}`$,
ii) the map from the set of irreducible characters to the semiring of characters $`Irr(G_2)R_0(G_1),\chi \chi ^{}=sum_{\psi Irr(G_1)}n_{\chi ,\psi }\psi `$, where $`n_{\chi ,\psi }0`$
such that $`\chi ^{}(C)=\chi (C^{})`$ for all $`\chi Irr(G_2),Ccl(G_1)`$.
We say that two finite groups $`G_1,G_2`$ have the same character tables if there are one-to-one mappings $`\chi \chi ^{}`$ and $`CC^{}`$ between the sets of irreducible characters and conjugasy classes, respectively, of $`G_1`$ and $`G_2`$, such that $`\chi ^{}(C^{})=\chi (C)`$ for all $`\chi ,C`$. For examples, see .
3. Semisimple monoidal categories. The monoidal category is a category G with a bifunctor
$$:𝒢\times 𝒢𝒢(X,Y)XY$$
which called tensor (or monoidal) product. This functor must be equiped with a functorial collection of isomorphisms (so-called associativity constraint)
$$\phi _{X,Y,Z}:X(YZ)(XY)Z\text{for any}X,Y,Z𝒢$$
which satisfies to the following pentagon axiom:
$$(X\phi _{Y,Z,W})\phi _{X,YZ,W}(\phi _{X,Y,Z}W)=\phi _{X,Y,ZW}\phi _{XY,Z,W}.$$
A quasimonoidal functor between monoidal categories $`𝒢`$ and $``$ is a functor $`F:𝒢`$ , which is equipped with the functorial collection of isomorphisms (the so-called quasimonoidal structure)
$$F_{X,Y}:F(XY)F(X)F(Y)\text{for any}X,Y𝒢.$$
We shal call it monoidal structure if
$$F_{X,YZ}(IF_{Y,Z})=F_{XY,Z}(F_{X,Y}I)$$
for any objects $`X,Y,Z𝒢`$.
The morphism $`\alpha :FG`$ between two monoidal functors $`F,G:𝒢`$ is monoidal if $`F_{X,Y}(\alpha _X\alpha _Y)=\alpha _{xY}G_{X,Y}`$ for any $`X,Y𝒢`$.
Monoidal category structures on $`𝒢`$ differed by the associativity constraint will be called twisted forms of each other.
The structures of monoidal functor for $`F:𝒢`$ will be called twisted forms of each other.
The monoidal category $`𝒢`$ is rigid if it is equipped with the dualization functor, which is a contravariant functor $`()^{}:𝒢𝒢`$ with a collections of morphisms $`\iota :1XX^{}`$ and $`ϵ\upsilon :x^{}X1`$ for any $`X𝒢`$ such that the compositions
$$X\stackrel{I\iota }{}X(X^{}X)\stackrel{\phi }{}(XX^{})X\stackrel{ϵ\upsilon I}{}X,$$
$$X^{}\stackrel{\iota I}{}(X^{}X)X^{}\stackrel{\phi ^1}{}X^{}(XX^{})\stackrel{Iϵ\upsilon }{}X^{}$$
are identical.
Let $`𝒢`$ be a semisimple monoidal $`k`$-linear category over the field algebraically closed $`k`$ with the set $`S`$ of isomorphism classes of simple objects. The collection of dimensions $`m_{y,z}^x=dim_kHom_𝒢(X,YZ)`$ form a semiring structure on the set $`S`$. Here $`X,Y`$ and $`Z`$ are some representatives of the classes $`x,y,zS`$. Note that the enveloping ring of semiring $`S`$ coincides with the Grothendieck ring $`K_0(𝒢)`$ of the category $`𝒢`$. The semiring $`S(𝒢)`$ is rigid for the rigid monoidal category $`𝒢`$.
Proposition 3. Semisimple monoidal categories are twisted forms of each other iff their semirings of simple objects are isomorphic. Isomorphism classes of quasimonoidal functors $`F:𝒢`$ between semisimple monoidal categories are in one-to-one correspondence with the homomorphisms $`S(𝒢)S()`$ of the semirings of simple objects. In particular, monoidal functors $`F,G:𝒢`$ between semisimple monoidal categories are twisted forms of each other iff they induce the same map $`K_0(𝒢)K_0()`$ of the Grothendieck rings.
Proof. The proposition follows from the fact that two fuctors $`F,G:𝒢`$ between semisimple categories are isomorphic iff they induce the same map $`S(𝒢)S()`$ of the semirings of simple objects. $`\mathrm{}`$
4. Drinfel’d algebras. A Drinfel’d algebra or quasi-bialgebra is an algebra $`H`$ together with a homomoprhisms of algebras
$$\mathrm{\Delta }:HHH\text{(coproduct)},\epsilon :Hk\text{(counit)}$$
and an invertible element $`\mathrm{\Phi }H^3`$ (associator) for which
$$(\mathrm{\Delta }I)(\mathrm{\Delta }(h))=\mathrm{\Phi }(I\mathrm{\Delta })(\mathrm{\Delta }(h))\mathrm{\Phi }^1hH\text{(coassociativity)},$$
$$(II\mathrm{\Delta })(\mathrm{\Phi })(\mathrm{\Delta }II)(\mathrm{\Phi })=(1\mathrm{\Phi })(I\mathrm{\Delta }I)(\mathrm{\Phi })(\mathrm{\Phi }1),$$
$$(\epsilon I)\mathrm{\Delta }=(I\epsilon )\mathrm{\Delta }=I.$$
Drinfel’d algebra is a generalization of the well-known notion of bialgebra which corresponds to the case of trivial associator $`\mathrm{\Phi }=1`$.
Drinfel’d algebras structures on the algebra $`H`$ which is differed only by associator will be called twisted forms of each other.
A quasi-homomorphism of quasi-bialgebras $`H_1`$ and $`H_2`$ is pair $`(f,F)`$ consisting of a homomorphism of algebras $`f:H_1H_2`$ and an invertible element $`FH_2^2`$ such that
$$\mathrm{\Delta }(f(h))=F(ff)(\mathrm{\Delta }(h))F^1.$$
It is a homomorphism of quasi-bialgebras if, additionly,
$$(\mathrm{\Delta }I)(F)(F1)(fff)(\mathrm{\Phi }_1)=\mathrm{\Phi }_2(I\mathrm{\Delta })(F)(1F).$$
Two homomorphisms of quasi-bialgebras are twisted forms of each other if they differ only by the invertible element $`F`$. We can define the morphism between two homomorphisms $`(f,H),(g,G):H_1H_2`$ as an element $`cH_2`$ such that $`cf(h)=g(h)c`$ for any $`hH_1`$ and $`\mathrm{\Delta }(c)G=F(cc)`$.
A homomorphism of bialgebras $`H_1,H_2`$ is a homomorphism of algebras $`f:H_1H_2`$ such that $`\mathrm{\Delta }f=(ff)\mathrm{\Delta }.`$
Now we discuss the connection between monoidal categories and quasi-bialgebras. Coproduct allows to define the structure of $`H`$-module on the tensor product $`M_kN`$ of two $`H`$-modules:
$$h(mn)=\mathrm{\Delta }(h)(mn),hH,mM,nN.$$
The associator $`\mathrm{\Phi }`$ defines the associativity constraint
$$\phi :LMNLMN,\phi (lmn)=\mathrm{\Phi }(lmn).$$
Thus the category $`(H)`$ of $`H`$-modules is a monoidal category.
The homomorphism of quasi-bialgebras $`f:H_1H_2`$ defines the monoidal functor
$$f^{}:(H_2)(H_1)$$
with the monoidal structure defined by the invertible element $`FH_2^2`$
$$f_{}^{}{}_{M,N}{}^{}:f^{}(MN)f^{}(M)f^{}(M)f_{}^{}{}_{M,N}{}^{}(mn)=F(mn).$$
The morphisms between homomorphisms $`f,g:H_1H_2`$ of quasi-bialgebras correspond to the monoidal morphisms between monoidal functors $`f^{},g^{}:(H_2)(H_1)`$.
The quasi-Hopf algebra $`H`$ will be called rigid if the monoidal category $`(H)`$ is rigid. The dualization functor for $`(H)`$ corresponds to the antihomomorphism $`S:HH`$ (antipode) with some additional properties (see ). For bialgebra these properties takes a form of the relation
$$\mu (SI)\mathrm{\Delta }=\mu (IS)\mathrm{\Delta }=\epsilon ,$$
where $`\mu :HHH`$ is the multiplication in $`H`$. Bialgebra with an antipode is called Hopf algebra.
Example 2. Group algebra $`k[G]`$ of the group $`G`$. As $`k`$-vector space $`k[G]`$ is spanned by the elements of the group $`G`$. The structure maps have the following forms on the basis:
$$\mathrm{\Delta }(g)=gg,\epsilon (g)=1,S(g)=g^1.$$
The homomorphism of the groups $`f:G_1G_2`$ defines the homomorphism of their group algebras and any homomorphism of bialgebras $`k[G_1]k[G_2]`$ is of this kind. The group algebra provides an example of so-called cocommutative Hopf algebra $`t\mathrm{\Delta }=\mathrm{\Delta }`$. Over the algebraically closed field of characteristic zero group algebras are characterized by this property (Kostant theorem): any cocommutative finite dimensional Hopf algebra is isomorphic to the group algebra.
For semisimple quasi-bialgebra $`H`$ denote by $`S(H)=S((H))`$ the semiring of simple modules. The semiring $`S(H)`$ is rigid for quasi-Hopf algebra $`H`$.
The next proposition is the direct consequence of the definitions and proposition Proposition 3.
Proposition 4. The homomorphisms of quasi-bialgebras $`f_1,f_2:H_1H_2`$ are twisted forms if and only if the monoidal functors $`(f_1)^{},(f_2)^{}`$ are twisted forms. In particular, the homomorphisms of semisimple quasi-bialgebras $`f_1,f_2:H_1H_2`$ induce the same homomorphism $`K_0(f_1),K_0(f_2):K_0(H_2)K_0(H_1)`$ of Grothendieck rings if and only if one is isomorphic to the twisted form of the other.
The generalization of the so-called Tannaka-Krein theory states that we can reconstruct a quasi-bialgebra from the monoidal category $`𝒢`$ and a quasimonoidal functor $`F:𝒢(k)`$ to the category of vector spaces as endomorphisms $`End(F)`$ of the functor $`F`$.
Theorem 1. Finite dimensional semisimple quasi-Hopf algebras $`H_1,H_2`$ are quasi-isomorphic if and only if their semirings of simple objects $`S(H_2),S(H_1)`$ are isomorphic.
Proof. Since twisting does not change the semiring of simple objects we need to prove the if statement. Let $`f^{}:S(H_2)S(H_1)`$ be an isomorphism of semirimgs of simple objects. By Proposition 3 we can costruct a quasi-monoidal functor (equivalence) $`F:(H_2)(H_1)`$ which induces the given homomorphisms $`f^{}`$. By Proposition 1 the composition $`d_1f^{}`$ coincides with $`d_2`$, where $`d_i`$ is a degree map for $`S(H_i)`$. Hence the composition $`F_1F`$ of functor $`F`$ with the forgetful funtor $`F_1:(H_1)(k)`$ is isomorphic to the forgetful funtor $`F_2:(H_2)(k)`$ as quasi-monoidal functor. Using Tannaka-Krein theory we can construct the isomomorphism of quasi-bialgebras $`f:H_1H_2`$ as
$$H_1=End(F_1)End(F_1F)End(F_2)=H_2.$$
$`\mathrm{}`$
Remark 1. The weak version of the theorem Theorem 1 for Hopf algebras was proved in where it was assumed that the isomorphism of (enveloping algebras of) semiring preserves the class of regular representation.
Corollary 1. The finite groups $`G_1,G_2`$ have the same character table if and only if their group algebras are isomorphic as quasi-Hopf algebras, e.g. there is an isomorphism of algebras $`f:k[G_1]k[G_2]`$ and an invertible element $`Fk[G_2]^2`$ such that $`F\mathrm{\Delta }_2(f(x))=(ff)(\mathrm{\Delta }_1(x))`$ for any $`xk[G_1]`$.
If we denote by $`\mathrm{\Delta }_F`$ the twisted by $`F`$ comultiplication on $`k[G_2]`$ $`\mathrm{\Delta }_F(x)=F\mathrm{\Delta }_2(x)F^1`$ then the map $`f`$ would be an isomorphism of Hopf algebras $`(k[G_1],\mathrm{\Delta }_1)`$ and $`(k[G_2],\mathrm{\Delta }_F)`$. The existence of such isomorphism is equivalent to the existence of an isomorphism of groups
$$G_1G(F)=G(k[G_2],\mathrm{\Delta }_F)=\{xk[G_2],F\mathrm{\Delta }_2(x)=(xx)F\}.$$
The cocommutativity of the coproduct $`\mathrm{\Delta }_F`$ is equivalent to the condition
$$t(F)=F$$
$`(2)`$
The coassociativity of the twisted coproduct $`\mathrm{\Delta }_F`$ is equivalent to the equation on $`F`$
$$(1F)(I\mathrm{\Delta })(F)=(F1)(\mathrm{\Delta }I)(F)\mathrm{\Phi },$$
$`(3)`$
where $`\mathrm{\Phi }`$ is some invertible $`G_2`$-invariant element of $`k[G_2]^3`$ (associator). In particular, such $`\mathrm{\Phi }`$ satisfy to the equation
$$(\mathrm{\Phi }1)(I\mathrm{\Delta }I)(\mathrm{\Phi })(1\mathrm{\Phi })=(\mathrm{\Delta }II)(\mathrm{\Phi })(II\mathrm{\Delta })(\mathrm{\Phi }).$$
$`(4)`$
We can replace $`F`$ by the product $`FC`$ for any $`G_2`$-invariant element $`Ck[G_2]^2`$ without changing the twisted coproduct $`\mathrm{\Delta }_{FC}=\mathrm{\Delta }_F`$. The $`G_2`$-invariance of $`C`$ allows to write the associator $`\mathrm{\Phi }^C`$ for the product $`FC`$ as
$$\mathrm{\Phi }^C=(\mathrm{\Delta }I)(C)^1(C1)^1\mathrm{\Phi }(1C)(I\mathrm{\Delta })(C).$$
$`(5)`$
Thus the element $`\mathrm{\Phi }`$ is defined up to the transformations (5).
We can also replace $`F`$ by $`F^c=(cc)F\mathrm{\Delta }(c)^1`$ for invertible $`ck[G_2]`$. Then the corresponding twisted coproducts will be connected by conjugation by $`c`$
$$\mathrm{\Delta }_{F^c}(cxc^1)=(cc)\mathrm{\Delta }_F(x)(cc)^1.$$
The preveous theorem reduces the problem of finding finite groups whose character tables coincide with the character table of $`G`$ to the problem of finding the solutions $`(\mathrm{\Phi },F)`$ to the equations (2), (3), (4) such that the order of the group $`G(F)`$ equals $`|G|`$. If the ground field $`k`$ is algebraically closed of characteristics zero, then we can ommite the condition $`|G(F)|=|G|`$ using Kostant theorem.
In V.Drinfeld suggested the following way of solving the equation (3) for general Hopf algebra. Introduce the new multiplication $`\mu _F`$ on the dual Hopf algebra $`k(G)=k[G]^{}`$
$$\mu _F(lm)(x)=(lm)(F\mathrm{\Delta }(x)),l,mk(G),xk[G].$$
This multiplication will be invariant under the action of the group $`G`$ on $`k(G)`$
$$(gl)(x)=l(xg).$$
By another words, elements of the group $`G`$ act as automorphisms of the algebra $`R_F=(k(G),\mu _F)`$. Moreover, the algebra $`R_F`$ is a so-called Galois $`G`$-algebra. It means, that the natural map of vector spaces
$$R_FR_FHom(k[G],R_F),lm(gg(l)m)$$
is an isomorphism. The group $`G(F)`$ appears as automorphisms group $`Aut_G(R_F)`$ of $`G`$-algebra $`R_F`$. It is not hard to see that if $`|G(F)|=|G|`$, then $`R_F`$ is also Galois $`G(F)`$-algebra, or bi-Galois $`GG(F)`$-algebra.
The algebras $`R_{F_1},R_{F_2}`$ are isomorphic as $`G`$-algebras iff there is an invertible $`ck[G]`$ such that $`F_1=F_2^c`$. This method is mostly applicable for the case of $`\mathrm{\Phi }=1`$, because of the following fact:
the algebra $`R_F`$ is associative iff $`\mathrm{\Phi }=1`$.
In general, it would be only $`\mathrm{\Phi }`$-associative in the following sense:
$$x(yz)=\mathrm{\Phi }(xy)z,x,y,zR,$$
where the product $`\mathrm{\Phi }(xy)z=\mu (\mu I)(\mathrm{\Phi }(xyz))`$ is defined by the action of $`k[G]^3`$ on $`R^3`$.
5. Galois algebras. Here we give brief description of bi-Galois associative algebras. As was explained above they correspond to the isomorphisms of character tables with trivial associators.
Let $`R`$ be an algebra with the action of the group $`G`$ ($`G`$-algebra). The cross product $`RG`$ is a vector space spanned by the elements $`agaR,gG`$ satisfying $`(a+b)g=ag+bg`$. The multiplication is given by the formula
$$(ag)(bf)=ag(b)gfa,bR,gG.$$
A $`G`$-algebra $`R`$ is Galois if the map
$$\theta :RGEnd(R)\theta (ag)(b)=ag(b)$$
is an isomorphism.
A Galois $`G`$-algebra $`R`$ has the following properties:
$`R`$ has no non-trivial $`G`$-invariant twosided ideals,
$`R`$ is semisimple,
$`G`$ acts transitively on the set of maximal twosided ideals of $`R`$.
Let $`S`$ be a stabilizer of some maximal ideal $`M`$ of $`R`$. Then the quotient algebra $`B=R/M`$ is simple Galois $`S`$-algebra and $`R`$ can be identified with the algebra of $`S`$-invariant functions
$$ind_S^G(B)=\{a:GB,a(sg)=s(a(g))sS,gG\}$$
with the $`G`$-action $`(fa)(g)=a(gf)`$.
The $`S`$-algebra $`B=End(V)`$ is Galois iff the multiplier of the projective representation $`SPGL(V)=Aut(B)`$ is a non-degenerated 2-cocycle. We call a 2-cocycle $`\alpha Z^2(G,k^{})`$ non-degenerated if for any $`sS`$ the map from the centralizer
$$C_S(s)k^{}t\alpha (s,t)\alpha (t,s)^1$$
is non-trivial.
Example 3. Let $`A`$ be a finite abelian group. Denote by $`\widehat{A}`$ the dual group $`Hom(A,k^{})`$. The 2-cocycle $`\alpha `$ on $`S=A\widehat{A}`$
$$\alpha ((a,\chi ),(b,\psi ))=\chi (b),a,bA,\chi ,\psi \widehat{A}$$
is non-degenerated.
Describe the automorphisms of Galois algebras.
The set of maximal ideals of $`G`$-Galois algebra $`R`$ can be identified as $`G`$-set with $`G/S`$ where $`S`$ is a stabilizer of some maximal ideal.
The action of automorphisms on maximal ideals defines the homomorphism $`Aut_G(R)N_G(S)/S`$.
The kernel of this homomorphism coincides with $`Aut_S(B)`$. The group of automorphisms of the simple Galois $`S`$-algebra is isomorphic to the character group $`\widehat{S}=Hom_{group}(S,k^{})`$.
The image of this homomorphism coincides with the stabilizer $`St_{N_G(S)/S}(\alpha )`$ of the cohomological class $`\alpha H^2(S,k^{})`$.
Thus we have a short exact sequence
$$\widehat{S}Aut_G(R)St_{N_G(S)/S}(\alpha ).$$
The class of this extension in $`H^2(St_{N_G(S)/S}(\alpha ),\widehat{S})`$ is the image of the class $`\alpha H^2(S,k^{})`$ by
$$d_2^{0,2}:H^0(St_{N_G(S)/S}(\alpha ),H^2(S,k^{}))H^2(St_{N_G(S)/S}(\alpha ),H^1(S,k^{}))$$
the differential of Hochschild-Serre spectral sequence corresponding to the extension $`SN_G(S)N_G(S)/S`$.
We can apply now the outlined description of automorphisms of Galois algebras to the investigation of bi-Galois algebras.
A biGalois $`G_1G_2`$-algebra is an algebra $`R`$ with the commuting actions of the groups $`G_1,G_2`$ such that $`R`$ is both Galois $`G_1`$-algebra and $`G_2`$-algebra.
The $`G_1G_2`$-biGalois algebra corresponds to the following data:
the normal inclusions $`SG_i`$ of the abelian group $`S`$ with the same quotient group $`F`$,
the non-degenerated class $`\alpha H^2(S,k^{})`$ such that
$$d(\alpha )=\gamma _1\gamma _2,$$
where $`\gamma _i`$ is the class of the extension $`SG_iF`$ in $`H^2(F,S)`$ and $`d`$ is the differential of Hochschild-Serre spectral sequence corresponding to the splitting extension of $`F`$ by $`S`$.
Functoriality of the differential $`d`$ allows to reduce consideration to the case of $`p`$-group $`S`$. The diffirential $`d`$ is trivial for abelian $`p`$-groups if $`p2`$.
Example (see, also $`[8]`$). Let $`S`$ be $`2n`$-dimensional vector space over the field $`𝐅_2`$ of two elements. The standard symplectic form $`<,>`$ on $`S`$ defines a non-degenerated 2-cocycle
$$\alpha Z^2(S,k^{}),\alpha (s,t)=(1)^{\beta (s,t)},$$
where $`\beta `$ is bilinear form on $`S`$ such that $`<s,t>=\beta (s,t)\beta (t,s)`$.
Let $`F=Sp_{2n}(2)`$ be the group of automorphisms of the form $`<,>`$. For $`n>1`$ the cohomology group $`H^2(F,\widehat{S})=H^2(Sp_{2n},𝐅_2^{2n})`$ is one dimensional $`𝐅_2`$-vector space generated by the class $`d(\alpha )`$. Thus the affine symplectic group $`AffSp_{2n}(2)`$ and the (unique) non-split extension of $`Sp_{2n}(2)`$ by $`𝐅_2^{2n}`$ have the same character tables.
For $`n=1`$ the cohomology group $`H^2(Sp_{2n},𝐅_2^{2n})`$ is trivial and the pair $`(S,\alpha )`$ defines the automorphism of character table of $`AffSp_2(2)=S_4`$ which doesn’t correspond to any group automorphism. This isomorphism intertwings the characters $`\chi _4`$ and $`\chi _5`$ and the classes $`2A`$ and $`4A`$.
| $`S_4`$ | 1 | 2A | 2B | 3A | 4A |
| --- | --- | --- | --- | --- | --- |
| $`\chi _1`$ | 1 | 1 | 1 | 1 | 1 |
| $`\chi _2`$ | 1 | -1 | 1 | 1 | -1 |
| $`\chi _3`$ | 2 | 0 | 2 | -1 | 0 |
| $`\chi _4`$ | 3 | 1 | -1 | 0 | -1 |
| $`\chi _5`$ | 3 | -1 | -1 | 0 | 1 |
6. Class-preserving automorphisms and permutation representations with the same character. We will call an automorphism $`\varphi Aut(G)`$ of the finite group $`G`$ by class-preserving if $`\varphi `$ preserves all conjugacy classes of $`G`$ $`\varphi (g)g^G`$ for all $`gG`$ (see ).
Proposition 5. For any class-preserving automorphism $`\varphi `$ of the finite group $`G`$ there is an invertible element $`ck[G]`$ such that $`\varphi (g)=cgc^1`$ for any $`gG`$ and $`F=\mathrm{\Delta }(c)^1(cc)`$ is $`G`$-invariant element of $`k[G]^2`$.
Proof. It follows directly from the Proposition 4 and the fact that class-preserving automorphism induces trivial automorphism of the character ring $`R(G)`$. $`\mathrm{}`$
Permutation representation of the group $`G`$ is a homomorphism $`\varphi :GS_n`$ to the group of automorphisms $`S_n=Aut(X)`$ of the finite set of order $`|X|=n`$. Define the character $`\chi _\varphi `$ of the permutation representation $`\varphi :GS_n`$ as $`\chi _\varphi (g)=|\{xX,\varphi (g)(x)=x\}|`$. So $`\chi _\varphi `$ is an image of the natural $`n`$-dimensional character of $`S_n`$ under the homomorphism $`\varphi ^{}:R(S_n)R(G)`$.
Proposition 6. For any two permutation representations $`\varphi ,\psi :GS_n`$ of the finite group $`G`$ there is an invertible element $`ck[S_n]`$ such that $`\varphi (g)=c\psi (g)c^1`$ for any $`gG`$ and $`F=\mathrm{\Delta }(c)^1(cc)`$ is $`\psi (G)`$-invariant element of $`k[G]^2`$.
Proof. It can be proved that the homomorphisms $`\varphi ^{},\psi ^{}:R(S_n)R(G)`$ of character tables corresponded to permutation representations $`\varphi ,\psi :GS_n`$ coincides if coincides the characters $`\chi _\varphi `$, $`\chi _\psi `$. Hence we can apply the Proposition 4. $`\mathrm{}`$
7. Concluding remarks. The cohomological nature of the sets of possible twistings was actively explored in the theory of quantum groups. Nonabelianity of those cohomology is probabily a major difficulty of the theory. In quantum group theory this difficulty was overcome by methods of tangent cohomology which are unapplicable for finite groups. In this case non-abelian cohomology sets of twistings can be abelianized by means of algebraic K-theory. Namely, the maps from the sets of twistings to some Hochschild cohomology of representation ring can be costructed . The detailed description of those maps would be the subject of subsequent paper. |
warning/0001/quant-ph0001113.html | ar5iv | text | # Actual computational time-cost of the Quantum Fourier Transform in a quantum computer using nuclear spins
## Abstract
There have been many proposed methods for the practical implementation of quantum computing. Now quantum computation has reached the turning point from being a conceptual system to becoming a physical one. In this paper, we discuss a practical elementary gate and the actual computational time-cost of the QFT in two physical implementations, namely the bulk spin resonance computer and the Spin Resonance Transistor. We show that almost all universal gates require different times for operation. The actual time-cost of the QFT is $`O(n2^n)`$ for large $`n`$. This differs drastically from the reported cost $`O(n^2)`$ based on ideal quantum computation.
Recent technological development has stimulated proposals for many quantum computers: Bulk Spin Resonance(BSR), trapped ions, cavityQED, Josephson junctions, coupled quantum dots and the Spin Resonance Transistor(SRT). BSR has been implemented experimentally in some organic molecules by the use of conventional nuclear magnetic resonance(NMR) equipment. the SRT is attractive from the viewpoint of it’s integration ability and it’s compatibility with silicon technology.
Complexity analysis classifies quantum algorithms according to a function that describes how a computational cost incurred in solving a problem scales up as larger problems are considered. The computational cost of a quantum algorithm has usually been estimated as the sum of the universal gates required in such ideal mathematical models as the Quantum Turing Machine(QTM) and the quantum circuit. The computational complexity is effective in estimating the essential performance of an algorithm to factor out the variations in performance experienced by different makes of computers with different amounts of Random Access Memory(RAM), swap space, and processor speeds. The above cost is proportional to an actual time-cost in the physical implementation where all quantum operations can be achieved in the same time. However, if the implementation being considered takes a different time for each quantum gate, there is a possibility that the actual time-cost will have a different behavior from the ideal cost. A hardware dependent time-cost is important to experimentalists who research into practical implementations. In this paper, we focus on the actual computational time-cost of the Quantum Fourier Transform (QFT) in two physical implementations that utilise nuclear spins: BSR and the SRT.
For the following discussion, we assumed that the quantum computers being considered are constructed from an array of n quantum registers labeled by $`j(1jn)`$ in one dimension. They calculate $`n`$ qubit data which corresponds to $`N=2^n`$ states. We defined the following matrices. The suffixes on each matrix represent the quantum registers on which it operates. In the definitions of $`C_{j,k}(\theta )`$ and $`D_{j,k}(\theta )`$, the first and the second suffixes represent the target bit and the controlled bit respectively.
$$H_j=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right),R_{yj}(\theta )=\left(\begin{array}{cc}\mathrm{cos}\theta /2& \mathrm{sin}\theta /2\\ \mathrm{sin}\theta /2& \mathrm{cos}\theta /2\end{array}\right)$$
$$R_{zj}(\alpha )=\left(\begin{array}{cc}e^{i\alpha /2}& 0\\ 0& e^{i\alpha /2}\end{array}\right),C_{j,k}(\theta )=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& e^{i\theta }\end{array}\right)$$
$$\mathrm{\Phi }_j(\delta )=\left(\begin{array}{cc}e^{i\delta }& 0\\ 0& e^{i\delta }\end{array}\right),D_{j,k}(\theta )=\left(\begin{array}{cccc}e^{i\theta }& 0& 0& 0\\ 0& e^{i\theta }& 0& 0\\ 0& 0& e^{i\theta }& 0\\ 0& 0& 0& e^{i\theta }\end{array}\right)$$
Firstly, we discuss the practical elementary gate in these implementations. The practical elementary gate is defined as a quantum gate that can be achieved directly by physical phenomena in the implementation being considered.
BSR consists of quantum registers that are the nuclear spins of each atom in an organic molecule. External RF-pulses and magnetic fields control the quantum states and the quantum correlations in the following way.
The nuclear spin under a strong magnetic field $`𝐁=(0,0,B_0)`$ is described by the Hamiltonian $`H=\gamma \mathrm{}B_0I_{jz}`$, where $`\gamma `$ is the gyro-magnetic ratio for the spin, and $`I_{jz}`$ is the z component of the $`j`$th nuclear spin. The time evolution of the system is $`\mathrm{exp}(iHt/\mathrm{})=\mathrm{exp}(i\gamma \mathrm{}B_0tS_{jz})`$ and this then corresponds to the rotation on the z-axis $`R_{zj}(\theta =\gamma \mathrm{}B_0t)`$. The RF-pulse enables the rotation on the other axis. The time evolution shows that the rotation angle $`\theta `$ can be controlled by two factors: the intensity $`B_0`$, and the duration $`t`$, of the external magnetic field. Consequently, there are two control modes. In the intensity control mode, each single qubit rotation takes the same time. In the duration control mode however, each single qubit phase rotation takes a different time, which is proportional to the rotation angle.
The exchange interaction between the $`j`$-th and the $`k`$-th registers is described by the Hamiltonian $`H_{exch.}=J_{jk}I_{jz}I_{kz}`$, where $`J_{jk}`$ is the exchange coupling constant. The time evolution of the system is described by
$$\mathrm{exp}(iJ_{jk}tI_{jz}I_{kz}/\mathrm{})=D_{j,k}(\theta =J_{jk}t/2\mathrm{}).$$
(1)
Sequence(1) shows that the exchange interaction can control the phase rotation angle. Any external fields, however, cannot control the interaction directly, because atoms in a molecule always interact with each other. The refocusing technique can control the angle $`\theta `$ effectively. Only the time duration between the refocusing pulses determines the angle $`\theta `$. The operation $`D_{j,k}(\theta )`$, therefore, takes a time that is proportional to the rotation angle $`\theta `$. An actual operation such as $`D_{j,k}`$ is effective only for adjacent registers because the exchange interaction $`J_{j,k}`$ between non-adjacent registers is very small.
The SRT is composed of quantum registers which are the nuclear spins of arrayed phosphorus ions in silicon, with a globally static magnetic field $`B`$ and an AC magnetic field $`B_{AC}`$. The implementation consists of two gates on the surface: the A-gate above each ion and the J-gate between adjacent ions.
The A-gate controls the strength of hyperfine interactions and the resonance frequency of the nuclear spin beneath it. A globally applied magnetic field $`B_{AC}`$ flips nuclear spins resonant with the field by the same process that occurs in BSR. In this case, only the duration of the resonance determines the rotation angle. It implies that each single qubit phase rotation always takes a different time.
The electron wave function extends over a large distance and makes an effective electron-mediated coupling for two nuclear spins sharing it in semiconductors. The J-gate controls the overlap of the electron wave functions bounded to two adjacent phosphorus atoms, and hence the electron-mediated exchange coupling $`J_{j,k}=J(t)`$ in the time evolution(1) directly. It therefore controls the phase rotation angle in the operation $`D_{j,k}`$. The operation $`D_{j,k}`$ in this implementation also operates only for adjacent registers.
In both implementations, all single qubit phase rotations and controlled phase rotations are practical elementary gates. They can make the quantum XOR in the sequence below, ordered from right to left
$`\sqrt{i}XOR(j,k)`$ (2)
$`=`$ $`R_{yj}({\displaystyle \frac{\pi }{2}})R_{zk}({\displaystyle \frac{\pi }{2}})R_{zj}({\displaystyle \frac{\pi }{2}})D_{j,k}({\displaystyle \frac{\pi }{4}})R_{yj}({\displaystyle \frac{\pi }{2}}).`$ (3)
The suffixes $`j`$ and $`k`$ represent the target bit and the controlled bit respectively. The sequence(3) shows that the quantum XOR depends on single qubit gates in this technique, and is then concerned with the time required for phase rotation.
Barenco and co-workers showed that other universal gates can be constructed by all single qubit gates and the quantum XOR. Each $`n(1)`$ qubit gate required a different number of them, and then a different time for execution. As discussed above, almost all universal gates take a different time in these implementations. It suggests the possibility that the actual time-cost of the quantum algorithm is different from the ideal cost, though the complexity is not affected in the QFT case.
Next we estimate the actual time-cost of the QFT that could be achieved by the above practical elementary gates, by considering the time resolutions of the controlling external fields in these implementations.
The QFT is the transform with base $`N=2^n`$, corresponding to the n qubit defined by
$$|x\frac{1}{\sqrt{q}}\underset{c=0}{\overset{N1}{}}e^{\frac{2\pi icx}{N}}|c.$$
(4)
Shor proposed the algorithm factoring a composite integer by the QFT. The QFT can be constructed by the sequence in the order (from right to left)
$`H_0C_{0,1}(\theta _1)C_{0,2}(\theta _2)\mathrm{}C_{0,n1}(\theta _{n1})H_1\mathrm{}H_{n3}`$ (5)
$`C_{n3,n2}(\theta _1)C_{n3,n1}(\theta _2)H_{n2}C_{n2,n1}(\theta _1)H_{n1}`$ (6)
followed by a bit reversal transformation, where $`\theta _j\pi /2^j`$. The $`n`$ qubit QFT requires $`n(n1)/2`$ controlled phase-shifter $`(C_{j,k}(\theta _{kj})`$s and $`n`$ Hadamard transformation $`H_j`$s, and then $`n(n+1)/2O(n^2)`$. This estimation is based on the complexity analysis method. It coincides with an actual time-cost in the case where all gates required in the sequence(6) are practical elementary ones.
The controlled phase rotation $`C_{j,k}(\theta _{kj})`$ can be achieved by the sequence
$`C_{j,k}(\theta _{kj})`$ $`=`$ $`R_{zk}(\theta _{kj+1})\mathrm{\Phi }_k(\theta _{kj+2})R_{zj}(\theta _{kj+1})`$ (8)
$`XOR(j,k)R_{zj}(\theta _{kj+1})XOR(j,k).`$
The intensity control mode can make all operations $`(C_{j,k}(\theta _{kj}),H_j`$) take almost the same time. Consequently, the actual time-cost coincides with the ideal cost in this mode.
The duration control mode, however, makes each operation $`C_{j,k}(\theta _{kj})`$ require the time $`\tau _{kj}`$ in proportion to the phase rotation angle $`\theta _{kj}`$. The operation $`C_{0,n1}(\theta _{n1})`$ rotates the minimum phase angle $`\theta _{n1}`$ and takes the minimum time $`\tau _{n1}`$ of all rotations in the QFT(6). The range of required phase rotations in the QFT(6) increases with $`2^n`$. For example, the time ratio $`\tau _0/\tau _{n1}`$ approximates to $`2^{100}10^{30}`$ in the 100 qubit QFT.
In general, the time resolution $`t_R`$ controlling the external field is determined by the response time of the system, the delay of the electronic signal and so on. We can only execute in the physical implementations that satisfy the relationship:
$$\tau _0>\tau _1>\mathrm{}>\tau _{n1}t_R.$$
(9)
It is important for the actual time-cost estimation to determine how we set up the unit time $`t_{unit}`$. The unit time $`t_{unit}`$ should be also greater than the time resolution $`t_R`$. The QFT’s in these implementations have various actual time-costs from $`t_{unit}=\tau _0`$ to $`t_{unit}=\tau _{n1}`$.
If we adopt the maximum rotation time $`\tau _0`$ as the unit time $`t_{unit}`$, the actual time-cost is $`O(n)`$, since
$`{\displaystyle \underset{j=0}{\overset{n2}{}}}{\displaystyle \underset{k=j+1}{\overset{n1}{}}}{\displaystyle \frac{\tau _{kj}}{\tau _0}}`$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{n2}{}}}{\displaystyle \underset{k=j+1}{\overset{n1}{}}}{\displaystyle \frac{\theta _{kj}}{\theta _0}}={\displaystyle \underset{j=0}{\overset{n2}{}}}{\displaystyle \underset{k=j+1}{\overset{n1}{}}}2^{jk}`$ (10)
$`=`$ $`n+2^{1n}2O(n).`$ (11)
On the other hand, the condition $`t_{unit}=\tau _{n1}`$ makes the actual time-cost $`O(n2^n)`$, since
$`{\displaystyle \underset{j=0}{\overset{n2}{}}}{\displaystyle \underset{k=j+1}{\overset{n1}{}}}{\displaystyle \frac{\tau _{kj}}{\tau _{n1}}}`$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{n2}{}}}{\displaystyle \underset{k=j+1}{\overset{n1}{}}}{\displaystyle \frac{\theta _{kj}}{\theta _{n1}}}`$ (12)
$`=`$ $`{\displaystyle \underset{j=0}{\overset{n2}{}}}{\displaystyle \underset{k=j+1}{\overset{n1}{}}}2^{n1+jk}`$ (13)
$`=`$ $`(n2)2^{n1}+1O(n2^n).`$ (14)
In this way, the actual time-cost varies from $`O(n)`$ to $`O(n2^n)`$, and it depends on which of these is adopted as the unit time $`t_{unit}`$. The former time-cost however, is not valid for any $`n`$ in the following way.
The former condition, $`t_{unit}=\tau _0`$, means that all phases are always rotated by the external field with constant intensity B for various data of magnitude $`n`$. In this case, the minimum time $`\tau _{n1}`$ decreases exponentially with increasing $`n`$. We cannot rotate the phase $`\theta _j`$ to satisfy the condition $`\tau _j<t_R`$. There exists the upper bound $`n_b`$, satisfying the relationship(9) for the intensity B under consideration. The estimated time-cost(11) is valid for any $`n`$ satisfying $`nn_b`$. It is, however, not valid for any $`n`$ which is greater than $`n_b`$.
The latter condition $`t_{unit}=\tau _{n1}`$ means that the intensity B decreases exponentially with $`n`$ satisfying the relation(9). In this case, we can rotate all phase angles $`\theta _j(0jn1)`$ in the QFT(6) accurately. A particular condition $`t_{unit}=\tau _{n1}=t_R`$ always achieves all phase rotations in the minimum total time, and then yields the best computing performance for each value of $`n`$. The actual time-cost always obeys eq.(14) for any $`n`$.
In this way, the actual time-cost varies from $`O(n)`$ to $`O(n2^n)`$ for any $`n(n_b)`$, and follows only the latter for all other values of $`n`$. These costs are estimated making the assumption that the QFT is always executed accurately.
An approximate QFT, (AQFT) can reduce the arbitrary numbers of the controlled phase shift gates by sacrificing the accuracy. We can select the AQFT with an actual time-cost between $`O(n)`$ and $`O(n2^n)`$, by considering the required accuracy for any $`n`$.
We have discussed the actual time-cost from the viewpoint of the phase rotations in the QFT(6). Almost $`C_{j,k}(\theta _{kj})`$ operations occur for non-adjacent registers in the QFT. We need to construct such non-adjacent gates using adjacent ones so here we estimate the actual time-cost required for constructing such non-adjacent gates from adjacent ones. We must construct any non-adjacent gates by use of the adjacent swap technique. The swap $`S_{j,k}`$ is an operation for exchanging data between two quantum registers simply. It is achieved via the sequence $`S_{j,k}=XOR(k,j)XOR(j,k)XOR(k,j)`$. The non-adjacent two qubit gate $`U_{j,k}`$ is achieved by adjacent swaps and adjacent $`U_{l,l+1}`$ through the following $`kj1`$ sequences.
$`U_{j,k}`$ $`=`$ $`S_{k,k1}S_{k1,k2}\mathrm{}S_{j+3,j+2}S_{j+2,j+1}U_{j,j+1}`$ (16)
$`S_{j+1,j+2}S_{j+2,j+3}\mathrm{}S_{k2,k1}S_{k1,k}`$
$`=`$ $`S_{k,k1}\mathrm{}S_{l+2,l+1}S_{j,j+1}\mathrm{}S_{l1,l}U_{l,l+1}`$ (18)
$`S_{l,l1}\mathrm{}S_{j+1,j}S_{l+1,l+2}\mathrm{}S_{k1,k}(j<l<k)`$
$`=`$ $`S_{j,j+1}S_{j+1,j+2}S_{j+2,j+3}\mathrm{}S_{k2,k1}U_{k1,k}`$ (20)
$`S_{k1,k2}\mathrm{}S_{j+3,j+2}S_{j+2,j+1}S_{j+1,j}.`$
Figure 1 represents the sequence(16). These sequences show that the non-adjacent operation $`U_{j,k}`$ requires both the time for the adjacent operation $`U_{j,j+1}`$ and another time in proportion to $`|kj|`$ for swaps transferring data.
Considering the construction (16$``$20), the QFT(6) requires more $`(n1)n(2n1)/6O(n^3)`$ adjacent swaps for data transfer besides the adjacent controlled phase rotations.
There is, however, a way to reduce some swaps. From sequences(6) and (20), we can obtain the actual quantum circuit shown by Fig.2. In the figure, swaps $`S_{j,j+1}\mathrm{}S_{k1,k}S_{k,k1}\mathrm{}S_{j+1,j}`$ inside each box can be reduced to the identity operation$`I`$. The $`n`$ qubit QFT requires $`(n1)(n2)O(n^2)`$ swaps as a result. It has the same polynomial order as the ideal cost. In this way, there are some cases where the swaps can be reduced due to the symmetry of the algorithm under consideration.
We can conclude that the actual time-cost of the QFT is dominated by the phase rotations and then are $`O(n2^n)`$ in the range $`n>n_b`$.
The range of required phase rotation angles increases exponentially with $`n`$ in the QFT if the accuracy is preserved. In the implementations that are considering, we need to obtain it by controlling the duration or the intensity of the external field. The duration control mode increases the actual time-cost drastically with $`2^n`$. The intensity control mode in BSR takes least time and seems to be the most efficient case for the QFT. This mode, however, leads to another burden on the equipment. The required range of intensity for the applied field increases with $`2^n`$. If the minimum phase rotation is implemented by the field $`B=10^3`$ T, the maximum phase rotation for the 100 qubit QFT requires a field intensity $`B10^{27}`$ T, which is far beyond the current feasible intensity for the magnetic field.
Our results lead to concerns about the feasibility of factoring huge numbers with polynomial order time-costs in the implementations we are considering.
We have shown that almost all universal gates do not expend the same time in BSR and the SRT. This causes the actual time-cost of the QFT to be drastically different from the ideal one. The ideal cost of the QFT is not effective in the range $`n>n_b`$ for practical cases of these implementations if the accuracy is preserved.
We believe that both ideal and actual discussions are important for the development of quantum computer science. The former discussion stimulates study of the characteristic of the algorithm itself, and the latter is also important for the current situation in quantum computation as we move from the conceptual system to the physical one. Our discussion shows the necessity of discussing the practical elementary gate in other proposed quantum computers, and of estimating the actual time-cost for other quantum algorithms. |
warning/0001/astro-ph0001313.html | ar5iv | text | # Metallicity in damped Lyman-𝛼 systems: evolution or bias?
## 1 Introduction
Damped Lyman-$`\alpha `$ systems (DLAs) are high column density (N(HI)$`>`$10<sup>20</sup> cm<sup>-2</sup>) absorbers detected in the optical spectra of quasars up to relatively high redshifts (up to $`z`$5). Their study constitutes a powerful means to investigate the properties of distant galaxies (or of their building blocks). In particular, DLA metal abundances have been widely used in the past few years (e.g. Prochaska and Wolfe 1999, Pettini et al. 1999, Edmunds and Philips 1997, and references therein) in order to probe the nature of DLAs (i.e. galactic discs vs. galaxies of other morphological types) or even to probe the so-called “cosmic chemical evolution” (i.e. the global evolution of gas, metallicity and star formation rate in the Universe). It is not clear, however, whether the observed abundances alone allow to probe the nature or the history of those systems. One reason is possible depletion of metals into dust (e.g. Pei and Fall 1995). In this work we investigate another factor, namely observational biases, along the lines suggested by Boissé et al. (1998). Assuming that (proto)galactic discs constitute a quite plausible model for such systems we show that observational biases may completely alter our interpretation of DLA metal abundances.
## 2 Observational biases of metal abundance estimates in DLAs
In Fig. 1 we present the empirical evidence on which this work is based. Observations of Zn abundances in DLAs obtained by several groups are plotted as a function of HI column density N(HI). This plot is an up-dated version of Fig. 19 of Boissé et al. (1998) and confirms the suggestion of those authors, namely that there seems to be an anti-correlation between the observed Zn abundance and N(HI), independently of redshift $`z`$. As Boissé at al. (1998) notice, this is not to be interpreted as a physical correlation (i.e that high metallicities are characteristic of low column density systems); indeed, the observed abundance gradients in spiral galaxies offer clear evidence that higher abundances are found in the inner disc, where gas column densities are also higher (e.g. Garnett et al. 1997).
The correlation of Fig. 1 is rather to be interpreted in terms of observational biases: no systems with a combination of metallicity and column density outside the shaded region of Fig. 1 are presently detected, even if such systems do exist. The lack of high metallicity and high column density systems should be attributed to extinction effects, since extinction depends on both those factors. The lack of low metallicity and low column density systems is attributed by Boissé et al. (1998) again to observational selection effects, i.e. to the fact that below some level the amount of Zn atoms along the line of sight is insufficient to allow for a proper spectroscopic detection.
How “realistic” are these empirically determined detectability constraints and how can they be quantitatively understood? We postpone a discussion of those matters to Sec. 4. We simply wish to illustrate here the effect of these constraints (assuming they are real) on the detectability of a Milky Way type disc through its metal absorption lines.
In a recent work (Boissier and Prantzos 1999, herefter BP99) we presented a detailed model for the chemical and spectrophotometric evolution of the disc of our Galaxy. The galactic disc is simulated as an ensemble of concentric, independently evolving rings, gradually built up by infall of primordial composition. The chemical evolution of each zone is followed by solving the appropriate set of integro-differential equations. The spectrophotometric evolution is followed in a self-consistent way, i.e. with the star formation rate $`\mathrm{\Psi }(t)`$ and metallicity $`Z(t)`$ of each zone determined by the chemical evolution, and the same stellar Initial Mass Function (from Kroupa et al. 1993). The adopted stellar yields, lifetimes, evolutionary tracks and spectra are metallicity dependent. Dust absorption is included according to the prescriptions of Guiderdoni et al. (1998) and assuming a “sandwich” configuration for the stars and dust layers (Calzetti et al. 1994). The star formation rate (SFR) is locally given by a Schmidt-type law, i.e proportional to some power of the gas surface density $`\mathrm{\Sigma }_g`$ and varies with galactocentric radius $`R`$ as:
$$\mathrm{\Psi }(t,R)=\alpha \mathrm{\Sigma }_g(t,R)^{1.5}V(R)R^1$$
(1)
where $`V(R)`$ is the circular velocity at radius $`R`$. This radial dependence of the SFR is suggested by the theory of star formation induced by density waves in spiral galaxies (e.g. Wyse and Silk 1989).
It turns out that the number of observables reproduced by the model is much larger than the number of free parameters. In particular the model reproduces present day “global” properties (amount of gas, SFR, and supernova rates), as well as the current Milky Way disc luminosities in various wavelength bands and the corresponding radial profiles of gas, stars, SFR and metal abundances; moreover, the adopted inside-out star forming scheme leads to a scalelength of $``$4 kpc in the B-band and $``$2.8 kpc in the K-band, in agreement with observations (see BP99).
At this point it should be noticed that, among all metals observed in DLAs through absorption line measurements, Zn is usually considered to be the most reliable tracer of metallicity, because its abundance is expected to suffer little from depletion into dust (e.g. Pettini et al. 1997) However, from the theoretical point of view, the nucleosynthesis of Zn is not well understood. The only known production site is massive stars, and the most detailed models of this site are those of Woosley and Weaver (1995, WW95), who give the only available models with yields of the various isotopes as a function of initial stellar metallicity. The Zn yields of these models show an unexplained behaviour: the dominant Zn isotope, <sup>64</sup>Zn, is underproduced, while the yields of the next two more important isotopes, <sup>66</sup>Zn and <sup>68</sup>Zn, increase strongly as metallicity increases from 0.1 Z to Z. The reason for this behaviour in the models is not yet well understood. As a result, the \[Zn/Fe\] ratio calculated with the metallicity dependent yields of WW95 and a chemical evolution model shows a rather abrupt and pronounced increase for metallicities \[Fe/H\]$`>1,`$ while it stays nearly constant at lower metallicities (Timmes et al. 1995a, Goswami and Prantzos 2000). This behaviour is not found in the observed pattern of the Zn/Fe ratio: both halo and disk stars of all metallicities show solar Zn/Fe (see Goswami and Prantzos 2000 and references therein), despite the fact that SNIa, the main Fe producers in the disk, do not produce significant Zn amounts.
This situation does not allow us to trust theoretical nucleosynthesis prescriptions for the evolution of Zn, since theory seems unable to reproduce observations in the Galaxy. Therefore we adopt an empirical approach for the nucleosynthesis of Zn, already suggested in Timmes et al. (1995b): based on the observed Zn/Fe pattern in the Galaxy, we assume that Zn traces Fe at all metallicities. Since the history of Fe in the Galaxy is observationally well constrained (at least in the solar neighborhood) and reasonably well reproduced by our models (see BP99), we assume that it can be safely used to trace the history of Zn in other places as well.
In Fig. 1 we present the evolution of the Zn abundance in three representative zones of the Milky Way model (corresponding to the inner disc, the solar neighborhood and the outer disc, respectively) as a function of the corresponding gas column densities of those zones. It should be noticed that in our models we do not distinguish between atomic and molecular gas and the curves in Fig. 1 are drawn by assuming that all the gas is in atomic form (the contribution of He, i.e. 10% by number, is properly removed). This approximation certainly affects the results (i.e. the curve corresponding to the inner disc, where most of the gas is known to be in molecular form, should be shifted to the left by 0.5 on that same scale), but for our illustration purposes the figure is quite appropriate.
Indeed, it is clearly seen that the actual chemical evolution of the disc cannot be revealed by observations, because of the empirically determined constraints discussed previously. At any given epoch (i.e. along the “isochrones” in Fig. 1) only a sub-sample of the disc can be probed. At early times this sub-sample is representative of the inner regions which reach rapidely relatively high column densities and moderate metallicities; at late times these regions develop such high densities and metallicities that they move into the “forbidden” part of the diagram (upper right). At late times then, it is the outer regions that can be probed through the absorption lines.
Obviously, in such conditions it is difficult to derive the real chemical history of the system, since only snapshots of different regions at different times can be available, not of the same regions at different times.
## 3 Metallicity evolution in DLAs
The discussion in the previous section suggests that our picture of DLA chemical evolution, obtained through quasar absorption line measurements, may be seriously biased. Indeed, if taken at face value, Fig. 1 suggests that at early times the lowest metallicities and at late times the highest metallicities are unobservable. Could this bias modify the picture to such an extent as to give the impression of no-evolution at all?
In order to answer quantitatively this question, a working model for the evolution of DLAs is needed. Some authors suggested that DLAs are proto-galactic discs (e.g. Prochaska and Wolfe 1996, Ferrini et al. 1997, Prantzos and Silk 1998) while others interpreted them as low surface brightness galaxies (Jimenez et al., 1998), dwarf irregulars (Matteucci et al. 1997), galactic halos (Valageas et al. 1999) or proto-galactic “building blocks” (Haehnelt et al. 1998, Ledoux et al. 1998). It has also been suggested that DLAs differ substantially from the galaxies that contribute mostly to the observed SFR in the Universe (e.g. Pettini et al. 1999).
We shall adopt in this work the hypothesis that DLAs are galactic discs. We shall use a detailed model we developed recently for the chemical and spectrophotometric evolution of galactic discs (Boissier and Prantzos 2000). It is essentially the Milky Way model presented in Sec. 2, extended to disc galaxies through “scaling properties” derived by Mo, Mao and White (1998) in the framework of the Cold Dark Matter (CDM) scenario for galaxy formation. In our simplified version of this scenario disc profiles can be expressed in terms of only two parameters: maximal rotational velocity $`V_C`$ (measuring the mass of the halo and, by assuming a constant halo/disc mass ratio, also the mass of the disc) and spin parameter $`\lambda `$ (measuring the specific angular momentum of the halo). If all discs are assumed to start forming their stars at the same time (but not at the same rate!), the profile of a given disc can be expressed in terms of the one of our Galaxy. We constructed a grid of 25 models caracterised by $`V_C`$ = 80, 150, 220, 290, 360 km/s and $`\lambda /\lambda _{MW}`$ = 1/3, 1, 5/3, 7/3, 3, respectively, where $`\lambda _{MW}`$ is the spin parametre of the Milky Way. Increasing values of $`V_C`$ correspond to more massive discs and increasing values of $`\lambda `$ to more extended ones. The SFR is calculated from Eq. 1, with the appropriate velocity profile $`V(R)`$. Notice that the efficiency $`\alpha `$ is not a free parameter, since it is the same as in the Milky Way model.
It turns out that this simple model reproduces fairly well most of the main properties of present days discs (Boissier and Prantzos 2000): disc sizes and central surface brigthness, Tully-Fisher relations in various wavelength bands, colour-colour and colour-magnitude relations, gas fractions vs. magnitudes and colours, abundances vs. local and integrated properties, as well as integrated spectra for different galactic rotational velocities. Moreover, as shown in Prantzos and Boissier (2000), it also reproduces the observed abundance gradients in disc galaxies.
In Fig. 2 we present the results of our models for the Zn evolution in DLAs (assumed to be galactic discs) as a function of redshift $`z`$. For clarity, only the evolution of the inner and outer disc is presented in each panel (thin curves, corresponding to zones located at 0.5 and 5.5 scalelengths from the centre, respectively). Star formation is assumed to start at redshift $`z`$=6 for all discs, but any value of $`z>`$4-5 would produce results similar to those displayed here. The resulting evolution is not very different from that calculated in e.g. Prantzos and Silk (1998) for the Milky Way, or in Ferrini et al. (1997) with multi-zone disc models. Notice that Malaney and Chaboyer (1996), Timmes et al. (1995b), Matteucci et al. (1997), Edmunds and Philipps (1997) and Lindner et al. (1999) have studied metallicity evolution in DLAs with one-zone models, differing by the star formation timescales or by the time of the beginning of star formation.
It can be clearly seen in Fig. 2 that between redshifts $`z`$3 and $`z`$1 there is substantial metallicity evolution in all galactic zones, typically an increase by a factor $``$10. Such an increase is certainly not observed in the available data, which is also displayed on each one of the panels in Fig. 2.
The region enclosed within thick solid curves in Fig. 2 is obtained by application of the “empirical constraints” of Fig. 1, i.e. by excluding all regions with a combination of metallicity and column density $`F`$(Zn,N(HI)) = \[Zn/H\] + log(N(HI)) such that $`F<`$18.8 or $`F>`$21. The resulting observational picture is now completely different from the real one: no sizeable evolution in metallicity is observed (except, perhaps, in the lowest redshift range, where metallicities are somewhat higher than average).
It should be noticed that the empirical “filter” has been applied to our models by assuming that discs are seen “face-on” and that all the gas is in the form of HI. These simplifying assumptions have opposite effects on the derived column density along the line of sight: adopting a different inclination would increase the column density of our disc models; taking into account that part of the gas is in the form of H<sub>2</sub>, would decrease it. The former factor can be treated statistically, but not the later. Taking all other uncertainties into account (i.e. possible variations in the H<sub>2</sub>/HI ratio with metallicity, see Combes 1999) we think that Fig. 2 gives a rather good first approximation to the real situation.
The results of Fig.2 are summarised in the upper panel of Fig. 3, wher we plot in the same diagram all the “filtered” zones of our models as a function of redshift (shaded area) and compare them to observations. The “no-evolution” picture is even more clearly seen now, especially in the $`z`$1-3 redshift range. A firm prediction of these models is that the Zn abundances of DLAs at higher redshifts, in the range $`z`$3-5, will be not too different from those already detected at $`z`$1-3.
Assuming that this first approximation is correct (i.e. that DLAS are indeed galactic discs) it is interesting to calculate the most probable metallicity values expected at a given redshift. Even if our basic assumptions are correct, this is by no means a trivial task, since it implies the knowledge of the appropriate statistical factors as a function of the redshift. For illustration purposes we adopt here the following simplified assumptions:
i) the distribution function of discs in the velocity space $`V_C`$ is time-independent and given by the expression suggested in Gonzalez et al. 1999 (in the following we simplify, for clarity, the notation $`V_C`$ to $`V`$, unless if explicitly stated otherwise):
$$F_V(V)dV=\stackrel{~}{\mathrm{\Psi }}_{}\left(\frac{V}{V_{}}\right)^\beta exp\left[\left(\frac{V}{V_{}}\right)^n\right]\frac{dV}{V_{}}.$$
(2)
The parameters $`\mathrm{\Psi }_{},V_{},\beta `$ and $`n`$ are determined in Gonzalez et al. (1999) on the basis of observed Tully-Fisher relationships and luminosity (Schechter-type) functions. We adopt here the set of parameters of their Table 4 (fifth row, LCRS-Courteau data) corresponding to the velocity interval covered by our models. Our results would not be affected much by the choice of another velocity function, since the form of $`F_V`$ always favours discs of low $`V`$.
ii) the distribution function in spin parameter $`\lambda `$-space is time-independent and given by:
$$F_\lambda (\lambda )d\lambda =\frac{1}{\sqrt{2\pi }\sigma _\lambda }exp\left[\frac{ln^2(\lambda /\overline{\lambda })}{2\sigma _\lambda ^2}\right]\frac{d\lambda }{\lambda }$$
(3)
with $`\overline{\lambda }`$=0.05 and $`\sigma _\lambda `$=0.5 (obtained by numerical simulations, see e.g. Mo, Mao and White, 1998) and $`\lambda _{MW}`$=0.06 for the Milky Way disc (Sommer-Larsen, private communication). The $`\lambda `$-function favours moderately “compact” discs (those with $`\lambda `$0.04-0.05).
iii) the probability that a line of sight to a QSO intercepts a disc in the radius interval $`[R,R+dR]`$ is proportional to the geometrical cross-section $`F_RdR=2\pi RdR`$, favouring the detection of the outer regions of the larger discs.
iv) the distributions $`F_V`$, $`F_\lambda `$ and $`F_R`$ are independent.
Applying the joint probability function $`F(V,\lambda ,R)=F_VF_\lambda F_R`$ to our models, we obtain the mean metallicities $`<[Zn/H]>`$ shown in Fig. 3. The two curves in the upper panel are obtained by:
$$\left[\frac{Zn}{H}\right]=\frac{_\lambda _V_0^{R_L}F(\lambda ,V,R)\mathrm{\Phi }(R)[\frac{Zn}{H}(R)]𝑑R𝑑V𝑑\lambda }{_\lambda _V_0^{R_L}F(\lambda ,V,R)\mathrm{\Phi }(R)𝑑R𝑑V𝑑\lambda }$$
(4)
where $`R_L`$ is the radius of the largest disc in our models. The mean value over the whole disc (i.e. without applying the empirical “filter”) corresponds to $`\mathrm{\Phi }(R)`$=1 in all zones and is given by the thin curve in Fig. 3. In that case, $`<[Zn/H]>_U`$ ($`U`$ for “Unfiltered”) increases by a factor $``$20 between redshifts $`z`$=3 and $`z`$=1 and is clearly below all observational data (because the outer, low metallicity, regions of the discs are favoured in that case). This shows the importance of properly taking into account various statistical factors, something that has not been done in previous studies of DLAs with multi-zone disc models (e.g. in Ferrini et al. 1997 and Prantzos and Silk 1998).
The mean value $`<[Zn/H]>_F`$ over the “filtered’ disc zones (i.e. those in the shaded aerea of Fig. 3) is obtained with $`\mathrm{\Phi }(R)`$=1 in those zones and $`\mathrm{\Phi }(R)`$=0 outside them. It is shown by the thick curve in Fig. 3. This $`<[Zn/H]>_F`$ value is in the lower range of the “filtered” values, again because of the geometrical factor $`F_R`$. $`<[Zn/H]>_F`$ increases by a factor of $``$2 between $`z`$=3 and $`z`$=1, an increase which is compatible with the observations.
Pettini et al. (1999) have also estimated the column density weighted average of the \[Zn/H\] values in DLAs, by binning their data in 5 redshift bins (the last bin, at $`z>3`$, being in fact an upper limit only). Again, no substantial evolution is seen in the data (lower panel in Fig. 3). We also calculated the corresponding average metallicity in our models, by folding with the gas column densities $`N_H(R)`$ of the “filtered” zones (assuming that all but 10 % by number - corresponding to He - is in the form of atomic hydrogen):
$$\left[\frac{Zn}{H}\right]_{FW}=\frac{_\lambda _V_0^{R_L}F(\lambda ,V,R)\mathrm{\Phi }[\frac{Zn}{H}]N_H𝑑R𝑑V𝑑\lambda }{_\lambda _V_0^{R_L}F(\lambda ,V,R)\mathrm{\Phi }N_H𝑑R𝑑V𝑑\lambda }$$
(5)
where $`\mathrm{\Phi }`$, \[Zn/H\] and $`N_H`$ depend on radius $`R`$. The resulting evolution is shown in the lower panel of Fig. 3. It can be seen that the weighted mean metallicity $`<[Zn/H]>_{FW}`$ of the “filtered” zones of our models evolves in a way which is certainly compatible with the data. A clear prediction of the model is that at low redshifts ($`z<`$1) there should be as much evolution as in the $`z`$=1-3 range (i.e. a factor of $``$2-3 increase in the weighted mean metallicity of DLAs in both cases).
In summary, by using a self-consistent model of galactic chemical evolution (i.e. reproducing in detail the properties of local galaxies) and incorporating “reasonable” statistics and appropriate empirical constraints, we have shown unambiguously that only a small degree of evolution should be expected for the obeserved mean metallicity of DLAs; according to our models, these systems may well be galactic discs.
## 4 Origin of biases
It is interesting to see how the empirically determined upper limit in the Zn vs N(HI) plane may be intrepreted in terms of extinction. We shall assume here that our model discs constitute gaseous screens, seen face-on and reducing the intensity of the light of background quasars. The corresponding extinction is calculated in the rest-frame of the absorber by the formula:
$$A_\lambda (z)=(A_\lambda /A_V)_{}(A_V/N_H)_{}N_H(Z/Z_{})^{1.6}$$
(6)
where $`(A_\lambda /A_V)_{}`$ is the local normalised extinction curve (Natta and Panagia,1984), $`(A_V/N_H)_{}`$ is the extinction in the V-band in the solar neighborhood (from Bohlin et al., 1978) and the exponent 1.6 in the metallicity term is introduced in order to reproduce the observed extinction curves in the Small and Large Magellanic Clouds (Guiderdoni and Rocca-Volmerange, 1987). The wavelength $`\lambda `$ depends on the redshift $`z`$ of the absorber. Since the background QSO light is observed in the visible ($`\lambda _{V,obs}`$=0.55 $`\mu `$), we have for the corresponding wavelength $`\lambda _{ABS}`$ of the absorber
$$\lambda _{ABS}(z)=\frac{\lambda _{V,obs}}{1+z}$$
(7)
For illustration purposes, we calculated the extinction at 3 different redshifts $`z`$=1,2,3 (i.e. in the absorber’s rest-frame $`\lambda _{ABS}`$= 0.275, 0.183 and 0.137 $`\mu `$, respectively) during the evolution of all the zones of our models. We plot it in Fig. 4 as a function of $`F`$=\[Zn/H\]+log(N(HI)) (upper panel), while in the lower panel we plot the corresponding fraction of background light filtered through the screen. It can be seen that for values of $`F>`$21 (i.e. above the empirically determined limit in Fig. 1) extinction increases rapidly, reaching 1 mag at $`z`$=1, 1.5 mag at $`z`$=2 and 2 mag at $`z`$=3). The background intensity drops below 40% of its initial value. Our results substantiate the claim of Boissé et al. (1998) that extinction is biasing the interpretation of metallicity abundance determinations in DLAs.
As for the lower value of the “empirical” constraint, it can be understood as follows: Taking into account that, by definition, DLAs correspond to column densities log(N(HI)) $`>`$20, the lower limit $`F<`$18.8 corresponds to \[Zn/H\]$`<`$-1, i.e. to less than 5 10<sup>11</sup> atoms of Zn per cm<sup>2</sup> along the line of sight (adopting a solar ratio (Zn/H)$`{}_{}{}^{}`$4.5 10<sup>-8</sup> by number, e.g. Anders and Grevesse 1989). Current surveys of DLAs typically reach N(Zn)$`>`$10<sup>12</sup> cm<sup>-2</sup> (Pettini et al. 1997), which explains, perhaps, why the lower left part of Fig. 1 is void.
In any case, it should be interesting to detect the stellar population responsible for the chemical enrichment of the gas in DLAs. In Fig. 5 we plot our model B-V values of that population, which span a narrow range 0.2$`<`$B-V$`<`$0.4 for redshifts 1$`<z<`$3. In fact, these values are lower limits, since our model discs are assumed to be seen face-on. Thus, one of the firm predictions of our model applied to DLAs is that the underlying stellar population should be somewhat redder than the curve in Fig. 4.
## 5 Summary
The observationally determined constraint 18.8 $`<`$ \[Zn/H\]+log(N(HI)) $`<`$ 21 (Boissé et al. 1998) has profound implications for the intrepretation of metal abundances detected in DLAs. Assuming that DLAs are galactic discs, and using detailed (and successful) models for disc evolution, we show that current observations cannot probe the true evolution of those systems: a no-evolution picture, compatible with available data, emerges when the empirical constraints are taken into account. We calculate average metallicities in the galactic zones of our models that are “filetered” by the empirical constraints, by taking into account appropriate statistical factors. We find that the resulting column density weighted average metallicity shows a small increase at low redshifts and is compatible with currently available data. These findings suggest that DLAs may well be galactic discs, as argued e.g. in Prochaska and Wolfe, 1999 (for a different view, see Pettini et al., 1999).
We also show quantitatively how extinction may be indeed responsible for the non-detection of metal rich DLAs of high column density, as suggested by Boissé et al. (1998); according to those authors, such DLAs may be one day detected in samples drawn from the observation of fainter QSOs, independently of the redshift. Moreover, metal poor DLAs of low column density should also be detected with more sensitive instruments. Finally, if our interpretation of currently observed DLAs as galactic discs is correct, we expect that the underlying stellar populations should have B-V$`>`$0.2 (by a small amount) in the redshift range 1$`<z<`$3.
Acknowledgements: We are grateful to Patrick Boissé and Patrick Petitjean for useful discussions and comments on this work. |
warning/0001/cond-mat0001277.html | ar5iv | text | # Phase Transition in the Takayasu Model with Desorption
## I Introduction
Many systems in nature, ranging from reaction-diffusion systems to fluctuating interfaces, exhibit nonequilibrium steady states with a wide variety of phases. Of particular interest are the self-organized critical systems where different physical quantities have power law distributions in the steady state over a wide region of the parameter space. Self-organized criticality has been studied in a variety of model systems ranging from sandpiles to earthquakes. A particularly simple lattice model due to Takayasu, where masses diffuse, aggregate upon contact and adsorb unit masses from outside at a constant rate, was shown to exhibit self-organized criticality: the steady-state mass distribution has a nontrivial power law decay for large mass in all dimensions. This model initially generated a lot of attention as it was a simple exactly solvable model of self organized criticality with close connections to other solvable models such as the Scheidegger river model, the voter model and the directed abelian sandpile model. Recently there has been a renewed interest in this model as simple variants of the Takayasu model have been found useful in modelling the dynamics of a variety of systems including force fluctuations in granular systems such as bead packs, river networks, voting systems, wealth distributions, inelastic collisions in granular gases, the generalized Hammersley process, particle systems in one dimension and various generalized mass transport models.
In the Takayasu model, each site of a lattice has a nonnegative mass variable. Starting from an initial random distribution of masses, each mass hops to a nearest neighbour site (chosen at random) and aggregates with the mass there with rate $`1`$. In addition, a unit mass is adsorbed at every site with rate $`q`$. While the first move tends to create big masses via diffusion and aggregation, the second move replenishes the lower end of the mass spectrum. At large time $`t`$, the mass distribution at any site has the scaling behaviour, $`P(m,t)m^\tau f(m/t^\delta )`$ with $`\delta =1/(2\tau )`$. The interesting point is that even though the average mass per site increases linearly with time, $`mt`$, the mass distribution $`P(m,t)`$ approaches a time-independent power law distribution $`P(m)m^\tau `$ for $`t\mathrm{}`$ (since $`f(0)O(1)`$) for any nonzero adsorption rate $`q`$. The exponent $`\tau `$ is independent of $`q`$ and is known exactly, $`\tau =4/3`$ in one dimension and $`\tau =3/2`$ within mean field theory.
The steady-state mass distribution in the Takayasu model has the same power law decay for any nonzero adsorption rate $`q`$ and does not undergo any phase transition. In this paper we show that if we introduce an additional process of desorption of unit masses with rate $`p`$ in the Takayasu model (we call this the In-out model), a rich steady-state phase diagram emerges in the $`pq`$ plane. In particular we show that the system undergoes a nonequilibrium phase transition across a phase boundary $`p_c(q)`$. Nonequilibrium phase transitions between steady states have been studied extensively in recent years in a variety of systems. Examples include, amongst others, active-absorbing phase transitions in reaction diffusion systems, roughening transitions in fluctuating interfaces, phase transitions in driven diffusive lattice gas models, wetting transitions in solid-on-solid models, boundary driven transitions in one dimensional asymmetric exclusion processes and Bose-Einstein like condensation in models of aggregation and fragmentation. However we show below that the mechanism of the phase transition and the associated critical properties in the In-out model are very different from those of other models mentioned above.
There are quite a few physical systems where our In-out model may find applications. In nature there exist a variety of systems ranging from colloids to polymer gels where the basic constituents of the system diffuse and coalesce upon contact. For example, in a polymer gel the basic constituents are polymers of different sizes which diffuse in a solution and when an $`m`$-mer comes in contact with an $`n`$-mer, they aggregate to form an $`(m+n)`$-mer. Similarly during the growth of a thin film on an amorphous substrate (such as Bismuth on Carbon), clusters or islands of atoms can diffuse as a whole and when two of them come closer they coalesce. A zeroth order approach to model the dynamics of these systems would be to replace each cluster by a point particle (ignoring its shape) carrying a positive mass which indicates its size or number of atoms. When two particles coalesce their masses add up. In addition many of these systems are open in the sense that they can exchange basic units with the adjoining environment. For example, during the growth of a film on a substrate, single adatoms may adsorb on the substrate from the outside vapour or desorb into the vapour from the substrate. We attempt to incorporate these processes on a lattice in the In-out model and show that even this simple model has a very rich steady-state phase diagram. We had introduced this model in an earlier publication and some results were briefly mentioned. In this paper we present a detailed analysis of the model.
The paper is organized as follows. In section II, we define the In-out model precisely and summarize the different phases and the transitions between them. In section III, we solve the model analytically within mean field theory. In section IV, we present the numerical results in one dimension and discuss a scaling theory which provides scaling relations between different critical exponents. We conclude in section V with a summary and a discussion of open questions.
## II The In-out Model
For simplicity we define the In-out model on a one-dimensional lattice with periodic boundary conditions; generalizations to higher dimensions are straightforward. Each site of a lattice has a nonnegative mass variable $`m_i0`$. Initially each $`m_i`$ is chosen independently from any well defined distribution. The dynamics proceeds as follows. A site $`i`$ is chosen at random and then one of the following events can occur:
1. Adsorption: With probability $`q/(p+q+1)`$, a single particle is adsorbed at site $`i`$; thus $`m_im_i+1`$.
2. Desorption With probability $`p/(p+q+1)`$, a single particle detaches from and leaves site $`i`$; thus $`m_im_i1`$ provided $`m_i1`$.
3. Diffusion and Aggregation: With probability $`1/(p+q+1)`$, the mass $`m_i`$ at site $`i`$ moves to a nearest neighbour site \[either $`(i1)`$ or $`(i+1)`$\] chosen at random. If it moves to a site which already has some particles, then the total mass just adds up; thus $`m_i0`$ and $`m_{i\pm 1}m_{i\pm 1}+m_i`$.
If the site chosen is empty, only adsorption can occur with probability $`q/(p+q+1)`$.
The In-out model has only two parameters $`p`$ and $`q`$. The question we would like to address is: For given $`p`$ and $`q`$, what is the single site mass distribution $`P(m)`$ in the steady state? Note that in the limit $`p=0`$ (i.e., without the desorption process) our model reduces to the Takayasu model mentioned in the introduction.
While the Takayasu model (zero desorption, $`p=0`$) does not have a phase transition in the steady state, we find that introducing a nonzero desorption rate $`p`$ induces a rich steady state behaviour in the $`pq`$ plane. In fact we find that there is a critical line $`p_c(q)`$ in the $`pq`$ plane. For fixed $`q`$, if we increase $`p`$ from $`0`$, we find that for all $`p<p_c(q)`$, the steady state mass distribution has the same large $`m`$ behaviour as in the Takayasu case, i.e., $`P(m)m^\tau `$ where the exponent $`\tau `$ is the Takayasu exponent and is independent of $`q`$. Thus the Takayasu phase is stable upto $`p_c`$. For $`p=p_c(q)`$, we find the steady state mass distribution still decays algebraically for large $`m`$, $`P(m)m^{\tau _c}`$ but with a new critical exponent $`\tau _c`$ which is bigger than the Takayasu exponent $`\tau `$. For $`p>p_c(q)`$, we find that $`P(m)\mathrm{exp}(m/m^{})`$ for large $`m`$ where $`m^{}`$ is a characteristic mass that diverges if one approaches $`p_c(q)`$ from the $`p>p_c(q)`$ side. The critical exponent $`\tau _c`$ is the same everywhere on the critical line $`p_c(q)`$. This phase transition occurs in all spatial dimensions including $`d=1`$.
It is easy to write down an exact evolution equation for the mean mass $`m(t)`$ per site. Since the diffusion and aggregation move does not change the total mass, the only contributions to the time evolution of $`m`$ come from the adsorption and desorption processes. It is then evident that,
$$\frac{dm}{dt}=qps(t)$$
(1)
where $`s(t)`$ is the probability that a site is occupied by a nonzero mass. The first term on the right hand side of the above equation clearly indicates the increase in mass per site due to the adsorption of unit mass. The second term quantifies the loss in mass per site due to the desorption of unit mass taking into acount the fact that the desorption can take place from a site only if the site is occupied by a nonzero mass. Let us fix $`p`$ and vary $`q`$. As long as $`q<q_c(p)`$, it turns out that in the long time limit $`t\mathrm{}`$, the two terms on the right hand side of the above equation cancel each other and the occupation density reaches the asymptotic time independent value, $`s=q/p`$. This indicates that the average mass per site, $`m`$ becomes a constant in the long time limit. In fact, we show below that in this phase, the steady state mass distribution $`P(m)\mathrm{exp}(m/m^{})`$ for large $`m`$ with a finite first moment $`m`$. We call this phase the “Exponential ” phase. However if $`q>q_c(p)`$, the occupation density reaches a steady state value $`s`$ such that $`s<q/p`$. As a result in the long time limit, the second term on the right hand side of Eq. (1) fails to cancel the first term and the mean mass per site $`m(t)`$ increases linearly with time, $`m(qps)t`$. However, as we show below, even though the mean mass diverges in this phase as $`t\mathrm{}`$, the mass distribution reaches a steady state, $`P(m)m^\tau `$ for large $`m`$ where $`\tau `$ is the Takayasu exponent (which is always less than $`2`$ so that the mean mass diverges). Hence we call this entire phase the “Takayasu” phase.
## III Mean Field Theory
We first analyze the model exactly within the mean field approximation, ignoring correlations in the occupancy of adjacent sites. In that case we can directly write down equations for $`P(m,t)`$, the probability that any site has a mass $`m`$ at time $`t`$.
$`{\displaystyle \frac{dP(m,t)}{dt}}=`$ $``$ $`(1+p+q+s)P(m,t)+pP(m+1,t)`$ (2)
$`+`$ $`qP(m1,t)+PP;m1`$ (3)
$`{\displaystyle \frac{dP(0,t)}{dt}}=`$ $``$ $`(q+s)P(0,t)+pP(1,t)+s(t).`$ (4)
Here $`PP=_{m^{}=1}^mP(m^{},t)P(mm^{},t)`$ is a convolution term that describes the coalescence of two masses and $`s(t)=_{m=1}P(m,t)`$ denotes the probability that a site is occupied by a nonzero mass.
The above equations enumerate the possible ways in which the mass at a site might change. The first term in Eq. (3) is the “loss” term that accounts for the probability that a mass $`m`$ might move as a whole or desorb or adsorb a unit mass, or a mass from the neighbouring site might move to the site in consideration. In this last case, the probability of occupation of the neighbouring site, $`s(t)`$ multiplies $`P(m,t)`$ within the mean-field approximation where one neglects the spatial correlations in the occupation probabilities of neighbouring sites. The remaining three terms in Eq. (3) are the “gain” terms enumerating the number of ways that a site with mass $`m^{}m`$ can gain or lose mass to make the final mass $`m`$. The second equation Eq. (4) is a similar enumeration of the possibilities for loss and gain of empty sites.
To solve the equations, we compute the generating function, $`Q(z,t)=_{m=1}^{\mathrm{}}P(m,t)z^m`$ from Eq. (3) and set $`Q/t=0`$ in the steady state. We also need to use Eq. (4) to write $`P(1,t)`$ in terms of $`s(t)`$. This gives us a quadratic equation for $`Q`$ in the steady state. Choosing the root that corresponds to $`Q(z=0)=0`$, we find
$$2zQ(z)=p(z1)+qz(1z)+2sz\sqrt{(z1)\mathrm{\Delta }(z)}.$$
(5)
where
$`\mathrm{\Delta }(z)=`$ $`p^2(z1)+q^2z^2(z1)2pqz(z1)`$ (6)
$``$ $`4qz(zsp/q).`$ (7)
Note that the occupation density $`s`$ in the above expression of $`Q(z)`$ is yet to be determined. The steady state mass distribution $`P(m)`$ can be formally obtained from $`Q(z)`$ in Eq. (5) by evaluating the Cauchy integral,
$$P(m)=\frac{1}{2\pi i}_{C_o}\frac{Q(z)}{z^{m+1}}𝑑z$$
(8)
over a contour $`C_o`$ encircling the origin in the complex plane. This expression for $`P(m)`$ however will contain the yet to be determined unknown quantity $`s`$. In fact, determining $`s`$ is the most nontrivial part of the mean field calculation as we show below.
In order to extract the large-$`m`$ behaviour of $`P(m)`$ from Eq. (8), one needs to deform the contour $`C_o`$ so that it goes around the branch cut singularities of the function $`Q(z)`$. From Eq. (5), it is evident that such singularities occur at $`z=1`$ and also at the roots of $`\mathrm{\Delta }(z)=0`$ where $`\mathrm{\Delta }(z)`$ is given by Eq. (7). Since $`\mathrm{\Delta }(z)`$ is a cubic polynomial in $`z`$, it has three roots $`z_1`$, $`z_2`$ and $`z_3`$, each of which can be determined in terms of the unknown quantity $`s`$.
We now analyse the large $`m`$ behaviour of $`P(m)`$ in different regions of the $`pq`$ plane. Let us fix the value of $`p`$ and increase $`q`$ from $`0`$. A similar analysis can be carried out for fixed $`q`$ as a function of $`p`$. As we increase $`q`$ from $`0`$, we encounter the following three regimes,
(i) For small $`q`$ (with a fixed $`p`$), we first assume that the mean mass $`m`$ reaches a time-independent constant as $`t\mathrm{}`$. This assumption will be justified a posteriori. Then from Eq. (1), it follows that the occupation density also reaches a steady state value, $`s=q/p`$. Substituting this in the expression for $`\mathrm{\Delta }(z)`$ in Eq. (7), the three roots of $`\mathrm{\Delta }(z)=0`$ are $`z_1=1`$ and $`z_{2,3}=(p+22\sqrt{p+1})/q`$. Then from Eq. (5), it follows that the only branch cut singularities of $`Q(z)`$ are at $`z_2`$ and $`z_3`$ with $`z_3>z_2>1`$ for small $`q`$. Therefore the branch cut at $`z_2`$ essentially controls the large-$`m`$ behaviour of $`P(m)`$ when the contour in Eq. (8) is deformed and by analysing the integral around this cut we find that for large $`m`$,
$$P(m)exp(m/m^{})/m^{3/2}$$
(9)
with $`m^{}=1/lnz_2`$. Since $`P(m)`$ decays exponentially in this phase, $`m`$ is also finite and nonzero thus justifying the assumption made in the begining. In this phase the unknown function $`s`$ is therefore exactly given as $`s=q/p`$. Note however that this analysis is valid as long as $`z_2>1`$ and the characteristic mass $`m^{}`$ diverges as $`z_2`$ approaches $`1`$ from above.
(ii) As the value of $`q`$ is increased (for fixed $`p`$) the roots $`z_2`$ and $`z_3`$ decrease, until at a critical value $`q_c(p)`$, the value of $`z_2`$ just reaches unity. The double root ($`z_1`$ and $`z_2`$) at $`z=1`$ of $`\mathrm{\Delta }(z)`$ then leads to a branch cut singularity of order $`3/2`$ in $`Q(z)`$ in Eq. (5), which in turn implies
$$P(m)m^{5/2}.$$
(10)
This power law decay characterizes the critical point and the condition $`z_2=1`$ determines the locus of the critical line in the $`pq`$ plane,
$$q_c(p)=p+22\sqrt{p+1}.$$
(11)
The value of $`s`$ is given exactly by $`s=q_c/p`$.
(iii) As $`q`$ is increased further ($`q>q_c(p)`$) for fixed $`p`$, the mean mass per site $`m`$ does not reach a time-independent value in the steady state, but increases indefinitely with time. Consequently we cannot use the relation $`s=q/p`$ anymore. However, $`P(m)`$ reaches a time-independent distribution. So the question is what is the selection principle that determines the unknown function $`s`$ in this regime?
Note that at $`q=q_c(p)`$, the two roots $`z_1`$ and $`z_2`$ of $`\mathrm{\Delta }(z)=0`$ coincided, $`z_1=z_2=1`$ and $`z_3>1`$. As $`q`$ increases further, since we do not know what $`s`$ is a priori, the exact locations of the three roots of $`\mathrm{\Delta }(z)=0`$ in the complex plane are also unknown. However since $`\mathrm{\Delta }(z)`$ is a polynomial with real coefficients, if $`z`$ is a root of $`\mathrm{\Delta }(z)=0`$, so must be its complex conjugate $`z^{}`$. Thus as $`q`$ increases beyond $`q_c(p)`$, there are two possibilities. The first possibility is that all the three roots of $`\mathrm{\Delta }(z)=0`$ are real and distinct. But in that case, as $`q`$ increases slightly beyond $`q_c(p)`$, at least one of them must become less than $`1`$. This however would lead to an exponential growth of $`P(m)`$ for large $`m`$ and hence is ruled out. The second and only possibility is that one of the three roots must be real while the other two are complex conjugates of each other, i.e., $`\mathrm{\Delta }(z)=(zz_c)(zz_{c}^{}{}_{}{}^{})(zz_3)`$ where $`z_3`$ is real and $`z_c`$ in general is complex with its real part less than $`1`$. However, if the imaginary part of $`z_c`$ is nonzero, this again can be shown to lead to an exponential divergence of $`P(m)`$ for large $`m`$. Therefore, we are led to the conclusion that $`z_c`$ must be real and thus $`\mathrm{\Delta }(z)=0`$ must have double roots at $`z_c`$, i.e., $`\mathrm{\Delta }(z)=(zz_c)^2(zz_3)`$ with $`z_c`$ real. In summary we conclude that for $`q>q_c(p)`$, $`z_3`$ remains $`>1`$ and the two roots $`z_1=z_2=z_c`$ continues to be coincident and real but the common value $`z_c`$ decreases below $`1`$ as $`q`$ increases beyond $`q>q_c(p)`$. This nontrivial ‘root sticking’ condition determines the unknown quantity $`s`$ for $`q>q_c(p)`$. This condition of double roots can be easily implemented by demanding the two conditions, $`\mathrm{\Delta }(z_c)=0`$ and $`\mathrm{\Delta }^{}(z_c)=0`$ where $`\mathrm{\Delta }^{}=d\mathrm{\Delta }(z)/dz`$. Also using the relation $`\mathrm{\Delta }(z)=(zz_c)^2(zz_3)`$ in Eq. (5), we find that the lowest branch cut singularity of $`Q(z)`$ is at $`z=1`$. This order $`1/2`$ singularity then leads to the following asymptotic behaviour of $`P(m)`$,
$$P(m)m^{3/2}.$$
(12)
Thus this entire phase, $`q>q_c(p)`$ is characterized by the same power-law decay of $`P(m)`$ as in the mean field Takayasu model which, as mentioned earlier, corresponds to the zero-desorption ($`p=0`$) limit of our model.
As mentioned above, the ‘root-sticking’ condition also determines quite non-trivially the occupation density $`s`$ for $`q>q_c(p)`$ for fixed $`p`$. To determine $`s`$ explicitly for $`q>q_c(p)`$ using this condition, let us fix $`p=1`$ for simplicity even though the calculation can be carried out for any arbitrary $`p`$. From Eq. (11), we find $`q_c=32\sqrt{2}`$ for $`p=1`$. We first substitute the expression for $`\mathrm{\Delta }(z)`$ from Eq. (7) in the ‘root-sticking’ conditions, $`\mathrm{\Delta }(z_c)=0`$ and $`\mathrm{\Delta }^{}(z_c)=0`$. We then eliminate $`z_c`$ from these equations and find $`s(q)`$ for $`q>32\sqrt{2}`$ as the only positive root of the cubic equation
$`16s^3`$ $``$ $`(q^212q+24)s^2(q^3+5q^2+57q+15)s`$ (13)
$`+`$ $`(q^3+5q^2+39q2)=0.`$ (14)
Thus we can determine the unknown quantity $`s`$ exactly everywhere in the $`pq`$ plane. In Fig. 1, we plot the function $`s(q)`$ for fixed $`p=1`$. For $`qq_c=32\sqrt{2}`$ we have $`s(q)=q`$ and for $`q>q_c=32\sqrt{2}`$, $`s(q)`$ is given by the real positive root of the cubic equation in Eq. (14).
Note that for fixed $`p`$, if $`q>q_c(p)`$, the steady state value of $`s(q)`$ (as determined from the ‘root sticking’ conditions) is less than $`q/p`$ and hence from Eq. (1), we find that the mean mass per site increases linearly with time, $`mvt`$ for large $`t`$. If one interprets the mass profile as the height of an interface (see Section V) then for $`q<q_c(p)`$, the average “height” of the interface becomes a constant as $`t\mathrm{}`$, while for $`q>q_c(p)`$, the average “height” $`m`$ increases linearly with velocity $`v`$. The ‘velocity’ $`v`$ defined more precisely as $`v=lim_t\mathrm{}\frac{m}{t}`$ is $`0`$ for $`q<q_c(p)`$ and nonzero for $`q>q_c(p)`$. For $`q`$ slightly bigger than $`q_c(p)`$, $`v[qq_c(p)]^y`$ where $`y`$ is a critical exponent independent of $`p`$. For example for $`p=1`$, we find from Eq. (14), $`v(qq_c(1))^2/(6\sqrt{2}8)`$ indicating that $`y=2`$ within mean field theory.
## IV Numerical results in one dimension and scaling theory
Having completed the mean field calculations we now turn to one dimension. While the Takayasu model ($`p=0`$) is exactly solvable in $`d=1`$, the same technique unfortunately does not work for $`p>0`$. Hence for nonzero $`p`$, we had to resort to numerical simulations in $`d=1`$. The qualitative predictions of mean field theory namely the existence of a power-law (Takayasu) phase ($`P(m)m^{\tau _T}`$) and a phase with exponential mass distribution, with a different critical behaviour at the transition ($`P(m)m^{\tau _c}`$), are found to hold in 1-d as well. Figure 2 shows the results of numerical simulations for the phase diagram along with the mean-field prediction (Eq. 11) and Figure 3 displays the numerical data for the decay of the mass distribution $`P(m)`$ in the two phases and at the transition point. The values obtained, $`\tau =4/3`$ (same as the exactly solvable $`p=0`$ case) and $`\tau _c1.833`$, are quite different from their mean-field values $`\tau =3/2`$ and $`\tau _c=5/2`$, reflecting the effects of correlations between masses at different sites.
If the phase boundary is crossed by increasing $`q`$ for fixed $`p`$, the Takayasu phase is obtained for $`q>q_c`$. As a function of the small deviation $`\stackrel{~}{q}qq_c`$ and large time $`t`$, the mass distribution $`P(m,\stackrel{~}{q},t)`$ is expected to display a scaling form for large $`m`$,
$$P(m,\stackrel{~}{q},t)\frac{1}{m^{\tau _c}}Y(m\stackrel{~}{q}^\varphi ,\frac{m}{t^\alpha })$$
(15)
in terms of three unknown exponents $`\varphi `$, $`\alpha `$, $`\tau _c`$ and the two variable scaling function $`Y`$. All other exponents then can be related to these three exponents via scaling relations. We give some examples below.
(a) Consider $`\stackrel{~}{q}>0`$ and $`t\mathrm{}`$ limit. Then $`P(m,\stackrel{~}{q})\frac{1}{m^{\tau _c}}Y(m\stackrel{~}{q}^\varphi ,0)`$. But we know that for $`\stackrel{~}{q}>0`$, in the steady state, $`P(m,\stackrel{~}{q})m^\tau `$ where $`\tau `$ is the known Takayasu exponent. This forces the scaling function $`Y(x,0)x^\gamma `$ for large $`x`$ such that $`P(m,\stackrel{~}{q})\stackrel{~}{q}^{\varphi \gamma }/m^{\tau _c\gamma }`$, indicating $`\gamma =\tau _c\tau `$.
(b) Consider again $`\stackrel{~}{q}>0`$ and finite but large $`t`$. The mean mass per site, $`m=mP(m,\stackrel{~}{q},t)𝑑m\stackrel{~}{q}^yt`$ where $`y`$ is the ‘velocity’ exponent. Using the scaling form of $`P`$, we find, $`y=\varphi [1\alpha (2\tau _c)]/\alpha `$.
(c) Next we consider the critical point, $`\stackrel{~}{q}=0`$. Using the scaling form, we find that the mean mass, $`mt^\zeta `$ for large $`t`$ where $`\zeta =\alpha (2\tau _c)`$ provided $`\tau _c<2`$. If $`\tau _c>2`$ (as in mean field theory), $`\zeta =0`$. Also, the root mean square mass fluctuations at the critical point, $`\sigma =\sqrt{(mm)^2}\sqrt{m^2}t^\beta `$ for large $`t`$ with $`\beta =\alpha (3\tau _c)/2`$. Note that for large $`t`$, $`m^2>>m^2`$ indicating that fluctuations grow faster than the mean as time increases.
Within mean field theory, by analysing $`P(m)`$ explicitly for $`\stackrel{~}{q}>0`$, we find $`P(m,\stackrel{~}{q})\stackrel{~}{q}/m^{3/2}`$ and also $`\tau _c=5/2`$. From (a) above, this immediately gives, $`\gamma \varphi =1`$ and $`\gamma =1`$ indicating $`\varphi =1`$. Also, we had shown before that the velocity exponent $`y=2`$ exactly within mean field theory. Using $`y=2`$, $`\tau _c=5/2`$ and $`\varphi =1`$ in (b) of the previous paragraph, we get $`\alpha =2/3`$. Since $`\tau _c=5/2>2`$, we note from (c) that $`\zeta =0`$. Also we find the fluctuation exponent $`\beta =1/6`$ from the scaling relation in (c). Thus within mean field theory, we find
$$P(m,\stackrel{~}{q},t)\frac{1}{m^{5/2}}Y(m\stackrel{~}{q},\frac{m}{t^{2/3}}).$$
(16)
We have determined the corresponding exponents in $`d=1`$ numerically. The critical exponent $`\tau _c1.83`$ has already been mentioned (see Fig 3). In Fig 4, we plot the velocity $`v`$ as a function of $`q`$ for fixed $`p=2.35`$. The velocity is zero for $`qq_c1.0`$ and increases as a power law, $`v(qq_c)^y`$ for small $`\stackrel{~}{q}=(qq_c)`$. We find $`y1.47`$. Note that since $`q_c`$ is not known exactly, this exponent is difficult to determine numerically and is subject to large error bars. We also find that at the critical point $`q_c1`$, the mean mass grows as, $`mt^\zeta `$ with $`\zeta 0.12`$. Note the difference from the mean field theory where $`m`$ does not grow with time at the critical point ($`\zeta =0`$). To measure the fluctuations at the critical point, we performed finite-size studies of the time-dependent ‘width’ $`W^2(t,L)=_{i=1}^L(m_i<m>)^2/L`$ at the critical point, where $`L`$ is the system size. This is expected to obey the scaling form $`Wt^\beta Z(t/L^z)`$; the value of $`z`$ is expected to be 2 as the movement of masses is diffusive. Figure 5 shows the scaling plot of $`W/t^\beta `$ versus $`t/L^z`$ for four different system sizes $`L=16`$, $`32`$, $`64`$ and $`128`$ at the critical point $`p2.35`$ for fixed $`q=1`$. We fix $`z=2`$ and find the best collapse of data for $`\beta 0.358`$. These exponent estimates in $`d=1`$ are consistent with the scaling relations mentioned in (a)-(c).
## V Summary and discussion
In this paper we have studied a simple lattice model where masses diffuse and aggregate with rate $`1`$, unit masses adsorb at any lattice site with rate $`q`$ and unit masses desorb from a site (provided the site is occupied by a mass) with rate $`p`$. For $`p=0`$ (without the desorption process), our model reduces to the well studied Takayasu model where the steady-state single site mass distribution has a power law decay, $`P(m)m^\tau `$ for large $`m`$ for any nonzero $`q`$. We show that varying the desorption rate $`p`$ induces a nonequilibrium phase transition at a critical value $`p=p_c(q)`$. For $`p<p_c(q)`$, $`P(m)m^\tau `$ for large $`m`$ as in the Takayasu ($`p=0`$) case. For $`p=p_c(q)`$, $`P(m)m^{\tau _c}`$ where $`\tau _c`$ is a new exponent and $`P(m)\mathrm{exp}(m/m^{})`$ for $`p>p_c(q)`$. We have solved the model analytically within the mean field theory and calculated all the mean field exponents exactly. In one dimension, we have computed the exponents numerically. We have also presented a general scaling theory.
Our model generalizes the Takayasu model and exhibits a nontrivial phase transition. There was an earlier generalization of the model where instead of carrying positive masses, the diffusing particles carried charges $`Q`$ of either sign while a random charge $`I`$, drawn from an arbitrary distribution, was added with rate $`q`$ to a lattice site. In this “charge” model, the steady-state single site charge distribution $`P(Q)`$ was found to have a power law tail (as in the mass case), $`P(Q)Q^\tau `$ for large positive $`Q`$ when the mean charge injected was positive, $`I>0`$ whereas for $`I=0`$, $`P(Q)Q^{\tau _1}`$ for large positive $`Q`$. It was shown that the exponent $`\tau _1=5/3`$ in $`d=1`$ and $`\tau _1=2`$ within mean field theory. Though this change of exponent at a critical value $`I=0`$ is similar to that in our model qualitatively, the exponent $`\tau _c`$ of the In-out model is very different from that of the “charge” model. This difference can be traced back to the mass positivity constraint in the In-out model, i.e., the desorption of an unit mass can take place from a lattice site only if the site has a nonzero mass.
In the In-out model, the total mass is not conserved due to the moves involving adsorption and desorption of unit mass. It is interesting to ask what would happen if the desorption of a unit mass from a site were followed by adsorption at a neighbouring site, so that the total mass would be conserved in every move. This was investigated using a lattice model and earlier, within a rate equation approach . In this conserved-mass model too there is a phase transition, but of a different character. It was found that there is an exponential phase (at high desorption-adsorption rate), separated by a critical line from a phase with a power-law mass distribution $`P(m)m^{\tau _{\mathrm{conserved}}}`$. This distribution coexists with an infinite aggregate which accommodates a finite fraction of the total mass — a real space analog of Bose-Einstein condensation . The exponent $`\tau _{\mathrm{conserved}}`$ was found to be 5/2 within mean field theoryand $`2.33`$ in 1-d , and the same exponent was found to describe $`P(m)`$ at the critical point. Evidently, the lack of mass conservation in the In-out model is responsible for the absence of the infinite aggregate in its high $`q`$ phase, as well as the change in the power to $`\tau `$ in the Takayasu phase and $`\tau _c`$ at criticality.
Another interesting difference between the conserved and the In-out model is the effect of a preferred direction for the motion of masses (a mass at site $`i`$ hops with a higher probability to $`i1`$ than $`i+1`$). We have checked that such a bias does not change the critical exponents of the In-out model. However, for the conserved mass model, the bias in direction changes the value of the exponents at the transition and in the aggregate phase .
The phase transition in the In-out model has some interesting implications for nonequilibrium wetting transitions if we interpret the configuration of masses as an interface profile regarding $`m_i`$ as a local height variable. Although the dynamics of the mass profile in our In-out model is not physical when interpreted as interface dynamics, nevertheless the phase transition in our model can be qualitatively interpreted as a nonequilibrium wetting transition of the interface. In the In-out model, the fact that the mass at each site is necessarily non-negative translates into the restriction that the mass profile is always above a wall at a fixed height (in our case $`0`$). The presence of this constraint is the key factor for the wetting transition. At fixed $`p`$, as we increase $`q`$, the mean height $`m`$ does not grow with time as long as $`q<q_c`$. This is our ‘Exponential’ phase where the mass or the interface profile is bound to the substrate at zero height. This phase is also ‘smooth’ as the mean square height fluctuation does not grow with system size. For $`q>q_c`$, the interface unbinds from the substrate and the mean height $`mvt`$ grows linearly with time with a velocity $`v`$. This phenomenon is similar to ‘wetting’ or ‘depinning’ of interfaces in general. In this ‘wet’ phase (Takayasu phase in our model), the interface is rough. Unlike recently studied models of nonequilibrium wetting, where the interface in the growing phase is self-affine, our model describes a much rougher interface for $`q>q_c`$. At the transition $`q=q_c`$, though, the interface is self-affine with a roughness exponent $`\chi =z\beta 0.7`$.
There are various open questions that remain to be settled. In this paper, we have only studied the phase transition in the steady-state single site mass distribution function. It would be very interesting to study the spatial correlations between masses at different sites and to track the behaviour of mass-mass correlation function as one crosses the phase boundary in the $`pq`$ plane.
Also in this paper we have only studied the simplest model where the rates of adsorption, desorption and hopping are constants and independent of particle mass. An important question is whether this phase transition would persist for general mass-dependent rates. In earlier work , a model with aggregation, adsorption and desorption was studied, but no transition to a power-law phase was found; the difference is traceable to the fact that in that model, the rate of removal of mass is proportional to the mass, unlike the unit-mass desorption process considered in the In-out model. It is therefore highly desirable to identify the class of models with mass-dependent rates where the phase transition described here will persist. |
warning/0001/cond-mat0001359.html | ar5iv | text | # Semi-classical magnetoresistance in weakly modulated magnetic fields
## Abstract
The semi-classical conductance of a two-dimensional electron gas is calculated in the presence of a one-dimensional modulated magnetic field with zero average. In the limit of small magnetic field amplitudes (B) the contribution of the magnetic modulation to the magnetoresistance increases as $`B^{3/2}`$ in the diffusive limit, while the increase is linear in $`B`$ in the ballistic regime. Temperature does not influence the power law behavior but it decreases the prefactor of this functional behavior.
In recent years there has been an increased interest in hybrid systems which promise to increase the functionality of present day semiconductor devices. One example of such type of systems are those in which semiconductors and magnetic materials are combined where the magnetic material provides a local magnetic field which influences locally the motion of the electrons in the semiconductor. The latter is usually a heterostructure which contains a two-dimensional electron gas (2DEG). The 2DEG acts as a detector measuring the magnetic state of the magnetic material. Previously, the coupling between such a non homogeneous magnetic field and the 2DEG was demonstrated in the case the magnetic field ($`B`$) was directed perpendicular to the 2D plane. In this case one has a modulation of the $`B`$-field on top of a homogeneous background field and the influence of the $`B`$-field modulation on the 2DEG is relatively weak .
When the magnetic field is directed parallel to the 2DEG the magnetic material is magnetized parallel to the 2D plane which leads to fringing fields near the edge of the magnetic material having a non zero magnetic field component perpendicular to the 2D layer. Those fields form a magnetic barrier for the electron motion in the 2D plane. Because now there is no background perpendicular magnetic field for the 2DEG the influence on the resistance of such magnetic barriers is much more pronounced and large increases in the magnetoresistance have been found .
In the present work we investigate the magnetotransport in weak modulated magnetic fields for which the average $`B`$-field is zero. For the case the typical magnetic energy is much smaller than the Fermi energy a semi-classical analysis is applicable. We find that in the diffusive regime the correction to the magnetoresistance exhibits a non analytical behavior in the limit of small magnetic field amplitudes which differs from the behaviour in the ballistic regime.
We consider electrons moving in a two-dimensional (2D) $`xy`$-plane. The magnetic field, directed along the $`z`$-direction, is periodic along the $`x`$-direction $`\stackrel{}{B}=B(0,0,b(x/l_0))`$ with period $`l_0`$, where $`b(x)`$ is a periodic ($`b(x+1)=b(x)`$) dimensionless function describing the magnetic field modulation with zero average value.
In a semi-classical analysis the electron motion in a magnetic field is described by the following Hamiltonian (or its energy):
$$\epsilon =\frac{1}{2m}\left\{p_x^2+\left[p_y\frac{eBl_0}{c}a(x/l_0)\right]^2\right\},$$
(1)
where $`m`$ is the electron effective mass, $`\stackrel{}{p}=(p_x,p_y)`$ is the electron canonical momentum, and the dimensionless periodic function
$$a(x)=^x𝑑x^{}b(x^{}),$$
(2)
characterizes the vector potential $`\stackrel{}{A}=Bl_0(0,a(x/l_0),0)`$ which is taken in the Landau gauge. The quantum energy spectrum of electrons in modulated magnetic fields was studied in Refs. .
We restrict our analysis to the case where the electron transition through a single period is ballistic, i. e. the mean free path $`l=v_F\tau l_0`$ ($`v_F`$ is the electron Fermi velocity, and $`\tau `$ the relaxation time) and the motion in the sample is diffusive, i.e. the mean free path is smaller than the size of the sample. ($`lL_x,L_y`$). For diffusive transport and in the limit of small magnetic fields ($`\omega _c\tau 1`$) the expression for the average conductivity tensor is given by the following integral over the electron phase space $`(x,p_x,p_y)`$
$`\sigma _{ij}={\displaystyle \frac{e^2}{(2\pi \mathrm{})^2L_x}}{\displaystyle _0^{L_x}}𝑑x{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑p_x{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑p_y\tau v_iv_j{\displaystyle \frac{f_F}{\epsilon }},`$ (3)
where the symbol $`f_F(\epsilon )=\{\mathrm{exp}\{(\mu \epsilon )/kT\}+1\}^1`$ stands for the equilibrium electron Fermi-Dirac distribution function with $`T`$ the temperature and $`\mu `$ the chemical potential which equals the Fermi energy $`E_F=\pi \mathrm{}^2n/m`$ in the zero temperature limit, where $`n`$ is the 2D electron density. Note that the coordinate $`y`$ is excluded from the phase space as the system is homogeneous in that direction. The expressions for the electron velocities follow from the Hamilton equations of motion
$`v_x`$ $`=`$ $`{\displaystyle \frac{\epsilon }{p_x}}={\displaystyle \frac{1}{m}}p_x,`$ (4)
$`v_y`$ $`=`$ $`{\displaystyle \frac{\epsilon }{p_x}}={\displaystyle \frac{1}{m}}\left\{p_y{\displaystyle \frac{eBl_0}{c}}a(x/l_0)\right\}.`$ (5)
First, let us consider the conductivity in the zero temperature limit when the derivative of the Fermi function reduces to a $`\delta `$-function. The component $`\sigma _{yy}`$ can be calculated straightforwardly, as the sample is homogeneuos along the $`y`$-direction all trajectories have to be taken into account. Inserting expression (4) into the conductivity expression (3) we obtain
$`\sigma _{yy}`$ $`=`$ $`{\displaystyle \frac{e^2}{(2\pi \mathrm{})^2L_x}}{\displaystyle _0^{L_x}}𝑑x{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑p_x{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑p_y\tau {\displaystyle \frac{1}{m^2}}`$ (8)
$`\times \left\{p_y{\displaystyle \frac{eBl_0}{c}}a(x/l_0)\right\}^2`$
$`\times \delta \left\{E_F{\displaystyle \frac{1}{2m}}\left\{p_x^2+\left(p_y{\displaystyle \frac{eBl_0}{c}}a(x/l_0)\right)^2\right\}\right\},`$
which leads to the well-known magnetic field independent result
$$\sigma _{yy}=\sigma _0=\frac{e^2E_F\tau }{2\pi \mathrm{}^2}.$$
(9)
We assumed that the relaxation time depends only on the electron energy and therefore for the case of zero temperature it can be replaced by the constant value $`\tau =\tau (E_F)`$. Thus the weak magnetic field modulation in the $`x`$-direction does not change the conductivity in the $`y`$-direction, as expected.
The calculation of $`\sigma _{xx}`$ is more complicated because for sufficiently strong magnetic fields (or small electron velocity, but such that $`\mathrm{}\omega _cE_F`$) some electrons can be forced into snake orbits. The electron motion on such orbits oscillates in the $`x`$-direction around the average value $`x_0`$. Therefore, such electrons do not contribute to $`j_x`$, and consequently, in the expression for the conductivity $`\sigma _{xx}`$ those snake orbits have to be excluded. The classification of all possible electron orbits are given in Fig. 1 where the Fermi surface $`\epsilon (p_x,p_y,x)=E_F`$ is plotted. In the case of zero temperature only electrons with trajectories on the Fermi surface contribute to the conductivity integral. As the energy (1) does not depend on $`y`$, the momentum $`p_y`$ is conserved, and consequently the trajectories are defined by the intersection of the above Fermi surface with the $`p_y=C^{te}`$ planes. It is apparent from Fig. 1 that there are two types of trajectories. The trajectories as indicated by symbol $`D`$ are able to run along the whole $`x`$-axis. Such electrons are moving along open trajectories and they will contribute to the current along the $`x`$-direction. The other trajectories, indicated by symbol $`E`$, are closed and they correspond to the snake orbits. Thus we can separate the conductivity into two parts
$$\sigma _{xx}=\sigma _0\sigma _{xx}^{\mathrm{s}.\mathrm{o}.},$$
(10)
where the symbol $`\sigma _{xx}^{\mathrm{s}.\mathrm{o}.}`$ stands for the snake orbit contribution. The latter term defines the decrement of the conductivity due to the modulated magnetic field, and in the limiting case of small magnetic fields it is proportional to the increase of the magnetoresistance due to the magnetic field modulation.
The snake orbits are located above the plane $`A`$ and below the plane $`B`$. Those planes are defined by $`p_y(A)=\sqrt{2mE_F}+(eBl_0/c)a_{\mathrm{min}}`$, and $`p_y(B)=\sqrt{2mE_F}+(eBl_0/c)a_{\mathrm{max}}`$, with $`a_{\mathrm{min}}=\mathrm{min}\{a(x)\}`$ and $`a_{\mathrm{max}}=\mathrm{max}\{a(x)\}`$ the extremal points of the vector potential $`a(x)`$. Thus we can write
$$\sigma _{xx}^{\mathrm{s}.\mathrm{o}.}=\sigma _A+\sigma _B,$$
(11)
where the contribution from the snake orbits above the $`A`$ plane is given by
$`\sigma _A=`$ $`{\displaystyle \frac{e^2}{(2\pi m\mathrm{})^2L_x}}{\displaystyle _0^{L_x}}𝑑x{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑p_xp_x^2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑p_y\tau `$ (15)
$`\times \delta \left\{E_F{\displaystyle \frac{1}{2m}}\left[p_x^2+\left(p_y{\displaystyle \frac{eBl_0}{c}}a(x/l_0)\right)^2\right]\right\}`$
$`\times \mathrm{\Theta }\left\{p_y{\displaystyle \frac{eBl_0}{c}}a_{\mathrm{min}}\sqrt{2\pi m\mu }\right\},`$
$`={\displaystyle \frac{e^2E_F\tau }{2\pi ^2\mathrm{}^2L_x}}{\displaystyle _0^{L_x}}𝑑x{\displaystyle _{\phi _0}^{\phi _0}}\mathrm{sin}^2(\phi )𝑑\phi .`$
Here the angular integration has to be performed over the white sections of the circle C in Fig. 1, or equivalently, along the bold arch shown in Fig. 2. The limiting angle is defined as $`\phi _0=\mathrm{arccos}(1\mathrm{\Delta }p_y/\sqrt{2mE_F})`$, with $`\mathrm{\Delta }p_y=(eBl_0/c)\left\{a(x/l_0)a_{\mathrm{min}}\right\}`$. In the asymptotic case of small magnetic field modulations it becomes $`\phi _0=\sqrt{2\mathrm{\Delta }p_y/\sqrt{2mE_F}}`$. In the latter case the integration in expression (15) leads to
$`\sigma _A={\displaystyle \frac{e^2E_F\tau }{2\pi ^2\mathrm{}^2L_x}}{\displaystyle _0^{L_x}}𝑑x{\displaystyle _{\phi _0}^{\phi _0}}\phi ^2𝑑\phi =\sigma _{yy}c_A\left({\displaystyle \frac{B}{B_0}}\right)^{3/2},`$ (16)
where $`B_0=(2\pi c\sqrt{mE_F})/(el_0)`$, and
$$c_A=\frac{82^{1/4}\sqrt{\pi }}{3}_0^1𝑑x\{a(x)a_{\mathrm{min}}\}^{3/2}.$$
(17)
The integration over the snake orbits below the plane $`B`$ leads to the same expression (16) except that now the coefficient $`c_A`$ has to be replaced by
$$c_B=\frac{82^{1/4}\sqrt{\pi }}{3}_0^1𝑑x\{a_{\mathrm{max}}a(x)\}^{3/2}.$$
(18)
Thus the resistance change due to the modulated magnetic field becomes $`\mathrm{\Delta }R_{xx}/R_0=(c_A+c_B)(B/B_0)^{3/2}`$ where $`R_0`$ is the resistance in the absence of any magnetic field modulation.
As a special case let us consider a simple cosine periodic magnetic field modulation (with period $`l_0`$) $`b(x)=\mathrm{cos}(2\pi x)`$, which leads to $`a(x)=\mathrm{sin}(2\pi x)/2\pi `$. In this case the coefficients can be easily evaluated
$`c_0=c_A+c_B`$ $`={\displaystyle \frac{82^{3/4}}{3\pi }}{\displaystyle _{1/4}^{1/4}}𝑑x\{\mathrm{sin}(2\pi x)+1\}^{3/2}`$ (20)
$`={\displaystyle \frac{642^{1/4}}{9\pi ^2}}0.86.`$
In the present classical ballistic situation an electron which passed through the first magnetic barrier will also pass through the other barriers. As a consequence the above result can also be applied to the one barrier situation. The periodic oscillating Fermi surface shown in Fig. 1 reduces now to a single step. The trajectories of $`D`$-type (i.e. the open orbits) contribute to the $`\sigma _{xx}`$ conductivity, but the snake orbits (trajectories of $`A`$-type) are replaced by trajectories which reflect from the barrier. The integral in expression (16) must now be evaluated over those reflected trajectories. When the barrier thickness $`l_oL_x`$, the integral $`_{L_x/2}^{L_x/2}𝑑x(a(x/l_0)a_{\mathrm{min}})^{3/2}`$ becomes the sample length multiplied by the total vector potential increment over the barrier: $`L_x\{a_{\mathrm{max}}a_{\mathrm{min}}\}^{3/2}`$, where $`a_{\mathrm{max}}=a(L_x/2)`$ and $`a_{\mathrm{min}}=a(L_x/2)`$.
A single magnetic field barrier in a 2DEG, created experimentally by parallel magnetization of magnetic strips placed on top of the 2DEG, can be represented by
$$B_z(x)=\frac{\mu _0M}{4\pi }\mathrm{ln}\frac{x^2+d^2}{x^2+(d+D)^2}.$$
(21)
The magnetic strip has a thickness $`D`$ with magnetization $`M`$ and is placed a distance $`d`$ from the 2D electron system. The vector potential is obtained by integrating (14) over $`(L_x/2,L_x/2)`$ which gives $`A_{\mathrm{max}}A_{\mathrm{min}}=\mu _0MD/2`$ where use has been made of $`L_xd,D`$ which is valid in typical experimental situations and $`A=Bl_0a`$. Finally, we obtain for the contribution of the reflected trajectories to the conductivity
$$\frac{\sigma _{xx}^{\mathrm{r}.\mathrm{t}.}}{\sigma _{yy}}=\frac{\mathrm{\Delta }R_{xx}}{R_0}=\frac{2^{1/4}}{3\pi }\left(\frac{e\mu _0MD}{c\sqrt{mE_F}}\right)^{3/2}.$$
(22)
The main feature of the obtained magnitoresistance is its non analytical behavior $`B^{3/2}`$ for small magnetic field amplitudes. It is remarkable that it does not depend on the actual form of the modulating field, but it is determined by the density of snake orbits (or reflected trajectories) at the Fermi surface.
For non zero temperature, the previous effect will be suppressed by thermal fluctuations. To generalize our results to non zero temperature we have to replace any function $`G(E_F)`$ which depends on the Fermi energy $`E_F`$ by the corresponding average over the derivative of the Fermi function
$$G(\mu )=\frac{1}{T}_0^{\mathrm{}}𝑑\epsilon \frac{\mathrm{exp}\left((\epsilon \mu )/k_BT\right)G(\epsilon )}{\left\{\mathrm{exp}\left((\epsilon \mu )/k_BT\right)+1\right\}^2}.$$
(23)
Consequently, taking into account $`\sigma _{yy}`$ (9) and the dependence of $`B_0`$ on $`E_F`$ we obtain for the snake orbit contribution to the conductivity
$`\sigma _{xx}^{\mathrm{s}.\mathrm{o}.}(T)=`$ $`\left({\displaystyle \frac{e^2\tau c_0}{2\pi \mathrm{}^2}}\right)\left({\displaystyle \frac{Bel_0}{c\sqrt{m}}}\right)^{3/2}{\displaystyle \frac{1}{T}}`$ (25)
$`\times {\displaystyle _0^{\mathrm{}}}d\epsilon {\displaystyle \frac{\mathrm{exp}\left((\epsilon \mu )/k_BT\right)\epsilon ^{1/4}}{\left\{\mathrm{exp}\left((\epsilon \mu )/k_BT\right)+1\right\}^2}},`$
where we assumed that the relaxation time does not depend on the energy. In the limit of small temperatures the integral can be evaluated analytically and we arrive at the final expression for the conductivity along the direction of the magnetic field modulation
$$\frac{\sigma _{xx}}{\sigma _0}=\left\{1c_0(T)\left(\frac{B}{B_0}\right)^{3/2}\right\},$$
(26)
where
$$c_0(T)=c_0\left\{1\frac{\pi ^2}{32}\left(\frac{k_BT}{\mu }\right)^2\right\}.$$
(27)
Notice that temperature decreases the nonanaliticity coefficient $`c_0`$ (with about 30% when $`k_BT=\mu `$) but does not influence the power law dependence.
Tunneling through magnetic barriers in the ballistic regime was studied numerically in Refs. . For weak barrier only electrons impinging on the magnetic barrier under an angle $`\theta \pi /2`$ (i.e. $`\mathrm{cos}(\theta )(\pi /2\theta )=\varphi `$) are reflected and thus we obtain from Eq. (6) of Ref. the conductance change due to tunneling through a weak magnetic barrier
$$\mathrm{\Delta }GG_0_0^{\phi _0}\phi 𝑑\phi ,$$
(28)
where $`G_0=e^2mv_FL_y/\mathrm{}^2`$. Notice that the difference with the diffusive case is the power of the angle $`\phi `$ in the expression for the conductance/conductivity. In the case of ballistic transport the conductance is proportional to $`v_x`$, while for diffusive transport, see Eq. (3), we have $`v_x^2`$. This difference leads to a linear $`B`$-dependence in the ballistic regime
$$\frac{\mathrm{\Delta }G}{G_0}=\frac{\mathrm{\Delta }R_{xx}}{R_0}=\frac{\mathrm{\Delta }p_y}{\sqrt{2mE_F}}=\frac{e\mu _0MD}{2c\sqrt{2mE_F}}.$$
(29)
The above expression is obtained for the single barrier (21) and agrees with Eq. (3) of Ref. .
Lets compare these results with the experiments of Refs. . In Ref. a Co strip of thickness $`D=90nm`$ was placed a distance $`d=35nm`$ from a 2DEG formed in a GaAs-heterostructure with $`E_F=15.7meV`$. The magnetization of the Co strip was $`\mu _0M9B`$ with a saturation magnetization of $`1.6T`$. In this experiment one is in the diffusive regime and inserting these values in Eq. (15) we obtain $`\mathrm{\Delta }R_{xx}/R_0=4.3B(T)^{3/2}`$. The estimated zero field resistance was $`R_0=2.35\mathrm{\Omega }`$ which results into $`\mathrm{\Delta }R_{xx}(\mathrm{\Omega })=10.1B(T)^{3/2}`$ and agrees with the experimental low magnetic field behavior $`\mathrm{\Delta }R_{xx}(\mathrm{\Omega })=9B(T)^{3/2}`$.
In contrast, the experiments of Vančura et al are closer to the ballistic regime (they have been performed on shorter samples, i.e. $`L_y=34\mu m`$). They placed rectangular Co dots of thickness $`D=100nm`$ above a 2DEG of width $`L_y=W=20nm`$ in a GaAs-heterostructure with $`E_F=19.7meV`$. Inserting these values in Eq. (19) we obtain $`\mathrm{\Delta }R_{xx}/R_0=3.6B(T)`$ where $`R_0=1/G_0=1.01\mathrm{\Omega }`$. This gives $`\mathrm{\Delta }R_{xx}(\mathrm{\Omega })=3.96B(T)`$ which is a factor 2.2 smaller than the experimental result $`\mathrm{\Delta }R_{xx}(\mathrm{\Omega })=8.75B(T)`$. At the saturation field of $`1.5T`$ we find theoretically $`\mathrm{\Delta }R_{xx}=5.94\mathrm{\Omega }`$ which compares with the experimental result $`\mathrm{\Delta }R_{xx}=3.5\mathrm{\Omega }`$ and which is a factor 1.7 smaller. Thus on the average we find a reasonable agreement between theory and experiment but it is clear that all details are not yet fully understood. The discrepancy may be due to the fact that the Co dots are not homogeneously magnetized. Nevertheless, the linear magnetic field dependence is nicely reproduced.
In conclusion we obtained the change in the magneto-resistance due to the presence of magnetic barriers with zero average magnetic field. We found that for diffusive transport the magneto-resistance increases as $`c(T)B^{3/2}`$ with the amplitude of the magnetic barrier ($`B`$) where the coefficient $`c(T)`$ decreases with increasing temperature. This result is different from the ballistic regime where the increase is linear in $`B`$.
Acknowledgments: This work is partially supported by the Flemish Science Foundation (FWO-Vl), IMEC and IUAP-IV. One of us (FMP) is a research director with FWO-Vl and he acknowledges discussions with V. Kubrak, T. Ihn and B. Gallagher. |
warning/0001/astro-ph0001264.html | ar5iv | text | # The giant X-ray outbursts from nearby, non-active galaxies: tidal disruption flares ?
## 1 Introduction
There is strong evidence for massive dark objects at the centers of many galaxies. Observations of the dynamics of stars have been used to derive constraints on the mass of the nucleus. An alternative approach to probe the conditions in the nuclear region, and to detect supermassive black holes (SMBHs) if they are there was suggested by Lidskii & Ozernoi (1979) and Rees (1988, 1990). They proposed to look for these SMBHs based on the flare of electromagnetic radiation, emitted by a star when it is tidally disrupted and accreted by the black hole. The first good candidates of these kind of events have been reported in the last few years.
## 2 X-ray variability of AGN
Many active galactic nuclei (AGN) are variable in X-rays with a range of amplitudes and on many different timescales (e.g., Mushotzky et al. 1993). The cause of variability is usually linked in some way or the other to the central engine; e.g., by changes in the accretion disk.
The source’s variability behavior provides important information on the emission mechanisms, the size and geometry of the central region, and the physical conditions in the illuminated circum-nuclear gas like the broad line region and the warm absorber.
Whereas X-ray variability by factors of about 2 – 3 is a common property of AGN, X-ray outbursts by factors of order 50 or more are extremely rare; only very few objects have been reported to show such behavior. Among these are E1615+061 (Piro et al. 1988) and IRAS 13224 with its repeated X-ray outbursts (Otani et al. 1996, Boller et al. 1997). Different mechanisms to account for the X-ray variability have been favored for these galaxies: a variable soft excess for E1615+061 (Piro et al. 1988, see also Piro et al. 1997), relativistic effects in the accretion disk for IRAS 13224 (Boller et al. 1997). The Narrow-line Seyfert 1 galaxy RXJ0134-4258 underwent a dramatic spectral change from ultra-soft to flat within 2 years, corresponding to huge X-ray variability in the hard band. The cause of this peculiar behavior is presently unclear. The model that has been studied in most detail so far is the presence of a time-variable warm absorber (Komossa & Meerschweinchen 2000, and references therein).
Tab. 1 gives some examples of large-amplitude variability (factors of 10 or more) of active galaxies. The list is not complete, but shows that large-amplitude variability occurs in all types of AGN. In addition, among those with the largest factors of variability are some galaxies which are not active at all; these will be discussed in more detail in the next sections of this contribution.
## 3 Flares from tidally disrupted stars
Questions of particular interest in the context of AGN evolution are: what fraction of galaxies have passed through an active phase, and how many now have non-accreting and hence unseen SMBHs at their centers (e.g., Rees 1989)?
Several approaches were followed to study these questions. Much effort has concentrated on deriving central object masses from studies of the dynamics of stars and gas in the nuclei of nearby galaxies. Earlier ground-based evidence for central quiescent dark masses in non-active galaxies (e.g., Tonry 1987, Dressler & Richstone 1988, Kormendy & Richstone 1992) has been strengthened by recent HST results (e.g., van der Marel et al. 1997, Kormendy et al. 1996; see Kormendy & Richstone 1995 for a review). There is now excellent evidence for a SMBH in our galactic center as well (Eckart & Genzel 1996).
On the other hand, X-rays trace the very vicinity of the SMBH. Lidskii & Ozernoi (1979) and Rees (1988, 1990) suggested to use the flare of electromagnetic radiation predicted when a star is tidally disrupted and accreted by a SMBH as a means to detect SMBHs in nearby non-active galaxies.
Depending on its trajectory, a star gets tidally disrupted after passing a certain distance to the black hole (e.g., Hills 1975, Lidskii & Ozernoi 1979, Diener et al. 1997), the tidal radius, given by
$$r_\mathrm{t}r_{}(\frac{M_{\mathrm{BH}}}{M_{}})^{\frac{1}{3}}.$$
(1)
The star is first heavily distorted, then disrupted. About half of the gaseous debris will be unbound and gets lost from the system (e.g., Young et al. 1977). The rest will be eventually accreted by the black hole (e.g., Cannizzo et al. 1990, Loeb & Ulmer 1997). The debris, first spread over a number of orbits, quickly circularizes (e.g., Rees 1988, Cannizzo et al. 1990) due to the action of strong shocks when the most tightly bound debris interacts with other parts of the stream (e.g., Kim et al. 1999). Most orbital periods will then be within a few times the period of the most tightly bound matter (e.g., Evans & Kochanek 1989; see also Nolthenius & Katz 1982, Luminet & Marck 1985).
A star will only be disrupted if its tidal radius lies outside the Schwarzschild radius of the black hole, else it is swallowed as a whole (this happens for black hole masses larger than about 10<sup>7</sup> M; in case of a Kerr black hole, tidal disruption may occur even for larger BH masses if the star approaches from a favorable direction (Beloborodov et al. 1992)). Larger BH masses may still strip the atmospheres of giant stars. Most theoretical work focussed on stars of solar mass and radius so far.<sup>1</sup><sup>1</sup>1Numerical simulations of the disruption process, the stream-stream collision, the accretion phase, the change in angular momentum of the black hole, the changes in the stellar distribution of the surroundings, and the disruption rates have been studied in the literature (e.g., Nduka 1971, Masshoon 1975, Nolthenius & Katz 1982, 1983, Carter & Luminet 1985, Luminet & Marck 1985, Evans & Kochanek 1989, Laguna et al. 1993, Diener et al. 1997; Lee et al. 1995, Kim et al. 1999; Hills et al. 1975, Gurzadyan & Ozernoi 1979, 1980, Cannizzo et al. 1990, Loeb & Ulmer 1997, Ulmer et al. 1998; Beloborodov et al. 1992; Frank & Rees 1976, Rauch & Ingalls 1998, Rauch 1999; Syer & Ulmer 1999, Magorrian & Tremaine 1999).
Explicit predictions of the emitted spectrum and luminosity during the disruption process and the start of the accretion phase are still rare (see Sect. 6.3 for details). The emission is likely peaked in the soft X-ray or UV portion of the spectrum, initially (e.g., Rees 1988, Kim et al. 1999, Cannizzo et al. 1990; see also Sembay & West 1993).
## 4 Tidal disruption events in active galaxies
Tidal disruption has occasionally been invoked to explain some exceptional events of variability in AGN or some general properties of AGN or LINERs.
The possibility of tidal disruption of a star by a SMBH was originally proposed as a means of fueling active galaxies (Hills 1975), but was later dismissed. Tidal disruption was invoked by Eracleous et al. (1995) in a duty cycle model to explain the UV brightness/darkness of LINERs.
Peterson & Ferland (1986) suggested this mechanism as possible explanation for the transient brightening and broadening of the HeII line observed in the Seyfert galaxy NGC 5548. Variability in the Balmer lines of some AGN (the appearance and disappearance of a broad component in H$`\beta `$ or H$`\alpha `$) has recently been interpreted in the same way.
Brandt et al. (1995) observed a giant X-ray outburst from the galaxy IC 3599 (Zwicky 159.034). The source was in its X-ray high-state during the ROSAT all-sky survey (RASS hereafter) and then declined in intensity within years. The peak luminosity exceeded $`L_\mathrm{x}`$ = 10<sup>43</sup> erg/s and the outburst spectrum was very soft (photon index $`\mathrm{\Gamma }_\mathrm{x}4`$ when fit by a powerlaw). An optical spectrum of the galaxy taken shortly after the X-ray outburst was characterized by strong emission lines from highly ionized species like FeX and HeII. These lines then declined in strength in subsequent years (Bade et al. 1995, their Fig. 8; Grupe et al. 1995) proving the association of the X-ray flare with the nucleus of this galaxy. Photoionization models for the outburst line emission were presented by Komossa & Bade (1999). The values of gas density, column density, and ionization parameter required to reproduce the observed emission line intensities are typical of a BLR or CLR, and the ionizing spectrum is characterized by a strong soft excess as observed, again confirming the association of the flare with IC 3599. Some uncertainties existed in the classification of this galaxy based on optical spectra: Brandt et al. (1995) noted that the outburst spectrum was Narrow-line Seyfert 1-like. The optical spectrum in quiescence is different, and was preliminary classified as starburst by Bade et al. (1995). Spectra of higher sensitivity and spectral resolution were then presented by Komossa & Bade (1999) who detected a broad component in H$`\alpha `$ that argues for a Seyfert 1.9 classification of IC 3599. This latter paper also presents several further arguments that IC 3599 shows permanent Seyfert activity.
## 5 Tidal disruption flares from non-active galaxies
In the UV spectral region, two UV spikes were detected at and near the center of the elliptical galaxy NGC 4552<sup>2</sup><sup>2</sup>2We include the case of NGC 4552 in this section, but note that there are several indications of very weak permanent activity in this galaxy (Renzini et al. 1995, Cappellari et al. 1999).. The central flare was interpreted by Renzini et al. (1995) as accretion event (the tidal stripping of a star’s atmosphere by a SMBH, or the accretion of a molecular cloud).
The discovery of a giant flare at soft X-ray energies from NGC 5905 was reported by Bade et al. (1996). The X-ray properties of the galaxy can be summarized as follows (see also Tab. 2, and Figs 1,2): (i) The X-ray spectrum during outburst was ultra-soft ($`kT_{\mathrm{bb}}`$ = 0.06 keV). (ii) The total amplitude of variability amounts to a factor of $``$200. (iii) The observed peak luminosity reached $`L_\mathrm{x}\text{ }>10^{4243}`$ erg/s. High quality optical spectra of this galaxy prior to the X-ray flare (Ho et al. 1995), and several years after the outburst (Schombert 1998, Komossa & Bade 1999) are of HII-type, with no signs of Seyfert activity. Komossa & Bade (1999) presented follow-up observations and discussed outburst scenarios. A summary of their results is given in the next section, plus an extended discussion of the possibility that the X-ray flare was due to a tidal disruption event.
A similar event was detected from the direction of the galaxy pair RXJ1242–1119 (Komossa & Greiner 1999). In this case, the flare luminosity was even higher. It reached nearly 10<sup>44</sup> erg/s in the ROSAT X-ray band (Tab. 2). Optical spectra taken of both galaxies reveal them to be non-active. No emission lines were detected.
## 6 Outburst scenarios
### 6.1 Alternatives to tidal disruption
Firstly, we note that based purely on a positional coincidence, an interlopper (flaring Galactic foreground object) could not be excluded, given the limited spatial positional accuracy of ROSAT of at least several arcseconds.
However, known populations of galactic flaring sources show different temporal properties. In addition, the growing number of X-ray flares detected at the locations of bright, nearby galaxies (NGC 5905, RXJ1242-11, RXJ1614+75; a further candidate is presented by Reiprich & Greiner 2000), makes a chance coincidence increasingly unlikely.
Other sources of the X-ray emission related to sources within the galaxies NGC 5905 and RXJ1242–11 were reviewed by Komossa & Bade (1999) in some detail, including some order of magnitude estimates: Most outburst scenarios do not survive close scrutiny, because they cannot account for the huge maximum luminosity (e.g., X-ray binaries within the galaxies, or a supernova in dense medium), require extreme fine-tuning (e.g., a warm-absorbed hidden Seyfert nucleus), are inconsistent with the optical observations (gravitational lensing), or predict a different temporal behavior (X-ray afterglow of a Gamma-ray burst).
### 6.2 Tidal disruption model
Intense electromagnetic radiation will be emitted in three phases of the disruption and accretion process: First, during the stream-stream collision when different parts of the bound stellar debris first interact with themselves (Rees 1988). Kim et al. (1999) have carried out numerical simulations of this process and find that the initial luminosity burst due to the collision may reach 10<sup>41</sup> erg/s, under the assumption of a BH mass of 10<sup>6</sup> M and a star of solar mass and radius. Secondly, radiation is emitted during the accretion of the stellar gaseous debris. Finally, the unbound stellar debris leaving the system may shock the surrounding interstellar matter like in a supernova remnant and cause intense emission.
The luminosity emitted if the black hole is accreting at its Eddington luminosity can be estimated by $`L_{\mathrm{edd}}1.3\times 10^{38}M/M_{}`$ erg/s. In case of NGC 5905, a BH mass of at least $`10^5`$ M would be required to produce the observed $`L_\mathrm{x}`$, and a higher mass if $`L_\mathrm{x}`$ was not observed at its peak value. For comparison, BH masses of $`M_{\mathrm{BH}}\text{ }<10^{67}\mathrm{M}_{}`$ have recently been reported by Salucci et al. (1999) for the centers of some late-type spiral galaxies. Alternatively, the atmosphere of a giant star could have been stripped instead of a complete disruption event. It is interesting to note that NGC 5905 possesses a complex bar structure (Friedli et al. 1996) which might aid in the fueling process by disturbing the stellar velocity fields.
Using the black body fit to the X-ray spectra of NGC 5905 and RXJ1242–11, we find the fiducial black body radius to be located between the last stable orbit of a Schwarzschild black hole, and inside the tidal radius.
We note that many details of the tidal disruption and the related processes are still unclear. In particular, the flares cannot be standardised. Observations would depend on many parameters, like the type of disrupted star, the impact parameter, the spin of the black hole, effects of relativistic precession, and the radiative transfer is complicated by effects of viscosity and shocks (Rees 1990). Uncertainties also include the amount of the stellar debris that is accreted (part may be ejected as a thick wind, or swallowed immediately). Related to this is the duration of the flare-like activity, which may be months or years to tens of years (e.g., Rees 1988, Cannizzo et al. 1990, Gurzadyan & Ozernoi 1979).
## 7 Search for further X-ray flares
We performed a search for further cases of strong X-ray variability (Komossa & Bade 1999) using the sample of nearby galaxies of Ho et al. (1995) and ROSAT all-sky survey and archived pointed observations. The sample of Ho et al. has the advantage of the availability of high-quality optical spectra, which are necessary when searching for ‘truly’ non-active galaxies. 136 out of the 486 galaxies in the catalogue were detected in pointed observations. For these, we compared the countrates with those measured during the RASS. We do not find another object with a factor $`\text{ }>`$50 amplitude of variability. Several sources show variability by a factor 10–20 but all of these are well-known AGN.
The absence of any further flaring event among the sample galaxies is entirely consistent with the expected disruption rates of one event in at least $``$10<sup>4</sup> years per galaxy (e.g., Magorrian & Tremaine 1999).
## 8 Future perspectives
Such X-ray outbursts provide important information on the presence of SMBHs in non-active galaxies, the accretion history of the universe, and the link between active and normal galaxies. Future X-ray surveys (like the one that was planned with ABRIXAS, or the one that will be carried out with MAXI) will be valuable in finding further of these outstanding sources.
In particular, rapid follow-up optical observations will be important in order to detect potential emission lines that were excited by the outburst emission. In case of a giant tidal disruption flare in an active galaxy, this would also provide an excellent chance to map the properties of the broad line region.
###### Acknowledgements.
It is a pleasure to thank Jules Halpern, David L. Meier, L.M. Ozernoi, Joachim Trümper and Weimin Yuan for fruitful discussions. The ROSAT project has been supported by the German Bundesministerium für Bildung, Wissenschaft, Forschung und Technologie (BMBF/DLR) and the Max-Planck-Society. The ‘ROSAT - ASCA workshop on AGN’ was funded by the Inter-Research Centers Cooperative Program of the JSPS and the Deutsche Forschungsgemeinschaft DFG.
Preprints of this and related papers can be retrieved from our webpage at http://www.xray.mpe.mpg.de/$``$skomossa/ |
warning/0001/hep-lat0001030.html | ar5iv | text | # References
UT-Komaba 00-3
Pions in Lattice QCD with the Overlap Fermions at Strong Gauge Coupling
Ikuo Ichinose<sup>1</sup><sup>1</sup>1e-mail address: ikuo@hep1.c.u-tokyo.ac.jpand Keiichi Nagao<sup>2</sup><sup>2</sup>2e-mail address: nagao@hep1.c.u-tokyo.ac.jp
Institute of Physics, University of Tokyo, Komaba, Tokyo, 153-8902 Japan
Abstract
In the previous paper we developed a strong-coupling expansion for the lattice QCD with the overlap fermions and showed that Lüsher’s “extended” chiral symmetry is spontaneously broken in some parameter region of the overlap fermions. In this paper, we obtain a low-energy effective action and show that there exist quasi-Nambu-Goldsone bosons which are identified as the pions. The pion field is a nonlocal composite field of quark and anti-quark even at the strong-coupling limit because of the nonlocality of the overlap fermion formalism and Lüsher’s chiral symmetry. The pions become massless in the limit of the vanishing bare-quark mass as it is desired.
One of the long standing problems in the lattice gauge theory is the formulation of lattice fermions. Recently a very promising formulation named overlap fermion was proposed by Narayanan and Neuberger and it has been studied intensively by both analytic and numerical methods. In the previous paper (which we shall refer to paper I hereafter), we slightly extended the overlap fermion by introducing a “hopping” parameter $`t`$ and studied the lattice QCD by using both the $`t`$-expansion and the strong-coupling expansion. There we calculated the effective potential of the chiral condensation and showed that Lüsher’s extended chiral symmetry is spontaneously broken at certain parameter region of the overlap fermion. In this paper we shall obtain an effective action of low-energy excitations and show that there exist quasi-Nambu-Goldsone bosons which are identified as the pions. As we show, the pion field is a nonlocal composite field of quark and anti-quark even at the strong-coupling limit.
Action of the overlap fermion on the $`d`$-dimensional lattice is given as follows,
$$S_F=a^d\underset{n,m}{}\overline{\psi }(m)D(m,n)\psi (n),$$
(1)
where the covariant derivative $`D(m,n)`$ is defined as
$`D`$ $`=`$ $`{\displaystyle \frac{1}{a}}\left(1+X{\displaystyle \frac{1}{\sqrt{X^{}X}}}\right),`$
$`X_{mn}`$ $`=`$ $`\gamma _\mu C_\mu (t;m,n)+B(t;m,n),`$
$`C_\mu (t;m,n)`$ $`=`$ $`{\displaystyle \frac{t}{2a}}\left[\delta _{m+\mu ,n}U_\mu (m)\delta _{m,n+\mu }U_\mu ^{}(n)\right],`$
$`B(t;m,n)`$ $`=`$ $`{\displaystyle \frac{M_0}{a}}+{\displaystyle \frac{r}{2a}}{\displaystyle \underset{\mu }{}}\left[2\delta _{n,m}t\delta _{m+\mu ,n}U_\mu (m)t\delta _{m,n+\mu }U_\mu ^{}(n)\right],`$ (2)
where $`r`$ and $`M_0`$ are dimensionless nonvanishing free parameters of the overlap lattice fermion formalism and $`U_\mu (m)`$ is gauge field on links. Other notations are standard. We have introduced a new parameter $`t`$.<sup>3</sup><sup>3</sup>3As we explained in the paper I, the $`t`$-dependence of the operator $`D(m,n)`$ is abosorbed by a redefinition of $`M_0`$. The original overlap fermion corresponds to $`t=1`$. For notational simplicity, we define
$$A\frac{1}{a}(drM_0),B\frac{rt}{2a},C\frac{t}{2a},$$
(3)
and
$`\mathrm{\Gamma }_\mu ^{}(m,n)`$ $``$ $`\delta _{m+\mu ,n}U_\mu (m)\delta _{m,n+\mu }U_\mu ^{}(n),`$
$`\mathrm{\Gamma }_\mu ^+(m,n)`$ $``$ $`\delta _{m+\mu ,n}U_\mu (m)+\delta _{m,n+\mu }U_\mu ^{}(n).`$ (4)
In terms of the above quantities,
$$X_{mn}=A\delta _{mn}+C\gamma _\mu \mathrm{\Gamma }_\mu ^{}(m,n)B\mathrm{\Gamma }_\mu ^+(m,n),$$
(5)
$$(X^{})_{mn}=A\delta _{mn}C\gamma _\mu \mathrm{\Gamma }_\mu ^{}(m,n)B\mathrm{\Gamma }_\mu ^+(m,n).$$
(6)
From Eq.(3), $`B,C=O(t)`$ and we consider $`A=O(1)`$ in the later discussion. Then it is rather straightforward to expand $`D(m,n)`$ in powers of $`t`$,
$`aD(m,n)`$ $`=`$ $`2\theta (A)\delta _{mn}+{\displaystyle \frac{C}{|A|}}{\displaystyle \gamma _\mu \mathrm{\Gamma }_\mu ^{}(m,n)}`$ (7)
$`+{\displaystyle \frac{BC}{2A|A|}}{\displaystyle \gamma _\mu \left(\mathrm{\Gamma }_\mu ^{}(m,l)\mathrm{\Gamma }_\nu ^+(l,n)+\mathrm{\Gamma }_\nu ^+(m,l)\mathrm{\Gamma }_\mu ^{}(l,n)\right)}`$
$`+{\displaystyle \frac{C^2}{2A|A|}}{\displaystyle \gamma _\mu \gamma _\nu \mathrm{\Gamma }_\mu ^{}(m,l)\mathrm{\Gamma }_\nu ^{}(l,n)}+O(t^3).`$
Higher-order terms of $`t`$ are nonlocal and the $`t`$-expansion corresponds to a kind of the hopping expansion. In the free field case or at the weak gauge coupling, the parameter region in which fermion propagator has no species doublers is easily identified in the $`(M_0,r)`$ plane. However in the strong-coupling theory like QCD, the parameter region of physical relevance should be determined by another requirement, because the pole in the quark propagator is not a physical observable. Therefore it is important to study the lattice QCD with the overlap fermions in rather wide region of the parameter space.
It is verified that the Ginsparg-Wilson (GW) relation
$$D\gamma _5+\gamma _5D=aD\gamma _5D,$$
(8)
is satisfied by the $`t`$-expanded $`D(m,n)`$ in Eq.(7) at each order of $`t`$. Action of the fermion $`S_F`$ in Eq.(1) is invariant under the following extended chiral transformation discovered by Lüscher ,
$$\delta \psi (m)=ϵ\gamma _5\left(\delta _{nm}aD(m,n)\right)\psi (n),\delta \overline{\psi }(m)=ϵ\overline{\psi }(m)\gamma _5,$$
(9)
where $`ϵ`$ is an infinitesimal transformation parameter.
Total action of the lattice QCD is given by
$`S_{tot}`$ $`=`$ $`S_G+S_{F,M},`$
$`S_G`$ $`=`$ $`{\displaystyle \frac{1}{g^2}}{\displaystyle \underset{pl}{}}\text{Tr}(UUU^{}U^{}),`$
$`S_{F,M}`$ $`=`$ $`S_FM_B{\displaystyle \overline{\psi }(m)\psi (m)},`$ (10)
where we have added the bare mass term of quarks. We shall consider the strong-coupling limit in this paper though a systematic strong-coupling expansion is possible. We consider the $`U(N)`$ gauge group for large $`N`$. It is easily verified that the following composite operators are covariant under the transformation (9),
$$\overline{q}(n)\overline{\psi }(n),q(n)\left(1\frac{a}{2}D(n,m)\right)\psi (n),$$
(11)
that is
$$\delta q(m)=ϵ\gamma _5q(m),\delta \overline{q}(m)=ϵ\overline{q}(m)\gamma _5.$$
(12)
Hereafter we often set the lattice spacing $`a=1`$. We consider the case of negative $`A`$ which is expected to have desired properties of QCD .
Partition function of the $`U(N)`$ QCD is given by the following functional integral,
$$Z[J]=D\overline{\psi }D\psi DU\mathrm{exp}\left\{S_{tot}+J(n)\widehat{Q}(n)\right\},$$
(13)
where $`[DU]`$ is the Haar measure and
$`J(n)\widehat{Q}(n)=J_\beta ^\alpha (n)\widehat{Q}_\alpha ^\beta (n)`$
$`\widehat{Q}_\alpha ^\beta (n)={\displaystyle \frac{1}{N}}{\displaystyle \underset{a}{}}q_{a,\alpha }(n)\overline{q}^{a,\beta }(n),`$ (14)
with color index $`a`$ and spinor-flavor indices $`\alpha `$ and $`\beta `$. It should be remarked that the source $`J`$ is coupled to the nonlocal operator $`\widehat{Q}`$ instead of $`_a\psi _{a,\alpha }(n)\overline{\psi }^{a,\beta }(n)`$. In Ref. we studied the gauged Gross-Neveu model with the overlap fermions and showed that $`\widehat{Q}_\alpha ^\beta (n)`$ is the proper composite fields for the extended chiral symmetry, i.e., the order parameter is given by $`\widehat{Q}_\alpha ^\beta (n)\delta _{\alpha \beta }`$ and the Nambu-Goldstone bosons correspond to tr$`{}_{S}{}^{}(\widehat{Q}(n)\gamma _5)`$ where tr<sub>S</sub> is the trace over spinor indices. In the rest of the present paper, we shall obtain the effective action $`S_{eff}(𝒬)`$ defined as
$$Z[J]=D𝒬e^{S_{eff}(𝒬)+J𝒬}$$
(15)
where integral over color-singlet elementary “meson” field $`𝒬_\beta ^\alpha `$ is defined as in Ref..
In order to evaluate the partition function (13), we make a change of variables as
$$(\overline{\psi },\psi )(\overline{q},q).$$
Then the measure of the functional integral is transformed as
$$[d\overline{\psi }d\psi ][d\overline{q}dq]e^{N\mathrm{Tr}\mathrm{ln}(1\frac{\mathrm{a}}{2}\mathrm{D})},$$
(16)
where Tr in (16) is the trace over the spinor-flavor as well as the real-space indices. The contribution from the Jacobian $`\text{Tr}\mathrm{ln}(1\frac{a}{2}D)`$ is easily evaluated by the $`t`$-expansion. In terms of $`\overline{q}`$ and $`q`$,
$`S_F^{}(q)`$ $`=`$ $`S_F(\psi )`$ (17)
$`=`$ $`\overline{q}(m)[{\displaystyle \frac{C}{|A|}}{\displaystyle }\gamma _\mu \mathrm{\Gamma }_\mu ^{}(m,n)+{\displaystyle \frac{BC}{2A|A|}}{\displaystyle }\gamma _\mu (\mathrm{\Gamma }_\mu ^{}(m,l)\mathrm{\Gamma }_\nu ^+(l,n)+\mathrm{\Gamma }_\nu ^+(m,l)\mathrm{\Gamma }_\mu ^{}(l,n))`$
$`+O(t^3)]q(n).`$
The term of $`O((C/A)^2)`$ in (7) has disappeared. As explained in paper I, this term is exactly determined from the term of $`O(C/A)`$ in (7) through the GW relation. We expect that a similar phenomenon occurs for higher-order terms of the $`t`$-expansion which are determined by lower-order terms by the GW relation. That is, the terms of higher-order of $`t`$ which are determined through the GW relation drop when the action is rewritten in terms of $`q`$ and $`\overline{q}`$. The mass term is also given by
$`M_B\overline{\psi }(m)\psi (m)`$ $`=`$ $`M_B\overline{q}(m)[\delta _{mn}{\displaystyle \frac{C}{2A}}{\displaystyle }\gamma _\mu \mathrm{\Gamma }_\mu ^{}(m,n)`$ (18)
$`{\displaystyle \frac{BC}{4A^2}}{\displaystyle \gamma _\mu \left(\mathrm{\Gamma }_\mu ^{}(m,l)\mathrm{\Gamma }_\nu ^+(l,n)+\mathrm{\Gamma }_\nu ^+(m,l)\mathrm{\Gamma }_\mu ^{}(l,n)\right)}`$
$`+O(t^3)]q(n).`$
Total action of $`q`$ is given by
$$S_{F,M}^{}(q)=S_{F,M}(\psi ).$$
(19)
Integral over the gauge field can be performed by the one-link integral,
$$e^{W(\overline{D},D)}=𝑑U_\mu \mathrm{exp}\left[\text{Tr}(\overline{D}_\mu U_\mu +U_\mu ^{}D_\mu )\right].$$
(20)
Explicit form of the one-link integral $`W(\overline{D},D)`$ for the present system is obtained as in paper I. Gauge fields in $`\mathrm{\Gamma }_\mu ^\pm (m,n)`$ are replaced with composite operators of $`q`$ and $`\overline{q}`$ after the integral over $`U_\mu (n)`$. (For details see paper I.) After some calculation,
$$DUe^{S_{F,M}^{}(q)+N\mathrm{Tr}\mathrm{ln}(1\frac{1}{2}\mathrm{D})+{\scriptscriptstyle \mathrm{J}\widehat{\mathrm{Q}}}}=e^{NM_B\text{tr}(\widehat{Q})S_2(\widehat{Q})+{\scriptscriptstyle J\widehat{Q}}},$$
(21)
where tr in (21) is the trace over spinor-flavor indices and $`S_2(\widehat{Q})`$ is some complicated function of $`\widehat{Q}`$ which is obtained in powers of $`t`$. Let us define the following operators;
$`ϵ`$ $`=`$ $`ϵ_\mu ^{\delta \sigma }(m)=\left({\displaystyle \frac{C}{A}}\right)^2\left(1{\displaystyle \frac{M_B}{2}}\right)^2\left(\widehat{Q}(m)\gamma _\mu \widehat{Q}(m+\mu )\gamma _\mu \right)^{\delta \sigma },`$
$`ϵ^{}`$ $`=`$ $`ϵ_\mu ^{{}_{}{}^{}\delta \sigma }(m)=\left({\displaystyle \frac{C}{A}}\right)^2\left(1{\displaystyle \frac{M_B}{2}}\right)^2\left(\widehat{Q}(m+\mu )\gamma _\mu \widehat{Q}(m)\gamma _\mu \right)^{\delta \sigma }.`$ (22)
Then in terms of $`ϵ`$’s, $`S_2(\widehat{Q})`$ is given as<sup>4</sup><sup>4</sup>4As we shall see, $`𝒬O(t^1)`$.
$`{\displaystyle \frac{1}{N}}S_2(\widehat{Q})`$ $`=`$ $`{\displaystyle \underset{m,\mu }{}}\text{tr}\left[g(ϵ_\mu (m))\right]`$ (23)
$`+{\displaystyle \frac{BC^3}{2A^4}}\left(1{\displaystyle \frac{M_B}{2}}\right)^3`$
$`\times {\displaystyle \underset{m,\mu ,\nu }{}}\{\text{tr}\left[𝒬(m+\mu )\gamma _\mu g^{}(ϵ_\mu (m))𝒬(m)\gamma _\mu 𝒬(m+\mu +\nu )\gamma _\nu g^{}(ϵ_\nu (m+\mu ))\right]`$
$`\text{tr}\left[𝒬(m+\mu +\nu )\gamma _\mu g^{}(ϵ_\mu (m+\nu ))𝒬(m+\nu )\gamma _\mu 𝒬(m+\mu )\gamma _\nu g^{}(ϵ_\nu ^{}(m+\mu ))\right]`$
$`+\text{tr}\left[𝒬(m)\gamma _\mu g^{}(ϵ_\mu ^{}(m))𝒬(m+\mu )\gamma _\mu 𝒬(m+\nu )\gamma _\nu g^{}(ϵ_\nu (m))\right]`$
$`\text{tr}\left[𝒬(m+\nu )\gamma _\mu g^{}(ϵ_\mu ^{}(m+\nu ))𝒬(m+\mu +\nu )\gamma _\mu 𝒬(m)\gamma _\nu g^{}(ϵ_\nu ^{}(m))\right]`$
$`+\text{tr}\left[𝒬(m+\nu )\gamma _\nu g^{}(ϵ_\nu (m))𝒬(m)\gamma _\mu 𝒬(m+\mu +\nu )\gamma _\mu g^{}(ϵ_\mu (m+\nu ))\right]`$
$`+\text{tr}\left[𝒬(m+\mu +\nu )\gamma _\nu g^{}(ϵ_\nu (m+\mu ))𝒬(m+\mu )\gamma _\mu 𝒬(m+\nu )\gamma _\mu g^{}(ϵ_\mu ^{}(m+\nu ))\right]`$
$`\text{tr}\left[𝒬(m)\gamma _\nu g^{}(ϵ_\nu ^{}(m))𝒬(m+\nu )\gamma _\mu 𝒬(m+\mu )\gamma _\mu g^{}(ϵ_\mu (m))\right]`$
$`\text{tr}\left[𝒬(m+\mu )\gamma _\nu g^{}(ϵ_\nu ^{}(m+\mu ))𝒬(m+\mu +\nu )\gamma _\mu 𝒬(m)\gamma _\mu g^{}(ϵ_\mu ^{}(m))\right]\}`$
$`+{\displaystyle \frac{N_{sf}C^4}{4A^4}}\left(1{\displaystyle \frac{M_B}{2}}\right)^2{\displaystyle \underset{m,\mu }{}}\text{tr}\left[𝒬(m+\mu )\gamma _\mu g^{}(ϵ_\mu (m))𝒬(m)\gamma _\mu g^{}(ϵ_\mu ^{}(m))\right]`$
$`+O(t^3),`$
where $`N_{sf}`$ is the dimension of the spinor-flavor index and
$$g(x)=1(14x)^{\frac{1}{2}}+\mathrm{ln}\left[\frac{1}{2}(1+(14x)^{\frac{1}{2}})\right].$$
(24)
Elementary meson fields $`𝒬`$ and their functional integral are introduced as in the previous case ,
$`{\displaystyle 𝑑\overline{q}𝑑q\mathrm{exp}\left(\frac{1}{N}J_\alpha ^\beta q_a^\alpha \overline{q}_\beta ^a\right)}`$ $`=`$ $`\left(\text{det}J\right)^N`$ (25)
$`=`$ $`{\displaystyle 𝑑𝒬\left(\text{det}𝒬\right)^Ne^{J𝒬}},`$
where the integral over $`𝒬`$ is defined by the contour integral, i.e., $`𝒬`$ is polar-decomposed as $`𝒬=RV`$ with positive-definite Hermitian matrix $`R`$ and unitary matrix $`V`$, and $`𝑑𝒬𝑑V`$ with the Haar measure of U($`N_{sf}`$). From (25), there appear additional terms like $`(N\text{Tr}\mathrm{log}𝒬)`$ in the effective action. Therefore the effective action is given by
$$S_{eff}(𝒬)=N\underset{n}{}\left[\text{tr}\mathrm{ln}𝒬(n)+M_B\text{tr}𝒬(n)\right]+S_2(𝒬),$$
(26)
Effective potential of the chiral condensate is obtained from $`S_{eff}`$ in Eq.(26) by setting
$$Q^{\alpha \beta }(n)=v\delta _{\alpha \beta }.$$
In paper I we obtained
$$v=\frac{|A|}{2C}\sqrt{\frac{2d1}{d^2}}+O(t^0)+O(M_B).$$
(27)
From the above result, we can expect that there appear quasi-Nambu-Goldstone bosons, i.e., pions. Naive expectation is that pions correspond to the composite operators like $`\overline{\psi }\gamma _5\psi `$ as in the continuum. We examined the effective action and found that there is no gapless excitation in the channel $`\overline{\psi }\gamma _5\psi `$. However we found that there are massless modes in the channel $`\text{tr}_F\left(\widehat{Q}\gamma _5\right)=\text{tr}_F\left(q\overline{q}\gamma _5\right)`$. It is straightforward to obtain the effective action of the pions by inserting the following expression of $`𝒬`$ into $`S_{eff}`$ in Eq.(26),
$$𝒬(m)=ve^{i\gamma _5\varphi _5(m)}.$$
(28)
For example from (22) and (28),
$$ϵ_\mu e^{i\gamma _5_\mu \varphi _5}.$$
(29)
We obtain
$$S_{eff}|_{𝒬(m)=ve^{i\gamma _5\varphi _5(m)}}=2^{\frac{d}{2}}N\left[C_\pi \underset{m,\mu }{}\text{tr}_F\left(_\mu \varphi _5(m)\right)^2\frac{M_Bv}{2}\underset{m}{}\text{tr}_F\left(\varphi _5(m)\right)^2\right],$$
(30)
where $`C_\pi `$ is some positive constant. Therefore the fields $`\varphi _5`$ are quasi-Nambu-Goldstone pions as expected.
If we introduce elementary “meson” fields $``$ as in paper I, i.e.,
$$_\alpha ^\beta (m)\frac{1}{N}\underset{a}{}\psi _{a,\alpha }(m)\overline{\psi }^{a,\beta }(m)$$
and parameterize them as
$$(m)=ve^{i\gamma _5\stackrel{~}{\varphi }_5(m)},$$
it is shown that there appears additional mass term of $`\stackrel{~}{\varphi }_5`$ which is finite for $`M_B0`$. This fact must be important for numerical studies of QCD with the overlap fermions.
In this paper we consider the leading-order of $`1/N`$. We expect that in the next-leading order of $`1/N`$ a mass term of the flavour singlet meson will appear from the Jacobian in (16). This is a solution to the $`U(1)`$ problem. The next-leading order terms in $`1/N`$ is under study and results will be reported in a future publication.
Finally let us comment on the domain wall fermions at strong coupling. Very recently lattice $`U(1)`$ gauge model with the domain-wall fermions was studied in the strong-coupling limit by the Hamiltonian formalism . There an effective Hamiltonian for low-lying color-singlet degrees of freedom is obtained by treating the terms proportional to the gauge fields $`U_\mu (m)`$ as perturbations. This idea is very close to ours in paper I and the present paper though we employ the Lagrangian formalism. From the effective Hamiltonian obtained there, it is concluded that domain-wall fermions at strong coupling suffer both the doubling problem and explicit breaking of chiral symmetry. Furthermore it is claimed that the result also applied to the overlap fermions. However we do not think that this is the case. First, the effective Hamiltonian obtained there is local. We expect that by integrating over heavy modes of the domain-wall fermions there appear nonlocal terms in the effective Hamiltonian or action of the light fermions. Therefore it is suspected that the effect of heavy fermions is not properly taken into account in the calculation in Ref.. Second, in order to prove the “equivalence” between the domain-wall and overlap fermions, bosonic Pauli-Villars fields must be introduced in the domain-wall fermion formalism . As a result of the introduction of the Pauli-Villars fields, the GW relation is satisfied by the overlap Dirac operator in the action of the light fermions. However in Ref., the bosonic Pauli-Villars fields are not included at all. From the above reasons, we do not think that the result in Ref. is applicable to the overlap fermions. Actually, the effective action obtained in this paper by integrating over the gauge fields is nonlocal even in terms of the chiral-covariant fields $`q(n)`$ and $`\overline{q}(n)`$ and also there exists the extended chiral symmetry for vanishing quark mass which is regarded as axial symmetry for small hopping parameter $`t`$. |
warning/0001/cond-mat0001274.html | ar5iv | text | # 1 The Problem
## 1 The Problem
One of the great successes in quantum field theory (QFT) in the last two decades has been the realisation of the unification of the disparate forces of quantum physics through spontaneous symmetry breaking. Even without a direct observation of the Higgs meson for electroweak unification, to date, supplementary evidence in support of the Standard Model of Glashow, Salam and Weinberg suggests that it is only a matter of time before Higgs particles will be found.
Analogy with condensed matter physics would suggest that this and other symmetries were not always broken but that, in the very early, and very hot universe, they were restored. With several levels of unification anticipated we expect several changes of phase to have occurred, sequential partial orderings of an initially disordered state. Although the effects of most of the transitions can only be inferred indirectly, the last change of phase relevant to particle physics, the hadronisation of the quark-gluon plasma, is accessible at the heavy-ion collider facility at Brookhaven (RHIC), just about to begin to produce data. Although this transition is too difficult to address in this article, it provides a huge incentive to understand the way in which phase transitions occur in QFT.
Consider a situation in which the symmetry group of the theory is broken, on cooling through the critical temperature $`T_c`$, by the degenerate groundstate manifold of the order-parameter fields. Using the tools of equilibrium thermal field theory (TFT) we can determine the nature of the transition. At late times after the transition the fields are ordered on large scales, in that they adopt a single value from this degenerate set over a large spatial region.
To understand how this is achieved requires that we go beyond equilibrium TFT. In practice, we often know remarkably little about the dynamics of thermal systems. In particular, we shall restrict our discussion to the onset of such phase transitions, the very early times after the implementation of a transition when the scalar fields are only just beginning to become ordered.
A simple question to ask is the following: In principle, the field correlation length diverges at a continuous transition. In practice, it does not. What happens? This is relevant for transitions that leave topological defects like walls, monopoles, vortices, or textures in their wake since we might expect ’defects’ to be just that, entities whose separation is characterised by the correlation length. If this were simply so, an understanding of correlation lengths would lead to a prediction for defect densities. Conversely, a measurement of defect densities would be a measurement of correlation lengths. Vortices, in particular (cosmic strings), can be important for structure formation in the universe. Because of their implications for astrophysics, estimates of early field ordering in their presence have been made by Kibble, using simple causal arguments.
There are great difficulties in converting such predictions for the early universe into experimental observations. Zurek suggested that similar arguments were applicable to condensed matter systems for which direct experiments on defect densities, in particular vortex densities, could be performed. This has lead to a burst of activity from condensed matter experimentalists developing experiments to test these predictions. To date, agreement with the Zurek predictions for a variety of condensed matter systems is good, but not total.
The problems are twofold, in that we need a better understanding of the predictions and of the experiments. In this article we shall only consider the first. Many of our observations have been discussed in more detail elsewhere, and our aim here is to give an overview. However, there is so much work in progress in this area that we do not claim to be exhaustive.
## 2 The Predictions
The predictions of Kibble and Zurek for the way in which the fields freeze in at the onset of a transition are simply expressed in terms of causality. As such, they are generic. However, to quantify them it is convenient to consider simple field theories in which they can be realised.
In principle, the simplest transitions are temperature quenches. In practice, the circumstances in which transitions occur are poorly understood. In the early universe we do not necessarily begin from a thermal distribution. More awkwardly, in condensed matter physics (and heavy-ion collisions) we do not have a spatially uniform temperature. However, for the sake of argument we assume a uniform temperature, and leave more realistic scenarios until later.
There is a further problem. In condensed matter systems the nature of the fields in the Landau-Ginzburg effective theory is prescribed, and their number is small. In the early universe we have little constraint on our model making beyond a feeling that, despite the ubiquity of gauge fields at later times, at early times scalar fields are important.
### 2.1 The models
As we shall see later, most experiments involve vortices and, with this in mind, we adopt the simplest theory, that of a single complex scalar field $`\varphi `$ in three spatial dimensions with a wine-bottle potential. The transition is continuous<sup>2</sup><sup>2</sup>2There is no evidence for strong first order, or discontinuous, transitions in QFT. The mechanism for such transitions, bubble nucleation, would lead to a very different onset than that considered here. .
That is, we assume that the qualitative dynamics are conditioned by the field’s equilibrium free energy, of the form
$$F(T)=d^3x\left(|\varphi |^2+m^2(T)|\varphi |^2+\lambda |\varphi |^4\right).$$
(2.1)
Intentionally, with the early universe in mind we have written $`F(T)`$ of Eq.2.1 in the form appropriate for a relativistic quantum field. The coefficient $`m^2(T)`$ has the interactions with the real particles of the heat bath taken into account, and vanishes at $`T=T_c`$. We are considering the $`\varphi `$ field as an open system, in which changes in the external environment lead to changes in the parameter $`m^2(T)`$, which takes the value $`M^2`$ when $`T=0`$. For relativistic QFT the change in temperature that leads to the change in the sign of $`m^2`$ is most simply understood as a consequence of the system expanding. Thus, in the early universe, once thermalisation is possible, a weakly interacting relativistic plasma at temperature $`TM`$ has an entropy density $`sT^3`$. As long as thermal equilibrium can be maintained, constant entropy $`S`$ per comoving volume, $`Ssa(t)^3`$, gives $`Ta(t)^1`$ and falling, for increasing scale factor $`a(t)`$.
Specifically, at time $`t`$ the temperature $`T(t)`$ satisfies
$$tT(t)^2=m_P\left(\frac{45}{16\pi ^3N^{}}\right)^{1/2}$$
(2.2)
where $`m_P`$ is the Planck mass, and $`N^{}`$ is the number of effective field degrees of freedom. Models that attempt to take inflation into account, however, lead to ’preheating’ that is not Boltzmannian. Nonetheless, even in such cases it is possible to isolate an effective temperature for long-wavelength modes. This is all that is necessary, but is too sophisticated for the simple scenarios that we shall present here. Even the inclusion of an FRW metric would complicate the issue at this stage, and we assume flat space-time, with decreasing temperature. Prior to the transition, we assume a uniform temperature $`T>T_c`$, for which $`m^2(T)>0`$. After the transition, $`m^2(0)=M^2<0`$ enforces the $`U(1)`$ symmetry-breaking, with field expectation values $`\varphi =\pm \eta `$, $`\eta ^2=M^2/\lambda `$. The Compton wavelength $`\xi _0=M^1`$ is the natural distance scale. For the sake of argument we assume mean-field behaviour $`m^2(T)=M^2(T/T_c1)`$, whereby the equilibrium correlation length $`\xi _{eq}(T)=|m(T)|^1=\xi _0(T/T_c1)^{1/2}`$.
Equally well, after rescaling, $`F`$ could be the Ginzburg-Landau free energy
$$F(T)=d^3x\left(\frac{\mathrm{}^2}{2m}|\varphi |^2+\alpha (T)|\varphi |^2+\beta |\varphi |^4\right)$$
(2.3)
for a non-relativistic condensed matter field, in which the chemical potential $`\alpha (T)=\alpha _0(T/T_c1)`$ vanishes again at the critical temperature $`T_c`$. In this case we envisage the change in $`\alpha (T)`$ as a consequence of an external cooling of the system or a change in the pressure of the system that leads to a change in $`T_c`$. In either case, we again assume circumstances in which, in a finite time, $`T/T_c`$ varies from greater than unity to less. The fundamental length scale $`\xi _0`$ is given from Eq.2.3 as $`\xi _0^2=\mathrm{}^2/2m\alpha _0`$ whereby $`\xi _{eq}(T)=\xi _0(T/T_c1)^{1/2}`$ as before. The Gross-Pitaevski theory suggests a natural time-scale $`\tau _0=\mathrm{}/\alpha _0`$. When, later, we adopt the time-dependent Landau-Ginzburg (TDLG) theory we find this still to be true, empirically, at order-of-magnitude level, and we keep it.
The minima of the final potential now constitute the circle $`\varphi =\eta e^{i\alpha }`$, where $`\eta ^2=M^2/\lambda `$ for Eq.2.1, and $`\eta ^2=\alpha _0/\beta `$ for Eq.2.3. When the transition begins $`\varphi `$ begins to fall into the valley of the potential, choosing a random phase. Although this randomly chosen phase will vary from point to point we expect domains across which the phase is roughly constant. How this collapse takes place determines the size of the first identifiable domains. It was suggested by Kibble and Zurek that this size is essentially the equilibrium field correlation length $`\xi _{eq}`$ at some appropriate temperature close to the transition. Two very different mechanisms have been proposed for estimating this size.
### 2.2 A first guess: Thermal activation
In early work on transitions it was assumed that initial domain size was fixed in the Ginzburg regime, identifiable from equilibrium theory. By this we mean the following. Suppose the temperature $`T(t)`$ varies sufficiently slowly with time $`t`$ that it makes sense to replace $`V(\varphi ,T)`$ by $`V(\varphi ,T(t))`$. Well away from the transition this is justified, but close to the transition it is not. Once we are below $`T_c`$, and the central hump in $`V(\varphi ,T(t))`$ is forming, the Ginzburg temperature $`T_G<T_c`$ signals the temperature above which there is a significant probability for thermal fluctuations over the central hump on the scale of the correlation length at that temperature. Most simply, it is determined by the condition
$$\mathrm{\Delta }V(T_G)\xi _{eq}^3(T_G)k_BT_G$$
(2.4)
where $`\mathrm{\Delta }V(T)`$ is the difference between the central maximum and the minima of $`V(\varphi ,T)`$. That is,
$$\xi _{eq}(T_G)=O\left(\frac{\xi _0}{(1T/T_c)^{1/2}}\right),$$
(2.5)
Whereas, above $`T_G`$ there will be a population of ’domains’, fluctuating in and out of existence, at temperatures below $`T_G`$ fluctuations from one minimum to the other become increasingly unlikely. For the relativistic theory of Eq.2.1, in units in which $`k_B`$ is unity, we find $`|1T_G/T_c|=O(\lambda )`$, whereby
$$\xi _{eq}(T_G)=O\left(\frac{1}{\lambda T_c}\right)$$
(2.6)
On the other hand, for non-relativistic condensed matter, we find
$$1T_G/T_c=(\beta k_BT_c/\alpha _0^2\xi _0^3)^2$$
(2.7)
It was originally suggested by Kibble that we identify $`\xi _{eq}(T_G)`$ with the scale at which stable domains begin to form. We shall show this to be incorrect, for quenches that are not too slow. However, we will find that strong thermal fluctuations do have a role, particularly in $`{}_{}{}^{4}He`$, for which the whole experiments take place within the Ginzburg regime. This is an issue that requires more than equilibrium physics. The most simple dynamical arguments invoke causality.
### 2.3 Causality in QFT and condensed matter
The first application of causality was again due to Kibble. Whatever the details of the microscopic physics, causality sets an upper limit over which the field can be correlated in the causal horizon with diameter $`r(t)2t`$. In the vicinity of the transition, whether at $`T=T_G`$ or not, this gives a correlation length $`\xi (t)r(t)`$ or, from Eq.2.2
$$\xi \left(\frac{45}{4\pi ^3N^{}}\right)^{1/2}\frac{m_P}{T_c^2}$$
(2.8)
where $`N^{}`$ the number of field degrees of freedom at $`T_c`$. This was used to bound monopole density in the early universe. Although the lack of correlation in the field phase in different causal horizons demonstrates the existence of a field structure that will naturally lead to defects, beyond that it is not a helpful guide.
If this invocation of causality, or the Ginzburg criteria, attempt to set scales once the critical temperature has been passed, other causal arguments attempt to set scales before it is reached. Again suppose that $`T/T_c`$ varies in time as a consequence of the change in the environment. We have seen that $`\xi _{eq}(T(t))`$, obtained by inserting the time-dependence of $`T/T_c`$ explicitly, diverges at $`T(t)=T_c`$, which we suppose happens at $`t=0`$. This cannot be the case for the true correlation length $`\xi (t)`$, which can only grow so far in a finite time. Initially, for $`t<0`$, when we are far from the transition, we again assume effective equilibrium, and the field correlation length $`\xi (t)`$ tracks $`\xi _{eq}(T(t))`$ approximately. However, as we get closer to the transition, $`\xi _{eq}(T(t))`$ begins to increase arbitrarily fast and the adiabatic approximation breaks down. This is the largest value that $`\xi _{eq}`$ attains. This largest value prior to the transition corresponds to the smallest value of the field correlation length after the transition. We shall argue later that the argument is too simple but, nonetheless it is a plausible starting point. The question then becomes: how large does the field correlation get in practice? Whether in QFT or condensed matter physics there is a maximum speed at which $`\xi (t)`$ can grow on purely causal grounds. As a crude upper bound, in QFT the true correlation length $`\xi (t)`$ fails to keep up with $`\xi _{eq}(T(t))`$ by the time $`\overline{t}`$ at which $`\xi _{eq}`$ is growing at c=1, the speed of light, $`d\xi _{eq}(T(\overline{t}))/dt=1`$. It was suggested, again by Kibble, that once we have reached this time $`\xi (t)`$ freezes in, remaining approximately constant until the time $`t+\overline{t}`$ after the transition when it once again becomes comparable to the now decreasing value of $`\xi _{eq}`$. The correlation length $`\overline{\xi }=\xi _{eq}(\overline{t})=\xi _{eq}(\overline{t})`$ is argued to provide the scale for the minimum domain size after the transition.
Specifically, if we assume a time-dependence $`m^2(t)=M^2t/t_Q`$ in the vicinity of $`t=0`$, when the transition begins to be effected, then the causality condition gives $`t_C=t_Q^{1/3}(2M)^{2/3}`$. As a result,
$$M\xi _{eq}(\overline{t})=(M\tau _0)^{1/3},$$
(2.9)
which we write as
$$\overline{\xi }=\xi _{eq}(\overline{t})=\xi _0\left(\frac{\tau _Q}{\tau _0}\right)^{1/3}$$
(2.10)
where $`\tau _0=\xi _0=M^1`$ are the natural time and distance scales. In contrast to Eq.2.5, Eq.2.10 depends explicitly on the quench rate, as we would expect. For $`\tau _Q\tau _0`$ the field is correlated on a scale of many Compton wavelengths.
This approach of Kibble was one of the motivations for a similar analysis by Zurek of transitions with scalar order parameters in condensed matter. Explicitly, in the Ginzburg-Landau free energy Eq.2.3, $`\alpha (T)`$ also vanishes at the critical temperature $`T_c`$. The only difference is that, in the causal argument, the speed of light should be replaced by the speed of (second) sound, with different critical index.
Explicitly, let us again assume the mean-field result $`\alpha (T)=\alpha _0ϵ(T)`$, where $`ϵ=(T/T_c1)`$, remains valid as $`T/T_c`$ varies with time $`t`$ as $`\alpha (t)=\alpha (T(t))=\alpha _0t/\tau _Q`$ in the vicinity of $`T_c`$. It follows that the equilibrium correlation length $`\xi _{eq}(t)`$ and the relaxation time $`\tau (t)`$ diverge when $`t`$ vanishes as
$$\xi _{eq}(t)=\xi _0\left|\frac{t}{\tau _Q}\right|^{1/2},\tau (t)=\tau _0\left|\frac{t}{\tau _Q}\right|^1,$$
(2.11)
in terms of $`\xi _0`$ and $`\tau _0`$ as given earlier.
The speed of sound is $`c(t)=\xi _{eq}(t)/\tau (t)`$, slowing down as we approach the transition as $`|t|^{1/2}`$. The causal counterpart to $`d\xi _{eq}(t)/dt=1`$ for the relativistic field is $`d\xi _{eq}(t)/dt=c(t)`$. This is satisfied at $`t=\overline{t}`$, where $`\overline{t}=\sqrt{\tau _Q\tau _0}`$, with corresponding correlation length
$$\overline{\xi }=\xi _{eq}(\overline{t})=\xi _{eq}(\overline{t})=\xi _0\left(\frac{\tau _Q}{\tau _0}\right)^{1/4}.$$
(2.12)
(cf. Eq.2.10). A variant of this argument that gives essentially the same results is obtained by comparing the quench rate directly to the relaxation rate of the field fluctuations. We stress that, yet again, the assumption is that the length scale that determines the initial correlation length of the field freezes in before the transition begins. Whatever, the field is already correlated on a scale of many Compton wavelengths when it begins to unfreeze.
### 2.4 Experimental predictions
The end result of the simple causality arguments is that, both for QFT and condensed matter, when the field begins to order itself its correlation length has the form
$$\overline{\xi }=\xi _0\left(\frac{\tau _Q}{\tau _0}\right)^\gamma .$$
(2.13)
for appropriate $`\gamma `$. In fact, the powers of Eq.2.10 and Eq.2.12 are mean-field results, changed on implementing the renormalisation group. Whereas, for $`{}_{}{}^{3}He`$, the critical behaviour of Eq.2.12 survives, for $`{}_{}{}^{4}He`$, $`\gamma =1/3`$. On the other hand, for QFT the relevant parameter is $`\overline{\xi }T_c`$, the ratio of the maximum correlation length to the thermal length $`\beta _c=T_c^1`$. In equilibrium theory, when the temperature is high enough the theory is essentially three-dimensional, with critical index $`\gamma `$ different from its four-dimensional mean-field value. At our level of discussion, it is sufficient to keep (four-dimensional) mean-field values.
As we said earlier, when the transition begins $`\varphi `$ begins to fall into the valley of the potential, choosing a random phase. This randomly chosen phase will vary from point to point leading to approximate domains across which the phase is roughly constant. Such domains will meet at defects, in this case vortices, tubes of ’false’ vacuum $`\varphi 0`$, around which the field phase changes by $`\pm 2\pi `$. In an early universe context these are ’cosmic strings’, but we shall not consider their properties here.
Correlation lengths in the early universe are not amenable to direct observation. Kibble made a second assumption, that the correlation length Eq.2.10 also sets the scale for the typical minimum intervortex distance at the time that vortices are produced.
That is, the initial vortex density $`n_{def}`$ is
$$n_{def}=O\left(\frac{1}{\overline{\xi }^2}\right)=\frac{1}{f^2\xi _0^2}\left(\frac{\tau _0}{\tau _Q}\right)^{2\gamma },$$
(2.14)
for $`\gamma =1/3`$, where $`f=O(1)`$ estimates the fraction of defects per ’domain’. Equivalently, the length of vortices in a box volume $`v`$ is $`O(n_{def}v)`$.
What is striking and suspect about these predictions is that they are universal. They do not use any information about the strength ($`\lambda `$ or $`\beta `$) of the interactions, and hence the magnitude of the order parameter after the transition or, in consequence, the existence of the Ginzburg regime.
Since $`\xi _0`$ also measures cold vortex thickness, $`\tau _Q\tau _0`$ corresponds to a measurably large number of widely separated vortices. For the early universe Kibble deduced
$$\overline{\xi }=\left(\frac{m_P}{\sqrt{N^{}}M^2T_c^2}\right)^{1/3},$$
(2.15)
where, as before, $`m_P`$ is the Planck mass, and $`N^{}`$ the number of field degrees of freedom at $`T_c`$<sup>3</sup><sup>3</sup>3We note that, despite their different origins, the predictions 2.6,2.8,2.15 may not differ so hugely from one another for very small coupling.. If it could be argued that this initial network behaves classically then, thereafter, the density will reduce due to the collapse of small loops, intersections chopping off loops which in turn collapse, and vortex straightening so as to reduce the gradient energy of the field.
However, even if cosmic strings were produced in so simple a way in the very early universe it is still not possible to compare the density Eq.2.14 with experiment. What is more amenable to experiment, in principle, is the length distributions of string networks, and their ability to show scaling behaviour. This only impinges indirectly on the correlation length of the field, but would have to be commensurate with any density calculations.
It was Zurek who first suggested that, if this relationship Eq.2.14 between defect density and correlation length were true, it could be tested directly in condensed matter systems, particularly in liquid helium.
## 3 Experiments in Condensed Matter
### 3.1 Vortices in superfluid helium
Vortex lines in both superfluid $`{}_{}{}^{4}He`$ and $`{}_{}{}^{3}He`$ are analogues of global cosmic strings. A crude but effective model is to treat the system as composed of two fluids, the normal fluid and the superfluid, which has zero viscosity. In $`{}_{}{}^{4}He`$ the bose superfluid is characterised by a complex field $`\varphi `$, whose squared modulus $`|\varphi |^2`$ is the superfluid density. The superfluid fraction is unity at absolute zero, falling to zero as the temperature rises to the lambda point at 2.17K. The Landau-Ginzburg theory for $`{}_{}{}^{4}He`$ has, as its free energy $`F(T)`$ of Eq.2.3.
The situation is more complicated, but more interesting, for $`{}_{}{}^{3}He`$, which becomes superfluid at the much lower temperature of 2 mK. The reason is that the $`{}_{}{}^{3}He`$ is a fermion. Thus the mechanism for superfluidity is very different from that of $`{}_{}{}^{4}He`$. Somewhat as in a BCS superconductor, these fermions form the counterpart to Cooper pairs. However, whereas the (electron) Cooper pairs in a superconductor form a $`{}_{}{}^{1}S`$ state, the $`{}_{}{}^{3}He`$ pairs form a $`{}_{}{}^{3}P`$ state. The order parameter $`A_{\alpha i}`$ is a complex $`3\times 3`$ matrix $`A_{\alpha i}`$. There are two distinct superfluid phases, depending on how the $`SO(3)\times SO(3)\times U(1)`$ symmetry is broken. If the normal fluid is cooled at low pressures, it makes a transition to the $`{}_{}{}^{3}HeB`$ phase, in which $`A_{\alpha i}`$ takes the form $`A_{\alpha i}=R_{\alpha i}(\omega )e^{i\mathrm{\Phi }}`$, where $`R`$ is a real rotation matrix, corresponding to a rotation through an arbitrary $`\omega `$ <sup>4</sup><sup>4</sup>4At large distance scales there is a complication in that the small spin-orbit coupling becomes important, to fix $`\omega `$ at $`arccos(1/4)`$, but this will not concern us here..
The Landau-Ginzburg free energy is, necessarily, more complicated, but the effective potential $`V(A_{\alpha i},T)`$ has the diagonal form $`V(A,T)=\alpha (T)|A_{ai}|^2+O(A^4)`$ for small fluctuations, and this is all that we need for the production of vortices at very early times. Beyond that it can be mimicked by Eq.2.3 for our purposes (e.g. see ). Thus the Zurek analysis leads to the prediction Eq.2.14, as before, for appropriate $`\gamma `$. However, for $`{}_{}{}^{3}He`$ the mean-field approximation is good and the mean-field critical index $`\gamma =1/4`$ is not renormalised, whereas for $`{}_{}{}^{4}He`$ a better value is $`\gamma =1/3`$, as for the naive relativistic theory.
### 3.2 Experiments in $`{}_{}{}^{3}He`$.
Although $`{}_{}{}^{3}He`$ is more complicated to work with, the experiments to check Eq.2.14 are cleaner in that, because the nucleus has spin $`1/2`$, even individual vortices can be detected by magnetic resonance. Further, because vortex width is many atomic spacings the Landau-Ginzburg theory is reliable.
So far, experiments have been of two types. In the Helsinki experiment superfluid $`{}_{}{}^{3}HeB`$ in a rotating cryostat is bombarded by slow neutrons. Each neutron entering the chamber releases 760 keV, via the reaction $`n+^3Hep+^3He+760keV`$. The energy goes into the kinetic energy of the proton and triton, and is dissipated by ionisation, heating a region of the sample above its transition temperature. The heated region then cools back through the transition temperature, creating vortices. Vortices above a critical size (dependent on the angular velocity of the cryostat) grow and migrate to the centre of the apparatus, where they are counted by an NMR absorption measurement. Suffice to say that the quench is very fast, with $`\tau _Q/\tau _0=O(10^3)`$. Agreement with Eq.2.14 and Eq.2.12 is very good, at the level of less than an order of magnitude. This is even though it is now argued that the Helsinki experiment should not show agreement because of the geometry of the heating event.
The second type of experiment has been performed at Grenoble and Lancaster. Rather than count individual vortices, the experiment detects the total energy going into vortex formation. As before, $`{}_{}{}^{3}He`$ is irradiated by neutrons. After each absorption the energy released in the form of quasiparticles is measured, and found to be less than the total 760 keV. This missing energy is assumed to have been expended on vortex production. Again, agreement with Zurek’s prediction Eq.2.14 and Eq.2.12 is good.
### 3.3 Experiments in $`{}_{}{}^{4}He`$.
The experiments in $`{}_{}{}^{4}He`$, conducted at Lancaster, follow Zurek’s original suggestion. The idea is to expand a sample of normal fluid helium, in a container with bellows, so that it becomes superfluid at essentially constant temperature. That is, we change $`1T/T_c`$ from negative to positive by reducing the pressure, thereby increasing $`T_c`$. As the system goes into the superfluid phase a tangle of vortices is formed, because of the random distribution of field phases. The vortices are detected by measuring the attenuation of second sound within the bellows. Second sound scatters off vortices, and its attenuation gives a good measure of vortex density. A mechanical quench is slow, with $`\tau _Q`$ some tens of milliseconds, and $`\tau _Q/\tau _0=O(10^{10})`$. Two experiments have been performed. In the first fair agreement was found with the prediction Eq.2.14, although it was not possible to vary $`\tau _Q`$. However, there were potential problems with hydrodynamic effects at the bellows, and at the capillary with which the bellows were filled. A second experiment, designed to minimise these and other problems has failed to see any vortices whatsoever.
Some care is needed. Not only is the Landau-Ginzburg effective theory more suspect for $`{}_{}{}^{4}He`$, but its Ginzburg regime is so wide, at $`O(1K)`$ that the transition takes place entirely within it.
### 3.4 Quenching in an annulus
Experiments for vortex densities are problematical in that a closely bound vortex-antivortex pair gives a count of 2 right up to annihilation, whereas their topological charge remains zero. It is possible to devise experiments that count topological charge. Consider a closed path in the bulk superfluid with circumference $`C\xi (t)`$. Naively, the number of ’regions’ through which this path passes in which the phase is correlated is $`𝒩=O(C/\xi (t))`$. Assuming an independent choice of phase in each ’region’, the r.m.s phase difference along the path is
$$\mathrm{\Delta }\theta _C\sqrt{𝒩}=O(\sqrt{C/\overline{\xi }}).$$
(3.16)
If we now consider a quench in an annular container of similar circumference $`C`$ of superfluid $`{}_{}{}^{4}He`$ and radius $`lC`$, Zurek suggested that the phase locked in is also given by Eq.3.16, with $`\overline{\xi }`$ of Eq.2.12. Since the phase gradient is directly proportional to the superflow velocity we expect a flow after the quench with r.m.s velocity
$$\mathrm{\Delta }v=O\left(\frac{\mathrm{}}{m}\sqrt{\frac{1}{C\overline{\xi }}}\right).$$
(3.17)
provided $`l=O(\overline{\xi })`$. Although in bulk fluid this superflow will disperse, if it is constrained to a narrow annulus it should persist, and although not large is measurable, in principle. Specifically, in the units of Zurek
$$\mathrm{\Delta }v=O((\tau _Q[\mu s])^{\nu /4}/\sqrt{C[cm]})$$
(3.18)
where $`\tau _Q`$ is, typically, tens of milliseconds and $`C`$ in centimetres. $`\nu =1/2`$ is the mean-field critical exponent above. In principle $`\nu `$ should be renormalised to $`\nu =2/3`$, but the difference to $`\mathrm{\Delta }v`$ is negligable. With such a small index the result is almost independent of quench rate.
In practice there are difficulties in performing annular measurements in $`{}_{}{}^{4}He`$. A similar, but easier, experiment can be performed on annular Type-II superconductors, on cooling through their critical temperatures. The relevant free energy is the extension of $`F`$ of Eq.2.3,
$$F(T)=d^3x\left(\frac{1}{4m}|i\mathrm{}\varphi \frac{2e}{c}𝐀|^2+\alpha (T)|\varphi |^2+\frac{1}{4}\beta |\varphi |^4\right)+\frac{B^2}{8\pi }.$$
(3.19)
A is the vector potential in the Coulomb gauge, and $`𝐁=𝐀`$.
If we can initially ignore the effects of the gauge field as a temperature quench is imposed through a change in $`\alpha (T)`$, as before, the result Eq.3.16 persists. For perimeter $`C`$ the variance in the number of flux quanta produced spontaneously is
$$\mathrm{\Delta }N_C=\frac{1}{2\pi }\mathrm{\Delta }\theta _C\frac{1}{2\pi }\sqrt{\frac{C}{\overline{\xi }}}.$$
(3.20)
In fact, it is more convenient to quench annular Josephson Junctions, in which two identical rings of superconductor are held apart by an oxide layer through which Cooper pairs can tunnel. If $`\theta _1`$ and $`\theta _2`$ are the phases of $`\varphi `$ on the upper and lower rings then, once the transition has taken place the tunneling current has the form
$$J=J_csin(\theta _1\theta _2)$$
(3.21)
and the theory is described by a dissipative Sine-Gordon equation.
The kinks of this equation are the ’fluxons’ of the Josephson Junction and are easy to observe experimentally. The variance in fluxon number at their formation is
$$\mathrm{\Delta }N_C=\frac{1}{2\pi }\mathrm{\Delta }(\theta _1\theta _2).$$
(3.22)
Since $`\theta _1,\theta _2`$ are independent as the transition begins, we would expect, from Zurek’s analysis, that
$$\mathrm{\Delta }\theta _1=\mathrm{\Delta }\theta _2\sqrt{\frac{C}{\overline{\xi }}},$$
(3.23)
whence
$$\mathrm{\Delta }N_C\frac{1}{2\pi }\sqrt{\frac{2C}{\overline{\xi }}}.$$
(3.24)
Although some caution is required in the interpretation of the experiments, which were not devised with this prediction in mind, it seems to be supported by experiment, for which, with $`\overline{\xi }=O(10^1)mm`$ <sup>5</sup><sup>5</sup>5Corresponding to $`\tau _Q=O(1)s`$, $`\xi _0=O(10^3)A`$. and $`C=0.5\times 10^1mm`$, say, we would expect $`\mathrm{\Delta }N_c=O(1)`$.
There is certainly no agreement, in this or any other experiment in which defects are seen, with the thermal fluctuation density that would be based on Eq.2.5.
## 4 The Kibble-Zurek picture for the freezing in of $`\xi `$ is correct.
Since all equations of motion have causality built into them we should be able to confirm the first predictions Eq.2.10 and 2.12 of Kibble and Zurek explicitly, as we shall now see.
### 4.1 Condensed matter: the TDLG equation
We assume that, for the condensed matter systems of interest to us, the dynamics of the transition can be derived from the explicitly time-dependent Landau-Ginzburg free energy
$$F(t)=d^3x\left(\frac{\mathrm{}^2}{2m}(\varphi _a)^2+\alpha (t)\varphi _a^2+\frac{1}{4}\beta (\varphi _a^2)^2\right).$$
(4.25)
in which we substitute $`T(t)`$ for $`T`$ directly in (2.3). In (4.25) $`\varphi =(\varphi _1+i\varphi _2)/\sqrt{2}`$ ($`a=1,2`$) is the complex order-parameter field, whose magnitude determines the superfluid density. As before, in a mean field approximation, the chemical potential $`\alpha (T)`$ takes the form $`\alpha (T)=\alpha _0ϵ(\overline{t})`$, where $`ϵ=(T/T_c1)`$. In a quench in which $`T_c`$ or $`T`$ changes it is convenient to shift the origin in time, to write $`ϵ`$ as
$$ϵ(t)=ϵ_0\frac{t}{\tau _Q}\theta (t)$$
(4.26)
for $`\mathrm{}<t<\tau _Q(1+ϵ_0)`$, after which $`ϵ(t)=1`$. $`ϵ_0=1T_0/T_c`$ measures the original temperature $`T_0`$ and $`\tau _Q`$ defines the quench rate. The quench begins at time $`t=0`$ but the transition from the normal to the superfluid phase only begins at time $`t_0=ϵ_0\tau _Q`$. When it is convenient to measure time from the onset of the transition we use the notation $`\mathrm{\Delta }t=tt_0`$.
Motivated by Zurek’s later numerical simulations, we adopt the time-dependent Landau-Ginzburg (TDLG) equation for $`F`$,
$$\frac{1}{\mathrm{\Gamma }}\frac{\varphi _a}{t}=\frac{\delta F}{\delta \varphi _a}+\eta _a,$$
(4.27)
where $`\eta _a`$ is Gaussian thermal noise, satisfying
$$\eta _a(𝐱,t)\eta _b(𝐲,t^{})=2\delta _{ab}T(t)\mathrm{\Gamma }\delta (𝐱𝐲)\delta (tt^{}).$$
(4.28)
This is a crude approximation for $`{}_{}{}^{4}He`$, and a simplified form of a realistic description of $`{}_{}{}^{3}He`$ but it is not a useful description of QFT, as it stands.
It is relatively simple to determine the validity of Zurek’s argument since it assumes that freezing in of field fluctuations occurs just before symmetry breaking begins. At that time the effective potential $`V(\varphi ,T)`$ is still roughly quadratic and we can see later that, for the relevant time-interval $`\overline{t}\mathrm{\Delta }t\overline{t}`$ the self-interaction term can be neglected ($`\beta =0`$).
In space, time and temperature units in which $`\xi _0=\tau _0=k_B=1`$, Eq.4.27 then becomes
$$\dot{\varphi }_a(𝐱,t)=[^2+ϵ(t)]\varphi _a(𝐱,t)+\overline{\eta }_a(𝐱,t).$$
(4.29)
where $`\overline{\eta }`$ is the renormalised noise. The solution of the ’free’-field linear equation is straightforward, giving a Gaussian equal-time correlation function
$$\varphi _a(𝐫,t)\varphi _b(\mathrm{𝟎},t)=\delta _{ab}G(𝐫,t)=d/^3ke^{i𝐤.𝐫}P(k,t).$$
(4.30)
in which the power spectrum $`P(k,t)`$ has a representation in terms of the Schwinger proper-time $`\tau `$ as
$$P(k,t)=_0^{\mathrm{}}𝑑\tau \overline{T}(t\tau /2)e^{\tau k^2}e^{_0^\tau 𝑑sϵ(ts/2)},$$
(4.31)
where $`\overline{T}`$ is the renormalised temperature. In turn, this gives
$$G(r,t)=_0^{\mathrm{}}𝑑\tau \overline{T}(t\tau /2)\left(\frac{1}{4\pi \tau }\right)^{3/2}e^{r^2/4\tau }e^{_0^\tau 𝑑sϵ(ts/2)}.$$
(4.32)
For constant $`ϵ`$, as happens at early times, on rescaling in Eq.4.32 we recover the usual Yukawa correlator
$$G(r,t)=\frac{T_0}{4\pi r}e^{r/\xi _0},$$
(4.33)
where $`T_0=T_c(1+ϵ_0)`$ is the initial temperature. For $`ϵ(t)`$ of Eq.4.26 a saddle-point calculation shows that, although we recover Eq.4.33 for $`r/\overline{\xi }>(\overline{\xi }/\xi _0)^3`$ at later times, the correlation function is then dominated by its smaller-r behaviour. At time $`t_0=ϵ_0\tau _0`$, when the transition begins, a saddle-point calculation shows that, provided the quench is not too fast,
$$G(r,t_0)\frac{T_c}{4\pi r}e^{a(r/\overline{\xi })^{4/3}},$$
(4.34)
where $`a=O(1)`$, confirming Zurek’s result.
Zurek’s prediction is robust, since further calculation shows that $`\xi (t)`$ does not vary strongly in the interval $`\overline{t}\mathrm{\Delta }t\overline{t}`$, where $`\mathrm{\Delta }t=tt_0`$.
### 4.2 QFT: Closed time-path ensemble averaging
For QFT the situation is rather different. In the previous section, instead of working with the TDLG equation, we could have worked with the equivalent Fokker-Planck equation for the probability $`p_t[\mathrm{\Phi }]`$ that, at time $`t>0`$, the measurement of $`\varphi `$ will give the function $`\mathrm{\Phi }(𝐱)`$. When solving the dynamical equations for a hot quantum field it is convenient to work with probabilities from the start.
Take $`t=0`$ as our starting time for the evolution of the complex field $`\varphi =(\varphi _1+i\varphi _2)/\sqrt{2}`$. Suppose that, at this time, the system is in a pure state, in which the measurement of $`\varphi `$ would give $`\mathrm{\Phi }_0(𝐱)`$. That is:-
$$\widehat{\varphi }(t=0,𝐱)|\mathrm{\Phi }_0,t=0=\mathrm{\Phi }_0|\mathrm{\Phi }_0,t=0.$$
(4.35)
The probability $`p_t[\mathrm{\Phi }]`$ that, at time $`t_f>0`$, the measurement of $`\varphi `$ will give the value $`\mathrm{\Phi }`$ is $`p_t[\mathrm{\Phi }]=|\mathrm{\Psi }_0|^2`$, where $`\mathrm{\Psi }_0`$ is the state-functional with the specified initial condition. As a path integral
$$\mathrm{\Psi }_0=_{\varphi (0)=\mathrm{\Phi }_0}^{\varphi (t)=\mathrm{\Phi }}𝒟\varphi \mathrm{exp}\left\{iS_t[\varphi ]\right\},$$
(4.36)
where $`S_t[\varphi ]`$ is the (time-dependent) action that describes how the field $`\varphi `$ is driven by the environment and spatial and field labels have been suppressed (e.g. $`𝒟\varphi =𝒟\varphi _1𝒟\varphi _2)`$. Specifically, for $`t>0`$ the action for the field is taken to be
$$S_t[\varphi ]=𝑑x\left(\frac{1}{2}_\mu \varphi _a^\mu \varphi _a\frac{1}{2}m^2(t)\varphi _a^2\frac{1}{4}\lambda (\varphi _a^2)^2\right).$$
(4.37)
where $`m(t)`$ describes the evolution of the action under external influences, to which the field responds.
It follows that $`p_t[\mathrm{\Phi }]`$ can be written in the closed time-path form
$$p_t[\mathrm{\Phi }]=_{\varphi _\pm (0)=\mathrm{\Phi }_0}^{\varphi _\pm (t)=\mathrm{\Phi }}𝒟\varphi _+𝒟\varphi _{}\mathrm{exp}\left\{i\left(S_t[\varphi _+]S_t[\varphi _{}]\right)\right\},$$
(4.38)
where $`𝒟\varphi _\pm =𝒟\varphi _{\pm ,1}𝒟\varphi _{\pm ,2}`$. Instead of separately integrating $`\varphi _\pm `$ along the time paths $`0tt_f`$, the integral can be interpreted as time-ordering of a field $`\varphi `$ along the closed path $`C_+C_{}`$ where $`\varphi =\varphi _+`$ on $`C_+`$ and $`\varphi =\varphi _{}`$ on $`C_{}`$. When we extend the contour from $`t_f`$ to $`t=\mathrm{}`$ either $`\varphi _+`$ or $`\varphi _{}`$ is an equally good candidate for the physical field, but we choose $`\varphi _+`$.
The choice of a pure state at time $`t=0`$ is too simple to be of any use. As we said earlier, we assume that $`\mathrm{\Phi }`$ is Boltzmann distributed at time $`t=0`$ at an effective temperature of $`T_0=\beta _0^1`$ according to the Hamiltonian $`H[\mathrm{\Phi }]`$ corresponding to the action $`S[\varphi ]`$, in which $`\varphi `$ is taken to be periodic in imaginary time with period $`\beta _0`$. We now have the explicit form for $`p_t[\mathrm{\Phi }]`$,
$$p_t[\mathrm{\Phi }]=_B𝒟\varphi e^{iS_C[\varphi ]}\delta [\varphi _+(t_f)\mathrm{\Phi }],$$
(4.39)
written as the time ordering of a single field along the contour $`C=C_+C_{}C_3`$, extended to include a third imaginary leg, where $`\varphi `$ takes the values $`\varphi _+`$, $`\varphi _{}`$ and $`\varphi _3`$ on $`C_+`$, $`C_{}`$ and $`C_3`$ respectively, for which $`S_C`$ is $`S[\varphi _+]`$, $`S[\varphi _{}]`$ and $`S_0[\varphi _3]`$.
To demonstrate how we can average without having to calculate $`p_t[\mathrm{\Phi }]`$ explicitly we see that $`G_{ab}(|𝐱𝐱^{}|;t)=\mathrm{\Phi }_a(𝐱)\mathrm{\Phi }_b(𝐱^{})_t`$ is given by
$$G_{ab}(|𝐱𝐱^{}|;t)=\varphi _a(𝐱,t)\varphi _b(𝐱^{},t),$$
(4.40)
the equal-time thermal Wightman function with the given thermal boundary conditions. Because of the non-equilibrium time evolution there is no time translation invariance in the double time label.
### 4.3 QFT: the free roll
Fortunately, as for the condensed matter case, the interval $`\overline{t}\mathrm{\Delta }t\overline{t}`$ occurs in the linear regime, when the self-interactions are unimportant. The relevant equation for constructing the correlation functions of this one-field system is now the second-order equation
$$\frac{^2\varphi _a}{t^2}=\frac{\delta F}{\delta \varphi _a},$$
(4.41)
for $`F`$ of Eq.2.1. This is solvable in terms of the mode functions $`\chi _k^\pm (t)`$, identical for $`a=1,2`$, satisfying
$$\left[\frac{d^2}{dt^2}+𝐤^2+m^2(t)\right]\chi _k^\pm (t)=0,$$
(4.42)
subject to $`\chi _k^\pm (t)=e^{\pm i\omega _{in}t}`$ at $`t0`$, for incident frequency $`\omega _{in}=\sqrt{𝐤^2+ϵ_0M^2}`$, for $`m^2(t)=ϵ(t)M^2`$, where $`ϵ(t)`$ is parameterised as for the TDLG equation above. This corresponds to a temperature quench from an initial state of thermal equilibrium at temperature $`T_0>T_c`$, where $`(T_0/T_c1)=ϵ_0`$. There is no reason to take $`ϵ_0`$ small. The diagonal correlation function $`G(r,t)`$ of Eq.4.30 is given as the equal-time propagator
$`G(r,t)`$ $`=`$ $`{\displaystyle d/^3ke^{i𝐤.𝐱}\chi _k^+(t)\chi _k^{}(t)C(k)}`$
$`=`$ $`{\displaystyle \frac{1}{2\pi ^2}}{\displaystyle 𝑑kk^2\frac{\mathrm{sin}kr}{kr}\chi _k^+(t)\chi _k^{}(t)C(k)},`$
where $`C(k)=\mathrm{coth}(\omega _{in}(k)/2T_0)/2\omega _{in}(k)`$ encodes the initial conditions.
An exact solution can be given in terms of Airy functions. Dimensional analysis shows that, on ignoring the k-dependence of $`C(k)`$, appropriate for large $`r`$ (or small $`k`$), $`\xi _{eq}(\overline{t})`$ of Eq.2.10 again sets the scale of the equal-time correlation function. Specifically,
$$G(r,\mathrm{\Delta }t=0)𝑑\kappa \frac{\mathrm{sin}\kappa (r/\overline{\xi })}{\kappa (r/\overline{\xi })}F(\kappa ),$$
(4.44)
where $`F(0)=1`$ and $`F(\kappa )\kappa ^3`$ for large $`\kappa `$. Kibble’s insight is correct, at least for this case of a single (uncoupled) field.
## 5 However, Defect Densities do not Determine $`\overline{\xi }`$ Directly (and Vice-Versa)
We have seen that there is no reason to disbelieve the causal arguments of Kibble for QFT and Zurek for condensed matter as to the correlation length $`\overline{\xi }`$ at the onset of the transition. The excellent agreement with the $`{}_{}{}^{3}He`$ experiments suggests that, for condensed matter, this length does, indeed, characterise vortex separation at the time when the defects form.
However, if we take the Lancaster experiment at face value, this cannot always be the case. This is not surprising. Numerical simulations of TDLG equations, by Zurek himself, and the related Langevin equations for non-equilibrium statistical fields show that, at early times after a transition, the field fluctuations remain on very small scales. In reality, different frequency modes freeze at different times, and the causal argument only applies strictly to the long wavelength modes. Of course, field ordering is controlled by the long wavelength modes, but we would like a more complete description of the freezing in of the field, if we are to make the second step of relating this to defect structure in the field . That is, although the field is correlated on long scales $`\overline{\xi }`$, this does not preclude structure on very much shorter scales. Moreover, it is this structure, that we might wish to think of as proto-defects, that will evolve into the defects, in this case vortices, once the transition is complete. Inspection of the field at early times shows no obvious relation between the separation of these proto-defects and $`\overline{\xi }`$. The question then becomes one of why the $`{}_{}{}^{3}He`$ experiments should be in agreement, rather than why the $`{}_{}{}^{4}He`$ experiments are not.
### 5.1 Classical vortices in condensed matter and QFT
The $`O(2)`$ string is, classically, a tube of false vacuum of radius $`O(M^1)`$ whose core is characterised by a line of string zeroes. Its winding number $`nZ`$ is the change in phase of the complex field around the core (in units of $`2\pi `$). We shall only consider strings with $`|n|=1`$ since all others can be considered as multiple zeroes.
It would be foolish to estimate the probability of finding vortices directly from $`p_t[\mathrm{\Phi }]`$. A starting-point for counting vortices in superfluids is to count line zeroes. Not all line zeroes are candidates for vortices since zeroes occur on all scales. However, a starting-point for counting vortices in superfluids is to count line zeroes of an appropriately coarse-grained field, in which structure on a scale smaller than $`\xi _0`$, the classical vortex size, is not present. This is also the unstated basis of the numerous numerical simulations of cosmic string networks built from Gaussian fluctuations (but see ).
Even then, there are several prerequisites before line zeroes can be identified with vortex cores, and $`n_{zero}(t)`$ with $`n_{def}(t)`$.
* The field, on average must have achieved its symmetry-broken ground-state equilibrium value
$$|\varphi |^2=\alpha _0/\beta \text{or}M^2/\lambda ,$$
(5.45)
non-perturbatively large (in $`\beta `$) This, in itself, is sufficient to show that the causal time $`\overline{t}`$ is not the time to begin looking for defects, since $`|\varphi |^2`$ is small at this time.
* Only when $`n_{zero}/l`$ is small in comparison to $`n_{zero}/l`$ at $`l=\xi _0`$ will the line-zeroes have the small-scale non-fractal nature of classical defects, although defects may behave like random walks on larger scales. As the power in the long wavelength modes increases the ’Bragg’ peak develops in $`k^2G(k,t)`$, moving in towards $`k=0`$. This condition then becomes the condition that the peak dominates its tail.
* The energy in field gradients should be commensurate with the energy in classical vortices with the same density as that of line zeroes.
We stress that these are necessary, but not sufficient, conditions for classical vortices. In particular, only the full nonlinearity of the system can establish classical profiles. We will term such zeroes as satisfy these conditions, proto-vortices. In fact, most (but not all) numerical lattice simulations cannot distinguish between proto-vortices and classical vortices.
Whereas the above are equally true for condensed matter and QFT there are further complications peculiar to QFT. In particular, in QFT we need to consider the whole density matrix $`\mathrm{\Phi }^{}|\rho (t)|\mathrm{\Phi }`$ rather than just the diagonal elements $`p_t[\mathrm{\Phi }]=\mathrm{\Phi }|\rho (t)|\mathrm{\Phi }`$. Classicality is understood in terms of ’decoherence’, manifest most simply by the approximate diagonalisation of the reduced density matrix on coarse-graining. By this we mean the separation of the whole into the ’system’, and its ’environment’ whose degrees of freedom are integrated over, to give a reduced density matrix. The environment can be either other fields with which our scalar is interacting or even the short wavelength modes of the scalar field itself . When interactions are taken into account this leads to quantum noise and dissipation.
In the Gaussian approximations for QFT that we shall adopt here, with $`\mathrm{\Phi }=0`$, integrating out short wavelengths with $`k>l^1`$ is just equivalent to a momentum cut-off at the same value. This gives neither noise nor dissipation and diagonalisation does not occur. Nonetheless, from our viewpoint of counting line-zeroes, fluctuations are still present when $`l=O(M^1)`$ that prevent us from identifying line-zeroes with proto-vortices easily.
For all these caveats, there are other symptoms of classical behaviour in QFT once $`G_l(0;t)`$ is non-perturbatively large. Instead of a field basis, we can work in a particle basis and measure the particle production as the transition proceeds. We shall see that, in the Gaussian approximation, $`n_{zero}`$, the density of line zeroes, is given in terms of the moments of $`G_l(r,t)`$. Whether we expand with respect to the original Fock vacuum or with respect to the adiabatic vacuum state, the presence of a non-perturbatively large peak in $`k^2G(k;t)`$ at $`k=k_0`$ signals a non-perturbatively large occupation number $`N_{k_0}1/\lambda `$ of particles at the same wavenumber $`k_0`$. With $`n_{zero}`$ of order $`k_0^2`$ this shows that the long wavelength modes can now begin to be treated classically.
From a slightly different viewpoint, the Wigner functional only peaks about the classical phase-space trajectory once the power is non-perturbatively large. More crudely, the diagonal density matrix elements are only then significantly non-zero for non-perturbatively large field configurations $`\varphi \lambda ^{1/2}`$ like vortices.
### 5.2 Line-zero density
Suppose, at some time, that the field has line zeroes $`𝐑_n(s)`$, where $`n=1,2,..`$ labels the zero, and $`s`$ measures the length along it. As a result the topological line density of zeroes $`\stackrel{}{\rho }(𝐫)`$ can be defined by
$$\stackrel{}{\rho }(𝐱)=\underset{n}{}𝑑s\frac{d𝐑_n}{ds}\delta ^3[𝐱𝐑_n(s)].$$
(5.46)
In (5.46) $`ds`$ is the incremental length along the line of zeroes $`𝐑_n(s)`$ ($`n`$=1,2,.. .) and $`\frac{d𝐑_n}{ds}`$ is a unit vector pointing in the direction which corresponds to positive winding number.
It follows that, in terms of the zeroes of the field $`\mathrm{\Phi }(𝐱)`$, $`\rho _i(𝐱)`$ can be written as
$$\rho _i(𝐱)=\delta ^2[\mathrm{\Phi }(𝐱)]ϵ_{ijk}_j\mathrm{\Phi }_1(𝐱)_k\mathrm{\Phi }_2(𝐱),$$
(5.47)
where $`\delta ^2[\mathrm{\Phi }(𝐱)]=\delta [\mathrm{\Phi }_1(𝐱)]\delta [\mathrm{\Phi }_2(𝐱)]`$. The coefficient of the $`\delta `$-function in (5.47) is the Jacobian of the more complicated transformation from line zeroes to field zeroes. What we want is not this, but the total line density $`\overline{\rho }(𝐱)`$,
$$\overline{\rho _i}(𝐱)=\delta ^2[\mathrm{\Phi }(𝐱)]|ϵ_{ijk}_j\mathrm{\Phi }_1(𝐱)_k\mathrm{\Phi }_2(𝐱)|.$$
(5.48)
The vanishing field expectation value and the independence of the field and its derivatives
$$\mathrm{\Phi }_a(𝐱)=0=\mathrm{\Phi }_a(𝐱)_j\mathrm{\Phi }_b(𝐱),$$
(5.49)
imply $`\rho _j(𝐱)=0`$ i.e. an equal likelihood of a string line-zero or an antistring line-zero passing through an infinitesimal area. However,
$$n(t)=\overline{\rho _i}(𝐱)_t>0$$
(5.50)
and measures the total line-zero density in the direction $`i`$, without regard to string orientation. The isotropy of the initial state guarantees that $`n`$ is independent of the direction $`i`$.
Whereas the correlation length $`\overline{\xi }=\xi _{eq}(\overline{t})`$ depends on the long-distance behaviour of $`G(r,t)`$, this is not the case for the line-zero density $`n_{zero}`$. In our Gaussian approximation of the previous section it is determined completely by the short-distance behaviour of $`G(r,t)`$ as
$`n_{zero}(t)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle \frac{G^{\prime \prime }(0,t)}{G(0,t)}}`$
$`=`$ $`{\displaystyle \frac{1}{2\pi }}f^{\prime \prime }(0,t),\text{where}f(r,t)={\displaystyle \frac{G(r,t)}{G(0,t)}}`$
Some caution is necessary. Since thermal and quantum fluctuations will give rise to zeroes on all scales, neither $`G(0,t)`$ nor its derivatives exist because of ultraviolet divergences. For the moment, we put in a cutoff $`l=O(\xi _0)`$ by hand, as
$$G_l(r,t)=d/^3ke^{i𝐤.𝐱}G(k,t)e^{k^2l^2}.$$
(5.52)
We note that the inclusion of a cut-off does not affect the long-distance correlation $`\overline{\xi }`$, which depends essentially on the weighting of the nearest singularities of $`G(k,t)`$ in the complex $`k`$-plane.
In the linear regime everything is calculable. If we define the line-zero separation $`\xi _{zero}(l)`$ by
$$n_{zero}(l,t)=\frac{1}{2\pi \xi _{zero}(l,t)^2}$$
(5.53)
it is apparent that $`\xi _{zero}(t)`$ has little, if anything, to do with $`\overline{\xi }`$ directly.
For non-Gaussian fields the situation is much more complicated. However, as long as there is a dominant wavenumber $`k_0`$ in $`G(k;t)`$ this sets a length scale $`\xi k_0^1`$ that characterises vortex separation. At the level of density calculations, as distinct from length distributions, this is essentially all that is needed.
We can already anticipate a mismatch between the defect density and the long-range correlation length that could derail the Kibble-Zurek predictions.
### 5.3 Annular experiments require a new correlation length
We conclude with a brief discussion of annular experiments, that have no counterpart in QFT. To the extent that an annulus can be treated as a one-dimensional system there has been good numerical work. However, to date there are no reliable calculations for quenching in an annulus that take realistic boundary conditions into account. Nonetheless, we have some hints as how to proceed. Consider a circular path in the bulk fluid (in the 1-2 plane), circumference $`C`$, the boundary of a surface $`S`$. For given field configurations $`\varphi _a(𝐱)`$ the phase change $`\theta _C`$ along the path can be expressed as the surface integral
$$\theta _C=2\pi _{𝐱S}d^2x\rho _3(𝐱),$$
(5.54)
Again we quench from an initial state with no rotation. The variance in the phase change around $`C`$, $`\mathrm{\Delta }\theta _C`$ is determined from
$$(\mathrm{\Delta }\theta _C)^2=4\pi ^2_{𝐱S}d^2x_{𝐲S}d^2y\rho _3(𝐱)\rho _3(𝐲)_t.$$
(5.55)
On using the conservation of charge it is not difficult to show, from the results of that, for Gaussian fields, with momentum cut-off at $`k=l^1`$, $`\mathrm{\Delta }\theta _C`$ satisfies
$$(\mathrm{\Delta }\theta _C)^2=_{𝐱S}d^2x_{𝐲S}d^2y𝒞_l(|𝐱𝐲|,t),$$
(5.56)
where $`𝐱`$ and $`𝐲`$ are in the plane of $`S`$, and
$$𝒞_l(r,t)=\frac{1}{r}\frac{}{r}\left(\frac{f_l^2(r,t)}{1f_l^2(r,t)}\right).$$
(5.57)
Since $`G_l(r,t)`$ is short-ranged $`𝒞_l(r,t)`$ is short-ranged also. With $`𝐱`$ outside $`S`$, and $`𝐲`$ inside $`S`$, all the contribution to $`(\mathrm{\Delta }\theta _C)^2`$ comes from the vicinity of the boundary of $`S`$, rather than the whole area. That is, if we removed all fluid except for a strip from the neighbourhood of the contour $`C`$ we would still have the same result. This supports the assertion by Zurek that the correlation length for phase variation in bulk fluid is also appropriate for annular flow. The purpose of the annulus (more exactly, a circular capillary of circumference $`C`$ with radius $`lC`$) is to stop this flow dissipating into the bulk fluid.
More precisely, suppose that $`C\xi (t)`$. Then, if we take the width $`2l`$ of the strip around the contour to be larger than the correlation length of $`𝒞_l(r,t)`$, Eq.5.56 can be written as
$$(\mathrm{\Delta }\theta _C)^22C_0^{\mathrm{}}𝑑rr^2𝒞_l(r,t).$$
(5.58)
The linear dependence on $`C`$ is purely a result of Gaussian fluctuations.
Insofar as we can identify the bulk correlation with the annular correlation, instead of Eq.3.17, we have
$$\mathrm{\Delta }v=\frac{\mathrm{}}{m}\sqrt{\frac{1}{C\xi _s(t)}}.$$
(5.59)
The step length $`\xi _s(t)`$ is given by
$$\frac{1}{\xi _s(t)}=2_0^{\mathrm{}}𝑑r\frac{f_l^2(r,t)}{1f_l^2(r,t)}.$$
(5.60)
If we quench in an annular capillary of radius $`l`$ much smaller than its circumference, we are, essentially, coarsegraining to that scale. That is, the observed variance in the flux along the annulus is $`\pi l^2\mathrm{\Delta }v`$ for $`\mathrm{\Delta }v`$ averaged on a scale $`l`$. We make the approximation that that is the major effect of quenching in an annulus. This cannot be wholly true, but it is plausible if the annulus is not too narrow for boundary effects to be important.
## 6 Nonetheless, Zurek’s Predictions are Correct When Thermal Fluctuations are Small
An integral part of the Zurek/Kibble predictions for $`n_{def}`$ is their universality. The ratio Eq.5.2 for $`n_{zero}`$ gives intimations as to how this could happen, insofar as the the specific effects of the interactions can be subsumed into prefactors that cancel. Although, in the context of Gaussian fluctuations, it can only be approximate, it is part of a general truth that defect density only depends on very limited attributes of the fluctuation power spectra.
We can see this in the models that we have introduced. The TDGL equation is highly dissipative, whereas the QFT equation is not, as it stands. Nonetheless, in each case the field will become ordered after the onset of the transition by the growth of unstable long wavelength modes. How this happens depends on the specifics of the equations. We would therefore expect the field behaviour as the transition is being implemented to differ in detail for condensed matter and QFT, but only weakly affecting the initial line-zero density, and it s ability to describe vortices.
### 6.1 TDLG condensed matter: vortex densities
We begin with condensed matter, which we will find to be easier. As the system evolves away from the transition time, the free equation Eq.4.29 ceases to be valid, to be replaced by the full equation
$$\dot{\varphi }_a(𝐱,t)=[^2+ϵ(t)+\overline{\beta }|\varphi (𝐱,t)|^2]\varphi _a(𝐱,t)+\overline{\eta }_a(𝐱,t),$$
(6.61)
where $`\overline{\beta }`$ is the rescaled coupling.
In order to retain some analytic understanding of the way that the density is such an ideal quantity to make predictions for, we adopt the approximation of preserving Gaussian fluctuations by linearising the self-interaction as
$$\dot{\varphi }_a(𝐱,t)=[^2+ϵ_{eff}(t)]\varphi _a(𝐱,t)+\overline{\eta }_a(𝐱,t),$$
(6.62)
where $`ϵ_{eff}`$ contains a (self-consistent) term $`O(\overline{\beta }|\varphi |^2)`$. Additive renormalisation is necessary, so that $`ϵ_{eff}ϵ`$, as given earlier, for $`tt_0`$.
Self-consistent linearisation is the standard approximation in non-equilibrium QFT, but is not strictly necessary here, since numerical simulations that identify line zeroes of the field can be made that use the full self-interaction. However, to date none address the questions we are posing here exactly, and until then there is virtue in analytic approximations provided they are not taken too seriously.
The solution for $`G(r,t)`$ is a straightforward generalisation of 4.32,
$$G(r,t)=_0^{\mathrm{}}𝑑\tau \overline{T}(t\tau /2)\left(\frac{1}{4\pi \tau }\right)^{3/2}e^{r^2/4\tau }e^{_0^\tau 𝑑sϵ_{eff}(ts/2)}.$$
(6.63)
Putting in the momentum cutoff $`k^1>l=\overline{l}\xi _0=O(\xi _0)`$ of Eq.5.52 by hand corresponds to damping the singularity in $`G(r,t)`$ at $`\tau =0`$ as
$$G_l(r,t)=_0^{\mathrm{}}\frac{d\tau \overline{T}(t\tau /2)}{[4\pi (\tau +\overline{l}^2)]^{3/2}}e^{r^2/4\tau }e^{_0^\tau 𝑑sϵ_{eff}(ts/2)},$$
(6.64)
making $`G_l(0,t)`$ finite. We stress that, for $`tt_0`$, the correlation length $`\xi `$ remains $`O(\overline{\xi })`$, independent of $`l`$.
Assuming a single zero of $`ϵ_{eff}(t)`$ at $`t=t_0`$, at $`r=0`$ the exponential in the integrand peaks at $`\tau =\overline{\tau }=2(tt_0)`$. Expanding about $`\overline{\tau }`$ to quadratic order gives
$$G_l(0,t)\overline{T}_ce^{2_{t_0}^t𝑑u|ϵ_{eff}(u)|}_0^{\mathrm{}}\frac{d\tau e^{(\tau 2(tt_0))^2|ϵ^{}(t_0)|/4}}{[4\pi (\tau +\overline{l}^2)]^{3/2}}.$$
(6.65)
For times $`t>ϵ_0\tau _Q`$ we see that, as the unfreezing occurs, long wavelength modes with $`k^2<t/\tau _Qϵ_0`$ grow exponentially.
The effect of the back-reaction is to stop the growth of $`G_l(0,t)G_l(0,t_0)=|\varphi |^2_t|\varphi |^2_0`$ at its symmetry-broken value $`\overline{\beta }^1`$ in our dimensionless units. A necessary condition for this is $`lim_u\mathrm{}ϵ_{eff}(u)=0`$. That is, we must choose
$$ϵ_{eff}(t)=ϵ(t)+\overline{\beta }(G_l(0,t)G_l(0,t_0)),$$
(6.66)
thereby preserving Goldstone’s theorem.
Beyond that, what is remarkable in this approximation is that the density of line zeroes uses no property of the self-mass contribution to $`ϵ_{eff}(t)`$, self-consistent or otherwise. With
$$G^{\prime \prime }(0,t)=\frac{1}{2}_0^{\mathrm{}}\frac{d\tau }{\tau }\overline{T}(t\tau /2)\left(\frac{1}{4\pi \tau }\right)^{3/2}e^{_0^\tau 𝑑sϵ_{eff}(ts/2)}.$$
(6.67)
all prefactors in $`n_{zero}`$ cancel<sup>6</sup><sup>6</sup>6Our ignoring prefactors in was fortuitous, leaving our conclusions obtained there unaffected., to give
$$n_{zero}(t)=\frac{1}{4\pi }\frac{_0^{\mathrm{}}\frac{d\tau }{(\tau +\overline{l}^2)^{5/2}}e^{(\tau 2(tt_0))^2/4\overline{t}^2}}{_0^{\mathrm{}}\frac{d\tau }{(\tau +\overline{l}^2)^{3/2}}e^{(\tau 2(tt_0))^2/4\overline{t}^2}}$$
(6.68)
on using the definition $`\tau _Q=\overline{t}^2`$ in natural units.
At $`t=t_0`$ both numerator and denominator are dominated by the short wavelength fluctuations at small $`\tau `$. Even though the field is correlated over a distance $`\overline{\xi }l`$ the density of line zeroes $`n_{zero}=O(l^2)`$ depends entirely on the scale at which we look. In no way would we wish to identify these line zeroes with prototype vortices. However, as time passes the peak of the exponential grows and $`n_{zero}`$ becomes increasingly insensitive to $`l`$. How much time we have depends on the magnitude of $`\overline{\beta }`$, since once $`G(0,t)`$ has reached this value it stops growing. Since $`G(0,t)=O(\mathrm{exp}(((tt_0)/\overline{t})^2)`$ at early times the backreaction is implemented extremely rapidly. We can estimate the time $`t^{}`$ at which this happens by substituting $`ϵ(u)`$ for $`ϵ_{eff}(u)`$ in the expression for $`G_l(0,t)`$ above.
For $`t>t^{}`$ the equation for $`n_{zero}(t)`$ is not so simple since the estimate above, based on a single isolated zero of $`ϵ_{eff}(t)`$, breaks down because of the approximate vanishing of $`ϵ_{eff}(t)`$ for $`t>t^{}`$. A more careful analysis shows that $`G_l(0,t)`$ can be written as
$$G_l(0,t)_0^{\mathrm{}}\frac{d\tau \overline{T}(t\tau /2)}{[4\pi (\tau +\overline{l}^2)]^{3/2}}\overline{G}(\tau ,t),$$
(6.69)
where $`\overline{G}(\tau ,t)`$ has the same peak as before at $`\tau =2(tt_0)`$, in the vicinity of which
$$\overline{G}(\tau ,t)=e^{2_{t_0}^t𝑑u|ϵ_{eff}(u)|}e^{(\tau 2(tt_0))^2/4\overline{t}^2},$$
(6.70)
but $`\overline{G}(\tau ,t)1`$ for $`\tau <2(tt^{})`$. Thus, for $`\tau _Q\tau _0`$, $`G_l(0,t)`$ can be approximately separated as
$$G_l(0,t)G_l^{UV}(t)+G^{IR}(t),$$
(6.71)
where
$$G_l^{UV}(t)=\overline{T}(t)_0^{\mathrm{}}𝑑\tau /[4\pi (\tau +\overline{l}^2)]^{3/2}$$
(6.72)
describes the scale-dependent short wavelength thermal noise, proportional to temperature, and
$$G^{IR}(t)=\frac{\overline{T}_c}{(8\pi (tt_0))^{3/2}}_{\mathrm{}}^{\mathrm{}}𝑑\tau \overline{G}(\tau ,t)$$
(6.73)
describes the scale-independent, temperature independent, long wavelength fluctuations. A similar decomposition $`G_l(0,t)G_l^{UV}(t)+G^{IR}(t)`$ can be performed as
$$G_l^{UV}(t)=2\pi \overline{T}(t)_0^{\mathrm{}}d\tau /[4\pi (\tau +\overline{l}^2)]^{5/2}$$
(6.74)
and
$$G^{IR}(t)=\frac{4\pi \overline{T}_c}{(8\pi (tt_0))^{5/2}}_{\mathrm{}}^{\mathrm{}}d\tau \overline{G}(\tau ,t).$$
(6.75)
In particular, $`G^{IR}(t)/G^{IR}(t)=O(t^1)`$.
Firstly, suppose that, for $`tt^{}`$, $`G^{IR}(t)G_l^{UV}(t)`$ and $`G^{IR}(t)G_l^{UV}(t)`$, as would be the case for a temperature quench $`\overline{T}(t)0`$. Then, with little thermal noise, we have widely separated line zeroes, with density $`n_{zero}(t)G^{IR}(t)/2\pi G^{IR}(t)`$. With $`n_{zero}/l`$ small in comparison to $`n_{zero}/l`$ at $`l=\xi _0`$ we identify such essentially non-fractal line-zeroes with prototype vortices, and $`n_{zero}`$ with $`n_{def}`$. Of course, we require non-Gaussianity to create true classical energy profiles. Nonetheless, the Halperin-Mazenko result may be well approximated for a while even when the fluctuations are no longer Gaussian. This is supported by the observation that, once the line zeroes have straightened on small scales at $`t>t^{}`$, the Gaussian field energy, largely in field gradients, is
$$\overline{F}_Vd^3x\frac{1}{2}(\varphi _a)^2=VG^{\prime \prime }(0,t),$$
(6.76)
where $`V`$ is the spatial volume. This matches the energy
$$\overline{E}Vn_{def}(t)(2\pi G(0,t))=VG^{\prime \prime }(0,t)$$
(6.77)
possessed by a network of classical global strings with density $`n_{zero}`$, in the same approximation of cutting off their logarithmic tails.
For times $`t>t^{}`$
$$n_{zero}(t)\frac{\overline{t}}{8\pi (tt_0)}\frac{1}{\xi _0^2}\sqrt{\frac{\tau _0}{\tau _Q}},$$
(6.78)
the solution to Vinen’s equation
$$\frac{n_{zero}}{t}=\chi _2\frac{\mathrm{}}{m}n_{zero}^2,$$
(6.79)
where $`\chi _2=4\pi \mathrm{}\mathrm{\Gamma }=4\pi \mathrm{}/\alpha _0\tau _0=4\pi `$ if, as earlier, we motivated $`\tau _0`$ from the Gross-Pitaevskii equation, in which $`\alpha _0\tau _0=\mathrm{}`$. More realistically, we find $`\chi _2>4\pi `$ both for $`{}_{}{}^{4}He`$ and $`{}_{}{}^{3}He`$. Taking $`\tau _08.0\times 10^{12}s`$ and $`\xi _05.6\AA `$) in the mean-field approximation for $`{}_{}{}^{4}He`$ gives $`\chi _25\times 4\pi `$. For $`{}_{}{}^{3}He`$ (with $`\xi _077nm,\tau _01ns`$) we find $`\chi _210\times 4\pi `$.
This decay law is assumed in the analysis of the Lancaster experiments, in which the density of vortices is inferred from the intensity of the signal of scattered second sound. RG improvement leaves $`{}_{}{}^{3}He`$ unchanged, but for $`{}_{}{}^{4}He`$ it redefines $`\chi _2`$ to $`\chi _2(1T/T_c)^{1/3}>\chi _2`$.
The first problem for the Lancaster experiments is that this makes an already large $`\chi _2`$ even larger. In the attenuation of second sound the signal to noise ratio is approximately $`O(1/\chi _2t)`$. The empirical value of $`\chi _2`$ used in the Lancaster experiments is not taken from quenches, but turbulent flow experiments. It is suggested that $`\chi _20.005`$, a good three orders of magnitude smaller than our prediction above. Although the TDLG theory is not very reliable for $`{}_{}{}^{4}He`$, if our estimate is sensible it does imply that vortices produced in a temperature quench decay much faster than those produced in turbulence. $`{}_{}{}^{3}He`$ experiments provide no check.
This is only one of our worries. We shall argue that, for early time at least, thermal fluctuations are large in the Lancaster experiments. However, for $`{}_{}{}^{3}He`$, with negligable UV contributions, we estimate the primordial density of proto-vortices as
$$n_{zero}(t^{})\frac{\overline{t}}{8\pi (t^{}t_0)}\frac{1}{\xi _0^2}\sqrt{\frac{\tau _0}{\tau _Q}},$$
(6.80)
in accord with the original prediction of Zurek. Because of the rapid growth of $`G(0,t)`$, $`(t^{}t_0)/\overline{t}=p>1=O(1)`$. With $`p`$ behaving as $`(\mathrm{ln}(1/\overline{\beta }))^{1/2}`$ there is very little variation. For $`{}_{}{}^{3}He`$ quenches $`p5`$ (and for $`{}_{}{}^{4}He`$ quenches $`p3`$). We note that the factor<sup>7</sup><sup>7</sup>7An errant factor of 3 appeared in the result of of $`f^2=8\pi p`$ gives a value of $`f=O(10)`$, in agreement with the empirical results of and the numerical results of <sup>8</sup><sup>8</sup>8The temperature quench of the latter is somewhat different from that considered here, but should still give the same results in this case.
### 6.2 Path integrals
Finally, we note that $`G(r,t)`$, and hence its derivatives, can be expressed as path integrals for the diffusion of a particle in a time-dependent potential $`ϵ_{eff}`$. Thus, for example,
$$G^{\prime \prime }(0,t)=\frac{T}{2}_0^{\mathrm{}}\frac{d\tau }{\tau }𝒟𝐱\mathrm{exp}\left[_0^\tau 𝑑s\frac{1}{4}(\frac{d𝐱}{d\tau })^2+ϵ_{eff}(ts/2)\right].$$
(6.81)
In Eq.6.81 the summation is over closed paths for the particle (mass 2, in our units) traversed in time $`\tau `$. We note that, if $`ϵ=ϵ_0`$ is fixed, then
$$G^{\prime \prime }(0,t)=\frac{T}{2}_0^{\mathrm{}}\frac{d\tau }{\tau }𝒟𝐱e^{L(\tau )ϵ_0}$$
(6.82)
where $`L(\tau )=_0^\tau 𝑑\tau ^{}\sqrt{\dot{𝐱}^2(\tau ^{})}`$ is the length of the loop. Eq.6.82 is just the free energy of a ’free gas’ of fluctuating relativistic Brownian loops of varying lengths and shapes with tension $`ϵ_0`$. In conventional statistical field theory the vanishing of $`ϵ_0`$ in Eq.6.82 signals a phase transition due to the proliferation of loops. The situation here is somewhat different in that the loops have a tension that varies with both the external time $`t`$ and $`\tau `$. The onset of the transition is signaled by this tension vanishing at some point on the loops.
## 7 For Slow Pressure Quenches or a Large Ginzburg Regime Thermal Fluctuations Cannot be Ignored
### 7.1 $`{}_{}{}^{4}He`$ experiments
The situation for the Lancaster $`{}_{}{}^{4}He`$ experiments is complex, since they are pressure quenches for which the temperature $`T`$ is almost constant at $`TT_c`$. Unlike temperature quenches, thermal fluctuations here remain at full strength<sup>9</sup><sup>9</sup>9Even for $`{}_{}{}^{3}He`$, $`T/T_c`$ never gets very small, and henceforth we take $`T=T_c`$ in $`G_l(0,t)`$ above. The necessary time-independence of $`G^{IR}(t)`$ for $`t>t^{}`$ is achieved by taking $`ϵ_{eff}(u)=O(u^1)`$. In consequence, as $`t`$ increases beyond $`t^{}`$ the relative magnitude of the UV and IR contributions to $`G_l(0,t)`$ remains approximately constant. Further, since for $`t=t^{}`$,
$$e^{2_{t_0}^t𝑑u|ϵ_{eff}(u)|}e^{(\mathrm{\Delta }t)^2/\overline{t}^2}1,$$
(7.83)
this ratio is the ratio at $`t=t^{}`$.
Nonetheless, as long as the UV fluctuations are insignificant at $`t=t^{}`$ the density of line zeroes will remain largely independent of scale. This follows if $`G^{IR}(t^{})G_l^{UV}(t^{})`$, since $`G_l(0,t)`$ becomes scale-independent later than $`G_l(0,t)`$. In we showed that this is true provided
$$(\tau _Q/\tau _0)(1T_G/T_c)<C\pi ^4,$$
(7.84)
where $`C=O(1)`$ and $`T_G`$ is the Ginzburg temperature. With $`\tau _Q/\tau _0=O(10^3)`$ and $`(1T_G/T_c)=O(10^{12})`$ this inequality is well satisfied for a linearised TDLG theory for $`{}_{}{}^{3}He`$ derived<sup>10</sup><sup>10</sup>10Ignoring the position-dependent temperature of from the full TDGL theory, but there is no way that it can be satisfied for $`{}_{}{}^{4}He`$, when subjected to a slow mechanical quench, as in the Lancaster experiment, for which $`\tau _Q/\tau _0=O(10^{10})`$, since the Ginzburg regime is so large that $`(1T_G/T_c)=O(1)`$. Effectively, the $`{}_{}{}^{4}He`$ quench is nineteen orders of magnitude slower than its $`{}_{}{}^{3}He`$ counterpart. It is satisfying to find the Ginzburg regime reappear as an essential ingredient in undermining universality. In particular, as in the original proposal of Kibble, large thermal noise inhibits vortex production.
When the inequality is badly violated, as with $`{}_{}{}^{4}He`$ for slow pressure quenches, then the density of zeroes $`n_{zero}=O(l^2)`$ after $`t^{}`$ depends exactly on the scale $`l`$ at which we look and they are not candidates for vortices. Since the whole of the quench takes place within the Ginzburg regime this is not implausible. However, it is possible that, even though the thermal noise never switches off, there is no more than a postponement of vortex production, since our approximations must break down at some stage. The best outcome is to assume that the effect of the thermal fluctuations on fractal behaviour is diminished, only leading to a delay in the time at which vortices finally appear. Even if we suppose that $`n_{zero}`$ above is a starting point for calculating the density at later times, albeit with a different $`t_0`$, thereby preserving Vinen’s law, we then have the earlier problem of the large $`\chi _2=O(f^2)`$. In the absence of any mechanism to reduce its value drastically, this would make it impossible to see vortices. As a separate observation, we note that the large value of $`f^2`$ in the prefactor of $`n_{zero}`$ is, in itself, almost enough to make it impossible to see vortices in $`{}_{}{}^{4}He`$ experiments, should they be present. This will be pursued elsewhere.
In summary, this work suggests that for slow pressure quenches in $`{}_{}{}^{4}He`$, we see no well-defined vortices at early times because of thermal fluctuations, and it is plausible that, if we do see them at later times, there are less than we would have expected because of their rapid decay and their initial low density. The situation is different for $`{}_{}{}^{3}He`$. That the density of such vortices as appear may agree with the Zurek prediction is essentially a consequence of dimensional analysis, given that the main effects of the self-interaction have a tendency to cancel in the counting of vortices, without introducing new scales. However, a numerical simulation that goes beyond the Gaussian approximation that is specifically tailored to the Lancaster parameters is crucial if we are to understand the late-time behaviour and see if this suggestion can be sustained. We shall pursue this elsewhere.
### 7.2 TDLG: annular physics
We saw for vortex formation in $`{}_{}{}^{4}He`$ that the dependence of the density on scale makes the interpretation of observations problematic. This is not the same if the $`{}_{}{}^{4}He`$ is confined to an annulus, since the annulus itself provides the coarse-graining. That the incoherent $`\xi _s`$ depends on its radius $`l`$ is immaterial. The end result is that
$$\mathrm{\Delta }v=\frac{\mathrm{}}{m}\sqrt{\frac{1}{C\xi _s(t^{};l)}}.$$
(7.85)
The time $`t^{}`$ for $`{}_{}{}^{4}He`$ at which we evaluate $`v`$, when $`|\varphi |^2=\alpha _0/\beta `$, depends weakly on $`l`$, varying from about $`3\overline{t}`$ to $`4\overline{t}`$ as $`l`$ varies from $`l=\xi _0\overline{\xi }`$ to $`l=10\overline{\xi }`$. For $`l4\overline{\xi }`$ the scale at which the coarse-grained field begins to occupy the ground states becomes largely irrelevant.
On expanding $`f^2(r,t;l)/(1f^2(r,t;l))`$ in powers of $`r`$ and keeping only the forward peak,
$$\frac{1}{\xi _s(t;l)}\frac{4G_2}{9G_1}\left(\frac{3G_3}{20G_2}\frac{G_2}{12G_1}\right)^{1/2},$$
(7.86)
where $`G_n(r,t)`$ is the $`(2n2)`$th derivative of $`G`$ with respect to $`r`$ at $`r=0`$ (and hence proportional to the $`2n`$th moment of $`G(k,t)`$).
If, as suggested by Zurek, we take $`l=O(\overline{\xi })`$ we recover Eq.3.17 qualitatively, although a wider bore would give a correspondingly smaller flow.
We assume that
$$2l\xi _{eff}(t;l)=\left(\frac{3G_3}{20G_2}\frac{G_2}{12G_1}\right)^{1/2},$$
(7.87)
since otherwise the correlations in the bulk fluid from which we want to extract annular behaviour are of longer range than the annulus thickness. The effect of this is to reduce the flow velocity for narrower annuli. The effect is largest for small radii $`l\overline{\xi }`$, for which the approximation of trying to read the behaviour of annular flow from bulk behaviour is most suspect.
Once $`l`$ is very large, so that the power in the fluctuations is distributed strongly across all wavelengths we recover our earlier result, that $`\xi _s(t^{};l)=O(l)`$. However, the change is sufficiently slow that annuli, significantly wider than $`\overline{\xi }`$, for which experiments are more accessible, will give almost the same flow as narrower annuli. This would seem to extend the original Zurek prediction of Eq.3.17 to thicker annuli, despite our expectations for incoherent flow. However, we stress again that caution is necessary, since in the approximation to characterise an annulus by a coarse-grained ring without boundaries we have ignored effects in the direction perpendicular to the annulus. In particular, the circular cross-section of the tube has not been taken into account. One consequence of this is that infinite (non-self-intersecting) vortices in the bulk fluid have no counterpart in an annulus.
Since $`\mathrm{\Delta }v`$ only depends on $`\xi _s^{1/2}`$ it is not sensitive to choice of $`l>2\overline{\xi }`$ at the relevant $`t`$. Given all these approximations our final estimate is (in the cm/sec units of Zurek)
$$\mathrm{\Delta }v0.2(\tau _Q[\mu s])^{\nu /4}/\sqrt{C[cm]}$$
(7.88)
for radii of $`2\overline{\xi }4\overline{\xi }`$, $`\tau _Q`$ of the order of milliseconds and $`C`$ of the order of centimetres. $`\nu =1/2`$ is the mean-field critical exponent above. In principle $`\nu `$ should be renormalised to $`\nu =2/3`$, but the difference to $`\mathrm{\Delta }v`$ is sufficiently small that we shall not bother. Given the uncertainties in its derivation the result Eq.7.88 is indistinguishable from Zurek’s (with prefactor $`0.4`$), but for the possibility of using somewhat larger annuli. The agreement is, ultimately, one of dimensional analysis, but the coefficient of unity could not have been anticipated easily, given the small prefactors of Eq.6.20. How experiments can be performed, even with wider annuli, is another matter.
## 8 QFT: The Appearance of Structure
It is because the formation of defects is an early-time occurrence that it is, in part, amenable to analytic solution. Again we revert to the mode decomposition of Eq.4.42. The field becomes ordered, as before, because of the exponential growth of long-wavelength modes, which stop growing once the field has sampled the groundstates. What matters is the relative weight of these modes (the ’Bragg’ peak) to the fluctuating short wavelength modes, since the contribution of these latter is very sensitive to the cutoff $`l`$ at which we look for defects. Only if their contribution to Eq.2.14 is small when field growth stops can a network of vortices be well-defined at early times, let alone have the predicted density. Since the peak is non-perturbatively large this requires small coupling, which we assume. However, there is a problem in that, in the absence of explicit damping of the type seen in FRW universes, rescattering of modes can rapidly undo the early-time appearance of structure. Thus,while we take small coupling, we do not take very small couplings of the magnitude (e.g. $`10^{12}`$) associated with inflationary models.
### 8.1 Mode Growth v Fluctuations: The free roll (continued)
We begin by extending the analysis of Section 4.2 to later times, still in the approximation of a free roll. This needs care for slow quenches since the backreaction serves to hold the field in the vicinity of the intermediate groundstates $`|\varphi ^2|=\varphi _0^2(t)`$ where, now
$$\varphi _0^2(t)=\frac{m^2(t)}{\lambda }=\frac{M^2}{\lambda }\frac{\mathrm{\Delta }t}{\tau _Q},$$
(8.89)
where $`\mathrm{\Delta }t=tt_0=tϵ_0\tau _Q`$ as before. Nonetheless, the free roll provides a basis for the more realistic picture.
Prior to the completion of the quench at $`\mathrm{\Delta }t=\tau _Q`$, the mode equation (4.42), now of the form
$$\left[\frac{d^2}{dt^2}+𝐤^2\frac{M^2\mathrm{\Delta }t}{\tau _Q}\right]\chi _k^\pm (t)=0,$$
(8.90)
is exactly solvable, as we saw earlier.
For this section it is convenient to redefine the origin of time at $`t=t_0`$, whereby we can drop the prefix $`\mathrm{\Delta }`$.
We are primarily interested in the exponentially growing modes that appear when
$$\mathrm{\Omega }_k^2(t)=𝐤^2+\frac{M^2t}{\tau _Q}>0.$$
(8.91)
For fixed $`k`$ this occurs when $`t>t_k^{}=\tau _Qk^2/M^2`$.
The WKB solution is adequate for our purposes. The coarse-grained $`G_l(r;t)`$ can be written as $`G_l(r;t)=G^{exp}(r;t)+G_l^{osc}(r;t)`$ where
$$G^{exp}(r;t)\frac{T_0}{M^2}_{|𝐤|<k_t}d/^3k\frac{MS_k(t)}{|k_t^2k^2|^{1/2}}e^{i𝐤.𝐱}|\alpha _k^+I_{1/3}(S_k(t))+\alpha _k^{}I_{1/3}(S_k(t))|^2$$
(8.92)
has exponentially growing long wavelength modes and
$$G_l^{osc}(r;t)\frac{T_0}{M^2}_{\mathrm{\Lambda }>|𝐤|>k_t}d/^3k\frac{MS_k(t)}{|k_t^2k^2|^{1/2}}e^{i𝐤.𝐱}|\alpha _k^+J_{1/3}(S_k(t))\alpha _k^{}J_{1/3}(S_k(t))|^2$$
(8.93)
has short wavelength oscillatory modes. In both cases
$$S_k(t)=_{t_k^{}}^t𝑑t^{}|\mathrm{\Omega }_k(t^{})|=\frac{2}{3}\frac{M}{\sqrt{\tau _Q}}|tt_k^{}|^{3/2}.$$
(8.94)
As in the case of condensed matter previously, we have coarse-grained the field by introducing a simple cut-off at $`k=\mathrm{\Lambda }=O(M)`$, or $`l=\mathrm{\Lambda }^1`$. For fixed $`t`$ the dividing monentum is $`k_t^2=M^2t/\tau _Q`$. Provided we are far from the transition we have incorporated the initial data into the $`\alpha ^\pm `$ in (8.92) and (8.93). The normalisation factor $`T_0/M^2`$ has been made visible. The remaining $`\alpha ^\pm `$ have no $`\lambda `$ dependence. Since $`G^{osc}(0;t)=O(T_0M)`$ for $`t=O(\tau _Q)`$ the $`\alpha ^\pm `$ have no $`\tau _Q`$ dependence.
For large $`t`$ the integrand in (8.92) will be peaked at some $`k_0(t)0`$ as $`t\mathrm{}`$, once the angular integrals have been performed. Assuming that $`k_0(t)k_t`$ the upper bound in the integral can be dropped and $`|k_t^2k^2|`$ approximated by $`|k_t^2|`$, knowing that there is no singularity at $`k=k_t`$. With nothing to stop $`|\alpha _k^++\alpha _k^{}|^2`$ behaving like a nonzero constant in the vicinity of $`k=0`$, it can be treated as slowly varying and the integral approximated as
$`G^{exp}(r;t)`$ $``$ $`{\displaystyle \frac{T}{M^2}}\left({\displaystyle \frac{\tau _Q}{t}}\right)^{1/2}e^{4Mt^{3/2}/3\sqrt{\tau _Q}}{\displaystyle _{|𝐤|<M}}d/^3ke^{i𝐤.𝐱}e^{2\sqrt{t\tau _Q}k^2/M}`$ (8.95)
$``$ $`{\displaystyle \frac{T}{M|m(t)|}}\left({\displaystyle \frac{M}{\sqrt{t\tau _Q}}}\right)^{3/2}e^{4Mt^{3/2}/3\sqrt{\tau _Q}}e^{r^2/\xi ^2(t)}`$
on performing the $`k^2`$ expansion of the exponent, where
$$\xi ^2(t)=\frac{2\sqrt{t\tau _Q}}{M}.$$
(8.96)
Although, like Eq.8.95, the expression Eq.8.96 is not supposed to be valid for small $`t`$, it does embody the Kibble freezeout condition Eq.2.10 in satisfying $`\dot{\xi }(\overline{t})=O(1)`$.
### 8.2 Comparison with Kibble’s results: First guess
The calculations we did on CM systems showed that, although the freeze-in time $`\overline{t}`$ was not relevant for the appearance of stable vortices, their subsequent density is determined by the scales at this time in the absence of strong fluctuations. We shall see that the same could be true here.
The calculations above were for a free roll. Let us suppose, provisionally, that the backreaction exerts its influence over such a short time that, in effect, it is if it were an instantaneous brake to domain growth. The provisional freeze-in time $`t^{}`$ is then, effectively, the time it takes to reach the transient groundstate $`\varphi _0^2(t)=m^2(t)/\lambda `$. That is, $`G(0;t^{})=O(\varphi _0^2(t^{}))`$, giving
$$(\sqrt{t^{}\tau _Q}M)^{3/2}e^{4M(t^{})^{3/2}/3\sqrt{\tau _Q}}=O\left(\lambda ^{1/2}(\frac{\tau _Q}{t^{}})^{3/2}\right).$$
(8.97)
Neglecting $`lnln`$ terms gives
$$Mt^{}\frac{1}{2}(M\tau _Q)^{1/3}\left(\mathrm{ln}(1/\lambda )+\mathrm{ln}(M\tau _Q)\right)^{2/3}.$$
(8.98)
The second term in the bracket of (8.98) can be ignored to a good approximation provided $`M\tau _Q<(1/\lambda )`$, preferably by a large margin, and we assume that this is so. In terms of the causal time $`\overline{t}`$ the relation is just
$$Mt^{}M\overline{t}(\mathrm{ln}(1/\lambda ))^{2/3}.$$
(8.99)
That is, the freeze in time $`t^{}`$ is (qualitatively) larger than the causal time $`\overline{t}`$. As far as the separation of scales is concerned, we have the same effect qualitatively if we had taken $`t^{}`$ as the time for the field to reach the final ground state as $`|\varphi ^2|=M^2/\lambda `$, rather than the provisional ground states $`\varphi _0^2(t)`$. Nonetheless, we shall remain with (8.97).
Whatever, the necessary condition that $`t^{}<\tau _Q`$ (since the previous calculations are for $`t<\tau _Q`$) reduces to
$$M\tau _Q>\mathrm{ln}(1/\lambda ).$$
(8.100)
where, and henceforth, we neglect coefficients of $`O(1)`$. That is, the quench time should be longer than the freeze-in time for the instantaneous quench, the time it takes the field to sample the ground-state in a free-roll<sup>11</sup><sup>11</sup>11In general, we recover the results of the instantaneous quench if we set $`M\tau _Q=O(\mathrm{ln}(1/\lambda ))`$, but shorter than the equilibriation time (or powers of it that preserve the separation of scales).
At this qualitative level the correlation length at the spinodal time is
$$M^2\xi ^2(t^{})(M\tau _Q)^{2/3}(\mathrm{ln}(1/\lambda ))^{1/3}.$$
(8.101)
The effect of the other modes is larger than for the instantaneous quench, giving, at $`t=t^{}`$
$$n_{zero}=\frac{M^2}{\pi (M\tau _Q)^{2/3}}(\mathrm{ln}(1/\lambda ))^{1/3}[1+E].$$
(8.102)
The error term $`E=O(\lambda ^{1/2}(M\tau _Q)^{4/3}(\mathrm{ln}(1/\lambda ))^{1/3})`$ is due to oscillatory modes, sensitive to the cut-off. In mimicry of Eq.2.14 it is helpful to rewrite Eq.8.102 as
$$n_{zero}=\left[\frac{1}{\pi \xi _0^2}\left(\frac{\tau _0}{\tau _Q}\right)^{2/3}\right](\mathrm{ln}(1/\lambda ))^{1/3}[1+E].$$
(8.103)
in terms of the scales $`\tau _0=\xi _0=M^1`$. The first term in Eq.8.103 is the Kibble estimate of Eq.2.14, the second is the multiplying factor, rather like that in Eq.6.80, that yet again shows that estimate can be correct, but for completely different reasons. As for condensed matter, the dependence on the interaction strength is only through a power of the logarithm of $`(1/\lambda ))^{1/3}`$.
The third term shows when it can be correct, since $`E`$ is also a measure of the sensitivity of $`n_{zero}`$ to the scale at which it is measured. The condition $`E^21`$, necessary for a vortex network to be defined, is then guaranteed if
$$(\tau _Q/\tau _0)^2(1T_G/T_c)<C,$$
(8.104)
where $`C=O(1)`$, on using the relation $`(1T_G/T_c)=O(\lambda )`$. This is the QFT counterpart to Eq.7.84 and concurs again, in principle, with Kibble’s earlier argument that large thermal fluctuations inhibit the appearance of vortices.
The easiest way to enforce $`E1`$ and $`M\tau _Q>\mathrm{ln}(1/\lambda )`$ is to take $`M\tau _Q=\mathrm{ln}(1/\lambda )^\alpha `$, for $`\alpha >1`$. The effect in Eq.8.103 is merely to renormalise the critical index. Of course, the Kibble prediction Eq.2.14 was only an estimate. Although it is good qualitatively it is misleading when considering definition of the network since a simple calculations shows that, at time $`\overline{t}`$, the string density is still totally sensitive to the definition of coarse-graining.
Finally, suppose that this approach is relevant to the local strings of a strong Type-II $`U(1)`$ theory for the early universe, in which the time-temperature relationship $`tT^2=\mathrm{\Gamma }M_{pl}`$ is valid, where we take $`\mathrm{\Gamma }=O(10^1)`$ in the GUT era. If $`G`$ is Newton’s constant and $`\mu `$ the classical string tension then, following , $`M\tau _Q10^1\lambda ^{1/2}(G\mu )^{1/2}`$. The dimensionless quantity $`G\mu 10^610^7`$ is the small parameter of cosmic string theory. A value $`\lambda 10^2`$ gives $`M\tau _Q(Mt^{})^a,a2`$, once factors of $`\pi `$, etc.are taken into account, rather than $`M\tau _Q1/\lambda `$, and the density of Eq.8.103 may be relevant.
### 8.3 Backreaction in QFT
To improve upon the free-roll result more honestly, but retain the Gaussian approximation for the field correlation functions, the best we can do is adopt a mean-field approximation along the lines of , as we did for the CM systems earlier. As there, it does have the correct behaviour of stopping domain growth as the field spreads to the potential minima. As before, only the large-$`N`$ expansion preserves Goldstone’s theorem.
$`G(r;t)`$ still has the mode decomposition of (4.3), but the modes $`\chi _k^\pm `$ now satisfy the equation
$$\left[\frac{d^2}{dt^2}+𝐤^2+m^2(t)+\lambda \mathrm{\Phi }^2(\mathrm{𝟎})_t\right]\chi _k^\pm (t)=0,$$
(8.105)
where we have taken $`N=2`$. Because $`\lambda \varphi ^4`$ theory is not asymptotically free, particularly in the Hartree approximation, the renormalised $`\lambda `$ coupling shows a Landau ghost. This means that the theory can only be taken as a low energy effective theory.
The end result is, on making a single subtraction at $`t=0`$, is
$$\left[\frac{d^2}{dt^2}+𝐤^2+m^2(t)+\lambda d/^3pC(p)[\chi _p^+(t)\chi _p^{}(t)1]\right]\chi _k^\pm (t)=0.$$
(8.106)
which we write as
$$\left[\frac{d^2}{dt^2}+𝐤^2\mu ^2(t)\right]\chi _k(t)=0.$$
(8.107)
On keeping just the unstable modes in $`\mathrm{\Phi }^2(\mathrm{𝟎})_t`$ then, as it grows, its contribution to (8.106) weakens the instabilities, so that only longer wavelengths become unstable. At $`t^{}`$ the instabilities shut off, by definition, and oscillatory behaviour ensues. Since the mode with wavenumber $`k>0`$ stops growing at time $`t_k^+<t^{}`$, where $`\mu ^2(t_k^+)=𝐤^2`$, the free-roll density at $`t^{}`$ must be an overestimate.
An approximation that improves upon the WKB approximation is
$$\chi _k(t)\left(\frac{\pi M}{2\mathrm{\Omega }_k(\eta )}\right)^{1/2}\mathrm{exp}\left(_0^t𝑑t\mathrm{\Omega }(t)\right)$$
(8.108)
when $`\eta =M(t_k^+t)>0`$ is large, and $`\mathrm{\Omega }_k^2(t)=\mu ^2(t)𝐤^2`$. On expanding the exponent in powers of $`k`$ and retaining only the quadratic terms we recover the WKB approximation when $`\mu (t)`$ is non-zero.
The result is that the effect of the back-reaction is to give a time-delay $`\mathrm{\Delta }t`$ to $`t^{}`$, corresponding to a decrease in the value $`k_0(t)`$ at which the power peaks of order
$$\frac{\mathrm{\Delta }t}{t^{}}=O\left(\frac{1}{ln(1/\lambda )}\right).$$
(8.109)
The backreaction has little effect for times $`t<t^{}`$. For $`t>t^{}`$ oscillatory modes take over the correlation function and we expect oscillations in $`G(k;t)`$.
In practice the backreaction rapidly forces $`\mu ^2(t)`$ towards zero if the coupling is not too small. For couplings that are not too weak the end result is a new power spectrum, obtained by superimposing oscillatory behaviour onto the old spectrum. As a gross oversimplification, the contribution from the earlier exponential modes alone can only be to contribute terms something like
$`G(r;t)`$ $``$ $`{\displaystyle \frac{T}{M|m(t^{})|}}e^{4M(t^{})^{3/2}/3\sqrt{\tau _Q}}{\displaystyle _{|𝐤|<M}}d/^3ke^{i𝐤.𝐱}e^{2\sqrt{t^{}\tau _Q}k^2/M}`$ (8.110)
$`\times `$ $`\left[\mathrm{cos}k(tt^{})+{\displaystyle \frac{\mathrm{\Omega }(k)W^{}(k)}{k}}\mathrm{sin}k(tt^{})\right]^2`$
to $`G`$, where $`\mathrm{\Omega }=M(t^{}t_k)^{1/2}/\tau _Q^{1/2}`$ and $`W^{}=1/4(t^{}t_k)`$. Fortunately, the details are almost irrelevant, since the density of line zeroes is independent of the normalisation, and only weakly dependent on the power spectrum.
The $`k=0`$ mode of Eq.8.110 encodes the simple solution $`\chi _{k=0}(t)=a+bt`$ when $`\mu ^2=0`$. As observed by Boyanovsky et al. this has built into it the basic causality discussed by Kibble. Specifically, for $`r,t\mathrm{}`$, but $`r/t`$ constant $`(2)`$,
$$G(r,t)\frac{C}{r}\mathrm{\Theta }(2t/r1).$$
(8.111)
It has to be said that this approximation should not be taken very seriously for large $`t`$, since we would expect rescattering to take place at times $`\mathrm{\Delta }t=O(1/\lambda )`$ in a way that the Gaussian approximation precludes.
Whatever, it follows directly that this causality, engendered by the Goldstone particles of the self-consistent theory, has little effect on the density of line-zeroes that we expect to mature into fully classical vortices. If Eq.8.104 is not satisfied, it is difficult to imagine how clean vortices, or proto-vortices, can appear later without some additional ingredient.
### 8.4 Really slow quenches
Finally, consider slowing the quench until the WKB approximation manifestly breaks down. This is almost certainly the case if the growing modes catch up with the moving minima within the Ginzburg regime by time
$$t_G=O(\lambda \tau _Q).$$
(8.112)
If we take (8.97) seriously it follows that domain growth has stopped by time $`t_G`$ provided $`M\tau _Q`$ is a power of $`\lambda ^1`$. Equation (8.97) only makes sense if $`Mt_G1`$ i.e $`M\tau _Q\lambda ^1`$. As an example, let us take $`M\tau _Q=\lambda ^{3/2}`$. In this case the dynamical correlation length $`\xi (t_G)`$ of (8.96) is $`O(1/M\lambda ^{1/2})`$, equal to the equilibrium correlation length $`\xi _{eq}(t_G)=|m^1(t_G)|`$, suggesting the string density $`n_G=O(M^2/\lambda )`$ that follows from (2.6), very much smaller than that of (8.103) above. However, a simple check shows that, in this case,
$$\frac{l}{n_l(t_G)}\frac{n_l(t_G)}{l}|_{l=M^1}=O(1).$$
(8.113)
Despite the low density of line zeroes at the Ginzsburg temperature, they do not provide a stable network of proto-vortices.
## 9 Conclusions
We examined the Kibble /Zurek predictions for the onset of phase transitions and the appearance of defects (in particular, vortices or global cosmic strings) as a signal of the symmetry breaking. Our results are in agreement with their prediction Eq.2.13 as to the magnitude of the correlation length at the time the transition truly begins, equally true for condensed matter and QFT.
However, this is not simply a measure of the separation of defects at the time of their appearance. The time $`\overline{t}`$ is too early for the field to have found the true groundstates of the theory. We believe that time, essentially the spinodal time, is the time at which proto-vortices can appear, which can later evolve into the standard classical vortices of the theory.
Even then, they may be frustrated by thermal field fluctuations. In TDLG condensed matter thermal noise is proportional to temperature. If temperature is fixed, but not otherwise, as in the pressure quenches of $`{}_{}{}^{4}He`$ this noise can inhibit the production of vortices, although there are other factors to be taken into account (such as their decay rate). On quenching from a high temperature in QFT there are always thermal fluctuations, and these can also disturb the appearance of vortices. The condition that thermal fluctuations are ignorable at the time that the field has achieved the true ground-states can be written
$$(\tau _Q/\tau _0)^\gamma (1T_G/T_c)<C,$$
(9.114)
where $`\gamma =1`$ for condensed matter and $`\gamma =2`$ for QFT. $`C=O(1)`$.
This restores the role of the Ginzburg temperature $`T_G`$ that the simple causal arguments overlooked. As in the earlier arguments, large fluctuations inhibit vortex production, but in their absence the equilibrium correlation length is not the relevant quantity. Quenches in $`{}_{}{}^{4}He`$ provide the major example for which Eq.9.114 is not satisfied. However, when Eq.9.114 is satisfied, the Kibble/Zurek prediction is recovered for the density of line zeroes, now seen as potential vortices. That no new scales (up to logarithms) are introduced by the interaction is a reflection of the fact that line-zero density only uses limited properties of the power spectrum of fluctuations (in the Gaussian approximation at least).
What happens at late time is unclear, although for TDLG numerical simulations can be performed (but have yet to address this problem exactly). On the other hand, not only is the case of a single self-interacting quantum scalar field in flat space-time a caricature of the early universe, but it is extremely difficult to go beyond the Gaussian approximation. To do better requires that we do differently. There are several possible approaches. One step is to take the FRW metric of the early universe seriously, whereby the dissipation due to the expansion of the universe can change the situation dramatically. Other approaches are more explicit in their attempts to trigger decoherence explicitly, as we mentioned earlier. Most simply, one treats the short wavelength parts of the field as an environment to be integrated over, to give a coarse-grained theory of long-wavelength modes acting classically in the presence of noise. However, such noise is more complicated than in TDLG theory, being multiplicative as well as additive, and coloured. This is an area under active consideration, but we shall stop here.
I would like to thank Alasdair Gill, Tom Kibble, Wojciech Zurek and Grisha Volovik for fruitful discussions. There is much work in this area and we have not included all references. We apologise to any authors who we may have missed inadvertently. This work is the result of a network supported by the European Science Foundation. |
warning/0001/hep-ph0001087.html | ar5iv | text | # Electroweak Phase Transition in the MSSM: 4-Dimensional Lattice Simulations
\[
## Abstract
Recent lattice results have shown that there is no Standard Model (SM) electroweak phase transition (EWPT) for Higgs boson masses above $``$ 72 GeV, which is below the present experimental limit. According to perturbation theory and 3-dimensional (3d) lattice simulations there could be an EWPT in the Minimal Supersymmetric Standard Model (MSSM) that is strong enough for baryogenesis up to $`m_h105`$ GeV. In this letter we present the results of our large scale 4-dimensional (4d) lattice simulations for the MSSM EWPT. We carried out infinite volume and continuum limits and found a transition whose strength agrees well with perturbation theory, allowing MSSM electroweak baryogenesis at least up to $`m_h=103\pm 4`$ GeV. We determined the properties of the bubble wall that are important for a successful baryogenesis.
PACS numbers: 11.10.Wx, 11.15.Ha, 12.60.Jv, 98.80.Cq
\]
The visible Universe is made up of matter. This statement is mainly based on observations of the cosmic diffuse $`\gamma `$-ray background, which would be larger than the present limits if boundaries between “worlds” and “anti-worlds” existed . The observed baryon asymmetry of the universe was eventually determined at the EWPT . On the one hand this phase transition was the last instance when baryon asymmetry could have been generated, around $`T100200`$ GeV. On the other hand at these temperatures any B+L asymmetry could have been washed out. The possibility of baryogenesis at the EWPT is particularly attractive, since the underlying physics can be—and has already largely been—tested in collider experiments.
The first detailed description of the EWPT in the SM was based on perturbative techniques , which resulted in $`𝒪(100\%)`$ corrections between different orders of the perturbative expansion for Higgs boson masses larger than about 60 GeV. The dimensionally reduced 3d effective model (e.g. ) was also studied perturbatively and gave similar conclusions. Large scale numerical simulations both on 4d and 3d lattices were needed to analyze the nature of the transition for realistic Higgs boson masses . These results are in complete agreement, and predict an end point for the first order EWPT at Higgs boson mass 72.0$`\pm `$1.4 GeV , above which only a rapid cross-over can be seen. The present experimental lower limit of the SM Higgs boson mass is by several standard deviations larger than the end point value, thus any EWPT in the SM is excluded. This also means that the SM baryogenesis in the early Universe is ruled out.
In order to explain the observed baryon asymmetry, extended versions of the SM are necessary. Clearly, the most attractive possibility is the MSSM. According to perturbative predictions the EWPT could be much stronger in the MSSM than in the SM , in particular if the stop mass is smaller than the top mass . At two-loop level stop-gluon graphs give a considerable strengthening of the EWPT (e.g. third and fourth paper of ). A reduced 3d version of the MSSM has recently been studied on the lattice (including $`\mathrm{SU}(3)\times \mathrm{SU}(2)`$ gauge fields, the right-handed stop and the “light” combination of the Higgses). The results show that the EWPT can be strong enough, i.e. $`v/T_c>1`$, up to $`m_h105`$ GeV and $`m_{\stackrel{~}{t}}165`$ GeV (where $`m_h`$ is the mass of the lightest neutral scalar and $`m_{\stackrel{~}{t}}`$ is that of the stop squark). The possibility of spontaneous CP violation for a successful baryogenesis is also addresed .
In this letter we study the EWPT in the MSSM on 4d lattices and carry out infinite volume and continuum limit extrapolations. Except for the U(1) sector and scalars with small Yukawa couplings, the whole bosonic sector of the MSSM is kept: SU(3) and SU(2) gauge bosons, two Higgs doublets, left-handed and right-handed stops and sbottoms. As it has been done in the SM case , fermions, owing to their heavy Matsubara modes, are included perturbatively in the final result. This work extends the 3d study in two ways:
a) We use 4d lattices instead of 3d. Note, that due to very soft modes—close to the end point in the SM—much more CPU time is needed in 4d than in 3d. However, this difficulty does not appear in the MSSM because the phase transition is strong and the dominant correlation lengths are not that large in units of $`T_c^1`$. Using unimproved lattice actions the leading corrections due to the finite lattice spacings are proportional to $`a`$ in 3d and only to $`a^2`$ in 4d. For O($`a`$) improvement in the 3d case cf. . In 4d simulations we also have direct control over zero temperature renormalization effects.
b) We include both Higgs doublets, not only the light combination. According to standard baryogenesis scenarios (see e.g. ) the generated baryon number is directly proportional to the change of $`\beta `$ through the bubble wall: $`\mathrm{\Delta }\beta `$. ($`\mathrm{tan}\beta =v_2/v_1`$, where $`v_{1,2}`$ are the expectation values of the two Higgses.)
The continuum lagrangian of the above theory in standard notation reads
$$=_g+_k+_V+_{sm}+_Y+_w+_s.$$
(1)
The gauge part, $`_g=1/4F_{\mu \nu }^{(w)}F^{(w)\mu \nu }+1/4F_{\mu \nu }^{(s)}F^{(s)\mu \nu }`$ is the sum of weak and strong terms. The kinetic part is the sum of the covariant derivative terms of the two Higgs doublets ($`H_1,H_2`$), the left-handed stop-sbottom doublet ($`Q`$), and the right-handed stop, sbottom singlets ($`U,D`$): $`_k=(𝒟_\mu ^{(w)}H_1)^{}(𝒟^{(w)\mu }H_1)+(𝒟_\mu ^{(w)}H_2)^{}(𝒟^{(w)\mu }H_2)+(𝒟_\mu ^{(ws)}Q)^{}(𝒟^{(ws)\mu }Q)+(𝒟_\mu ^{(s)}U^{})^{}(𝒟^{(s)\mu }U^{})+(𝒟_\mu ^{(s)}D^{})^{}(𝒟^{(s)\mu }D^{}).`$ The potential term for the Higgs fields reads $`_V=m_{12}^2[\alpha _1|H_1|^2+\alpha _2|H_2|^2(H_1^{}\stackrel{~}{H_2}+h.c.)]+g_w^2/8(|H_1|^4+|H_2|^42|H_1|^2|H_2|^2+4|H_1^{}H_2|^2),`$ for which two dimensionless mass terms are defined, $`\alpha _1=m_1^2/m_{12}^2`$ and $`\alpha _2=m_2^2/m_{12}^2.`$ One gets $`_{sm}=m_Q^2|Q|^2+m_U^2|U|^2+m_D^2|D|^2`$ for the squark mass part, and $`_Y=h_t^2(|QU|^2+|H_2|^2|U|^2+|Q^{}\stackrel{~}{H_2}|^2)`$ for the dominant Yukawa part. The quartic parts containing the squark fields read $`_w=g_w^2/8[2\{Q\}^4|Q|^4+4|H_1^{}Q|^2+4|H_2^{}Q|^22|H_1|^2|Q|^22|H_2|^2|Q|^2]`$ and $`_s=g_s^2/8\left[3\{Q\}^4|Q|^4+2|U|^4+2|D|^46|QU|^26|QD|^2+6|U^{}D|^2+2|Q|^2|U|^2+2|Q|^2|D|^22|U|^2|D|^2\right],`$ where $`\{Q\}^4=Q_{i\alpha }^{}Q_{j\beta }^{}Q_{i\beta }Q_{j\alpha }.`$ The scalar trilinear couplings have been omitted for simplicity. It is straightforward to obtain the lattice action, for which we used the standard Wilson plaquette, hopping and site terms.
The parameter space of the above Lagrangian is many-dimensional. We analyze the effect of the strong sector on the EWPT by using three specific sets of parameters. In one case the strong coupling has its physical value, whereas in the two other cases it is somewhat larger and smaller. The experimental values are taken for the weak and Yukawa couplings, and $`\mathrm{tan}\beta =6`$ is used. For the bare soft breaking masses our choice is $`m_{Q,D}=250`$ GeV, $`m_U`$=0 GeV. Lattice renormalization effects on these masses will be discussed later.
The simulation techniques are similar to those of the SU(2)-Higgs model (overrelaxation and heatbath algorithms are used for each scalar and gauge field); some new methods will be published elsewhere . The analysis is based on finite temperature simulations (in which the temporal extension of the lattice $`L_t`$ is much smaller than the spatial extensions $`L_{x,y,z}`$), and zero temperature ones (with $`L_tL_{x,y,z}`$). For a given $`L_t`$, we fix all parameters of the Lagrangian except $`\alpha _2`$. We tune $`\alpha _2`$ to the transition point, $`\alpha _{2c}`$, where we determine the jump of the Higgs field, the shape of the bubble wall, and the change of $`\beta `$ through the phase boundary. Using $`\alpha _{2c}`$ and the parameter set of the finite temperature case, we perform $`T=0`$ simulations and determine the masses (Higgses and W) and couplings (weak and strong) there. Extrapolations to the continuum limit and to infinite volumes are based on simulations at temporal extensions $`L_t=2,3,4,5`$ and at various lattice volumes for each $`L_t`$, respectively. Approaching the continuum limit, we move on an approximate line of constant physics (LCP), on which the renormalized quantities (masses and couplings) are almost constant, but the lattice spacing approaches zero. Our theory is bosonic, therefore the leading corrections due to finite lattice spacings are expected to be proportional to $`a^2`$. This lattice spacing dependence is assumed for physical quantities in $`a0`$ extrapolations.
We compare our simulation results with perturbation theory. We used one-loop perturbation theory without applying high temperature expansion (HTE). A specific feature was a careful treatment of finite renormalization effects, by taking into account all renormalization corrections and adjusting them to match the measured $`T=0`$ spectrum . We studied also the effect of the dominant $`T0`$ two-loop diagram (“setting-sun” stop-gluon graphs, cf. fifth ref. of ), but only in the HTE. We observed less dramatic enhancement of the strength of the phase transition due to two-loop effects than in . Since the infrared behavior of the setting-sun graphs is not understood, we use the one-loop technique with the $`T=0`$ scheme defined above. This type of one-loop perturbation theory is also applied to correct the measured data to some fixed LCP quantities, which are defined as the averages of results at different lattice spacings, (i.e. our reference point, for which the most important quantity is the lightest Higgs mass, $`m_h`$45 GeV).
Fig. 1. shows the phase diagram in the $`m_U^2`$–T plane. One identifies three phases. The phase on the left (large negative $`m_U^2`$ and small stop mass) is the “color-breaking” (CB) phase. The phase in the upper right part is the “symmetric” phase, whereas the “Higgs” phase can be found in the lower right part. The line separating the symmetric and Higgs phases is obtained from $`L_t=3`$ simulations, whereas the lines between these phases and the CB one are determined by keeping the lattice spacing fixed while increasing and decreasing the temperature by changing $`L_t`$ to 2 and 4, respectively. The shaded regions indicate the uncertainty in the critical temperatures. The phase transition to the CB phase is observed to be much stronger than that between the symmetric and Higgs phases. The qualitative features of this picture are in complete agreement with perturbative and 3d lattice results ; however, our choice of parameters does not correspond to a two-stage symmetric-Higgs phase transition. In this two-stage scenario there is a phase transition from the symmetric to the CB phase at some $`T_1`$ and another phase transition occurs at $`T_2<T_1`$ from the CB to the Higgs phase. It has been argued that in the early universe no two-stage phase transition took place, therefore we do not study this possibility and the features of the CB phase any further.
The bare squark mass parameters $`m_Q^2,m_U^2,m_D^2`$ receive quadratic renormalization corrections. As it is well known, one-loop lattice perturbation theory is not sufficient to reliably determine these corrections, thus we use the following method. We first determine the position of the non-perturbative CB phase transitions in the bare quantities (e.g. the triple point or the T=0 transition for $`m_U^2`$ in Fig. 1). These quantities are compared with the prediction of the continuum perturbation theory, which gives the renormalized mass parameters on the lattice.
Fig. 2 contains the continuum limit extrapolation for the normalized jump of the order parameter ($`v/T_c`$: upper data) and the critical temperature ($`T_c/m_W`$: lower data). The shaded regions are the perturbative predictions at our reference point (see above) in the continuum. Their widths reflect the uncertainty of our reference point, which is dominated by the error of $`m_h`$. Note that $`v/T_c`$ is very sensitive to $`m_h`$, which results in the large uncertainties. Results obtained on the lattice and in perturbation theory agree reasonably within the estimated uncertainties. (It might well be that the $`L_t`$=2 results are not in the scaling region; leaving them out from the continuum extrapolation the agreement between the lattice and perturbative results is even better.)
Based on this agreement we use one-loop perturbation theory without HTE to determine cosmologically allowed regions in the $`m_{\stackrel{~}{t}_R}`$ vs. $`m_h`$ plane of the full MSSM (including fermions, $`m_A=500`$ GeV), see Fig. 3. The two lines for each $`m_{\stackrel{~}{t}_L}`$ (which intersect for lower values of $`m_{\stackrel{~}{t}_L}`$) correspond to upper bounds resulting from $`v_n/T_c=1`$ (steeper curves, B1) and the T=0 maximum MSSM Higgs mass (B2). $`v_n`$ is the non-perturbative Higgs expectation value, assumed to be larger than the perturbative one by 14%, a correction factor obtained in the bosonic model (cf. Fig. 2). For large $`m_Q`$ (e.g. 600 GeV, $`m_{\stackrel{~}{t}_L}`$=630 GeV) the region below B1 is below B2. Decreasing $`m_Q`$ B2 decreases and B1 increases. At $`m_Q`$=560 GeV ($`m_{\stackrel{~}{t}_L}`$=590 GeV) B1 and B2 intersect at the CB value of $`m_{\stackrel{~}{t}_R}`$. Since B2 is almost constant this yields the overall maximum Higgs mass for a successful baryogenesis. For even smaller $`m_Q`$ ($`m_{\stackrel{~}{t}_L}`$) both B1 and B2 are relevant. Note that the maximum Higgs mass corresponds to a finite value of $`m_Q`$560 GeV, yielding $`m_h=103\pm 4`$ GeV (including also the uncertainties from Fig. 2 and the difference between the one and two-loop maximum Higgs mass calculations ).
In order to produce the observed baryon asymmetry, a strong first order phase transition is not enough. According to standard MSSM baryogenesis scenarios the generated baryon asymmetry is directly proportional to the variation of $`\beta `$ through the bubble wall separating the Higgs and symmetric phases. By using elongated lattice ($`2L^2192`$), $`L`$=8,12,16 at the transition point we study the properties of the wall. In our simulation procedure we restrict the length of one of the Higgs fields to a small interval between its values in the bulk phases. As a consequence, the system fluctuates around a configuration with two bulk phases and two walls between them. In order to have the smallest possible free energy, the wall is perpendicular to the long direction. We eliminate the effect of the remaining zero mode by shifting the wall of each configuration to some fixed position. Fig.4 gives the bubble wall profiles for both Higgs fields. The measured width of the wall is \[A+B$`\mathrm{log}(aLT_c)]/T_c`$ A=10.8$`\pm `$.1 and B=2.1$`\pm `$.1. This behavior indicates that the bubble wall is rough and without a pinning force of finite size its width diverges very slowly (logarithmically) . For the same bosonic theory the perturbative approach predicts $`(11.2\pm 1.5)/T_c`$ for the width.
Transforming the data of Fig. 4 to $`|H_2|^2`$ as a function of $`|H_1|^2`$, we obtain $`\mathrm{\Delta }\beta =0.0061\pm 0.0003`$. The perturbative prediction at this point is $`0.0046\pm 0.0010`$. Thus perturbative studies such as are confirmed by non-perturbative results.
To summarize, we presented 4d lattice results on the EWPT in the MSSM. Our simulations were carried out in the bosonic sector of the MSSM. We found quite a good agreement between lattice results and our one-loop perturbative predictions. Using this agreement together with a careful analysis of its uncertainties, we determined the upper bound for the lightest Higgs mass for a successful baryogenesis in the full (bosonic+fermionic) MSSM, which turned out to be ($`103\pm 4`$ GeV) consistent with the 3d analysis ($`105`$ GeV). We analyzed the bubble wall profile separating the Higgs and symmetric phases. The width of the wall and the change in $`\beta `$ is in fairly good agreement with perturbative predictions for typical bubble sizes. Both the upper bound for $`m_h`$ and the smallness of $`\mathrm{\Delta }\beta `$ indicate that experiments allow just a small window for MSSM baryogenesis.
Details of the present analysis will be discussed in a forthcoming publication .
This work was partially supported by Hungarian Science Foundation Grants OTKA-T22929-29803-M28413/FKFP-0128/1997. The simulations were carried out on the 46G PC-farm at Eötvös University. |
warning/0001/gr-qc0001015.html | ar5iv | text | # Dynamics of a string coupled to gravitational waves — Gravitational wave scattering by a Nambu-Goto straight string —
preprint: YITP-00-1, gr-qc/0001015
We study the perturbative dynamics of an infinite gravitating Nambu-Goto string within the general-relativistic perturbation framework. We develop the gauge invariant metric perturbation on a spacetime containing a self-gravitating straight string with a finite thickness and solve the linearized Einstein equation. In the thin string case, we show that the string does not emit gravitational waves by its free oscillation in the first order with respect to its oscillation amplitude, nevertheless the string actually bends when the incidental gravitational waves go through it.
There are significant interests in topological defects formed during phase transition in the early universe. In particular, it has been thought that these defects radiate gravitational wave by their rapid oscillation. Thus, it is crucially important to study the precise dynamics of the defects and their gravitational effects.
In the simplest case, the defects are idealized by the infinitesimally thin Nambu-Goto membranes, i.e., their dynamics is governed by the minimization of their world hyper-volume. If the self-gravity of membranes is ignored (test membrane case), the Nambu-Goto action admits oscillatory solutions. It is considered that the membranes gradually lose their kinetic energy by the gravitational wave emission. However, by taking into account the self-gravity of the Nambu-Goto wall, it is shown that a self-gravitating wall coupled to gravitational wave behaves in a quite different manner. The dynamical degree of freedom concerning the perturbative oscillations around a spherical one is given by that of gravitational waves and self-gravitating spherical walls do not oscillate spontaneously unlike test walls.
How about Nambu-Goto strings? It is also considered that the cosmic strings oscillate rapidly and gradually lose their kinetic energy by the gravitational wave emission. The energy momentum tensor of an oscillating infinite test string and the gravitational wave emission are studied by several authors. However, it is not clear how a Nambu-Goto string behaves when one takes into account of its self-gravity.
In this paper, we consider the perturbative oscillation of an infinite self-gravitating string using a exactly solvable model within the general-relativistic perturbation framework and show that a self-gravitating infinite string behaves in the same manner as the above self-gravitating wall in the first order with respect to the oscillation amplitude of the string.
It is known that the mathematical description of a thin string is more delicate than a domain wall because the support of a string is a surface of co-dimension two. There is no simple prescription of an arbitrary line source where a metric becomes singular. In this paper, we consider, first, a straight string with a finite thickness so that the singularity is regularized. Then, the metric junction formalism is applicable on the surface of the thick string. Next, we consider gravitational wave emission by the thick string motion which is excited by incident gravitational wave, i.e., the scattering problem of gravitational wave by a thick string. We analyze this problem by the gauge invariant linear perturbation theory and show the perturbative velocity of the string is given by the variable of gravitational waves. In the thin string case, we show that there is neither resonance nor phase shift in the gravitational wave scattered by a Nambu-Goto string, nevertheless the string is bent by gravitational waves. This shows that self-gravitating infinite strings do not oscillate spontaneously unlike test string at least in the first order with respect to its oscillation amplitude.
As the background for the perturbation, we consider a spacetime $`(,g_{\mu \nu })`$ containing a straight thick string. The surface $`𝒮`$ of the thick string divides $``$ into two regions: $`_{ex}`$ and $`_{in}`$. Note that $`_{in}`$ describes the ‘thick’ worldsheet of the string. We assume that the spacetime $``$ is static and cylindrically symmetric. Then we divide $``$ into two submanifolds so that $`=_1\times _2`$ and write the background metric on $``$ in the form
$$ds^2=\gamma _{ab}dy^ady^b+\eta _{pq}dz^pdz^q,$$
(1)
where $`\gamma _{ab}`$, the metric on $`_1`$, and $`\eta _{pq}`$, that on $`_2`$, are given by
$$\gamma _{ab}dy^ady^b=d\rho ^2+r(\rho )^2d\varphi ^2,\eta _{pq}dz^pdz^q=dt^2+dz^2,0\varphi 2\pi .$$
(2)
We shall use the indices $`a,\mathrm{},d`$ for tensors on $`_1`$ and $`p,\mathrm{},s`$ for those on $`_2`$. The string thickness is given by the circumference radius $`r_{}`$ of $`𝒮`$.
Since a Nambu-Goto string is characterized by a constant string tension $`\sigma _0`$, we consider the following energy-momentum tensor,
$$T_{\mu \nu }=\sigma \eta _{\mu \nu },$$
(3)
where $`\sigma =\sigma _0`$ for $`r<r_{}`$ and $`\sigma =0`$ for $`rr_{}`$ . Here $`\eta _{\mu \nu }`$ is the four dimensional extension of $`\eta _{pq}`$.
The Einstein equations for the metric (1) and (2) are reduced to the single equation
$$=2\frac{_\rho ^2r}{r}=16\pi G\sigma ,$$
(4)
where $``$ is the Ricci curvature on $`_1`$. The solution on $`_{in}_1`$ is given by
$$\gamma _{ab}dy^ady^b=\frac{dr^2}{1\widehat{\alpha }^2r^2}+r^2d\varphi ^2,\widehat{\alpha }^2=\frac{}{2},$$
(5)
and that on $`_{ex}_1`$ is
$$\gamma _{ab}dy^ady^b=\frac{dr^2}{(1\alpha )^2}+r^2d\varphi ^2,$$
(6)
where $`\alpha `$ is a deficit angle on $`_{ex}`$. These two solutions are joined along the surface $`𝒮`$ by Israel’s junction condition: $`[K_\nu ^\mu ]:=K_{\nu +}^\mu K_\nu ^\mu =0,`$ where $`K_{\nu \pm }^\mu `$ are the extrinsic curvature of $`𝒮`$ facing to $`_{ex}`$ and $`_{in}`$, respectively. For the solutions (5) and (6), this junction condition is reduced to $`\alpha =1\sqrt{1\widehat{\alpha }^2r_{}^2}`$. The global geometry of $`_1`$ is illustrated in Fig.1.
We consider the metric perturbations on the background geometry given by (1), (5) and (6). Let $`h_{\mu \nu }`$ be a perturbative metric and $`t_\nu ^\mu `$ be a perturbed energy-momentum tensor, which can be expanded by the harmonics on $`_2`$ as follows:
$`h_{ab}={\displaystyle f_{ab}S},h_{ap}={\displaystyle \left\{f_{a(o1)}V_{(o1)p}+f_{a(e1)}V_{(e1)p}\right\}},`$ (7)
$`h_{pq}={\displaystyle \left\{f_{(o2)}T_{(o2)pq}+f_{(e0)}T_{(e0)pq}+f_{(e2)}T_{(e2)pq}\right\}},`$ (8)
$`t_b^a={\displaystyle s_b^aS},t_p^a={\displaystyle \left\{s_{(o1)}^aV_{(o1)p}+s_{(e1)}^aV_{(e1)p}\right\}},`$ (9)
$`t_q^p={\displaystyle \left\{s_{(o2)}T_{(o2)}^{}{}_{q}{}^{p}+s_{(e0)}T_{(e0)}^{}{}_{q}{}^{p}+s_{(e2)}T_{(e2)}^{}{}_{q}{}^{p}\right\}}.`$ (10)
Here $`:=𝑑\omega 𝑑k_z`$ and
$`S:=e^{i\omega t+ik_zz},`$ $`V_{(o1)}^p:=ϵ^{pq}\widehat{D}_qS,`$ $`V_{(e1)}^p:=\eta ^{pq}\widehat{D}_qS,`$ (11)
$`T_{(e0)pq}:={\displaystyle \frac{1}{2}}\eta _{pq}S,`$ $`T_{(e2)pq}:=\left(\widehat{D}_p\widehat{D}_q{\displaystyle \frac{1}{2}}\eta _{pq}\widehat{D}^r\widehat{D}_r\right)S,`$ $`T_{(o2)pq}:=ϵ_{r(p}\widehat{D}_{q)}\widehat{D}^rS,`$ (12)
are independent tensor harmonics on $`_2`$, $`ϵ_{rs}`$ is a two-dimensional antisymmetric tensor on $`_2`$, and $`\widehat{D}_p`$ denotes the covariant derivative associated with $`\eta _{pq}`$. The symbols $`(o)`$ and $`(e)`$ refer to odd and even parity modes with respect to the inversion of $`(t,z)`$, respectively. The expansion coefficients are tensors on $`_1`$. The perturbative energy momentum tensor $`t_\nu ^\mu `$ has its support only on $`_{in}`$.
We define $`\kappa ^2:=\omega ^2k_z^2`$, which is the eigen value of a differential operator $`\eta ^{pq}\widehat{D}_p\widehat{D}_q`$. We note that the mode with $`\kappa =0`$, which propagates along the string, is not included in the expansion (7)-(10). The $`\kappa =0`$ mode will be discussed in Ref..
Here we consider the gauge-transformation of $`h_{\mu \nu }`$ and $`t_\nu ^\mu `$ associated with $`x^\mu x^\mu +\xi ^\mu `$, where $`\xi ^\mu `$ is expanded as
$$\xi _a:=\zeta _aS,\xi _p:=\left\{\zeta _{(o1)}V_{(o1)p}+\zeta _{(e1)}V_{(e1)p}\right\}.$$
(13)
Inspecting the gauge transformed variables $`h_{\mu \nu }\mathrm{\pounds }_\xi g_{\mu \nu }`$ and $`t_\nu ^\mu \mathrm{\pounds }_\xi T_\nu ^\mu `$, we find simple gauge-invariant combinations of the expansion coefficients: for odd modes,
$$F_a:=f_{a(o1)}\frac{1}{2}D_af_{(o2)},$$
(14)
and for even modes,
$`F_{ab}:=f_{ab}D_aX_bD_bX_a,F:=f_{(e0)}\kappa ^2f_{(e2)},`$ (15)
where $`D_a`$ is a covariant derivative associated with $`\gamma _{ab}`$ and the variable $`X^a:=f_{(e1)}^a\frac{1}{2}D^af_{(e2)}`$ is transformed to $`X^a\zeta ^a`$ by the gauge transformation.
We also introduce the gauge invariant variables for the perturbations of $`T_\nu ^\mu `$ by
$$\mathrm{\Sigma }:=16\pi G(s_{(e0)}+2X^aD_a\sigma ),V^a:=16\pi G(s_{(e1)}^a\sigma X^a).$$
(16)
Note that all the expansion coefficients except for $`s_{(e0)}`$ and $`s_{(e1)}^a`$ are gauge invariant by themselves. In this article, we consider the perturbative motion of a Nambu-Goto string in the first order with respect to its oscillation amplitude. Within this order, $`\mathrm{\Sigma }`$ is the energy density perturbation which is equal to the tangential tension of a string and $`V^a`$ corresponds to the momentum perturbation. The other coefficients, within the same order, $`s_b^a`$, $`s_{(o1)}^a`$, $`s_{(o2)}`$ and $`s_{(e2)}`$ are regarded as the tension normal to the string worldsheet, spin of the string, Lorentz boost along the string and the energy density perturbation which is not equal to the tangential tension, respectively.
As derived in , $`\mathrm{\Sigma }`$ and $`V^a`$ are induced by the motion of an infinite Nambu-Goto string within the first order of oscillation amplitude, while the others are induced in the higher order. Their results show that the energy momentum perturbations $`s_b^a`$, $`s_{(o1)}^a`$, $`s_{(o2)}`$ and $`s_{(e2)}`$ are irrelevant to the perturbation of an infinite Nambu-Goto string within the first order. This corresponds to the fact that a Nambu-Goto string is only charactorized by its energy density equal to its tangential tension and does not have any other properties such as the tension normal to its worldsheet, spin or boost along itself. Hence, in this paper, we concentrate only on $`\mathrm{\Sigma }`$ and $`V^a`$ and drop the other coefficients in the perturbative energy momentum tensor (9) and (10), since we consider the dynamics of an infinite Nambu-Goto string within the first order of its oscillation amplitude.
In terms of the gauge invariant variables, we write the perturbed Einstein equations: for odd modes,
$`D^aF_a=0,(\mathrm{\Delta }+\kappa ^2)F_aD^cD_aF_c=0,`$ (17)
and for even modes,
$`(\mathrm{\Delta }+\kappa ^2)F_{ab}=F_{ab}+2D_{(a}V_{b)}\gamma _{ab}D_cV^c,(\mathrm{\Delta }+\kappa ^2)F=0,`$ (18)
$`D^cF_{ac}{\displaystyle \frac{1}{2}}D_aF=V_a,F_c^c=0,`$ (19)
where $`\mathrm{\Delta }:=D^aD_a`$. The perturbative divergence of the energy momentum tensor are reduced to
$$\kappa ^2V^a+\frac{1}{2}D^aF=0,D_aV^a+\frac{1}{2}\mathrm{\Sigma }=0.$$
(20)
The first equation corresponds to the Euler equation which coincides with the equation of the perturbative string motion derived from the Nambu-Goto action and the second equation corresponds to the continuity equation for the energy density.
We find that (17)-(20) for even and odd modes on $`_{in}`$ are reduced to the wave equations for two scalar variables $`\mathrm{\Phi }_{(o)}`$ and $`\mathrm{\Phi }_{(e)}`$
$$(\mathrm{\Delta }+\kappa ^2)\mathrm{\Phi }_{(o),(e)}=0,$$
(21)
respectively, and all gauge invariant variables are given by $`\mathrm{\Phi }_{(o)}`$ and $`\mathrm{\Phi }_{(e)}`$ without loss of generality as follows:
$`F_a=ϵ_{ab}D^b\mathrm{\Phi }_{(o)},`$ (22)
$`F_{ab}=\left(D_aD_b{\displaystyle \frac{1}{2}}\gamma _{ab}\mathrm{\Delta }\right)\mathrm{\Phi }_{(e)},F=\mathrm{\Delta }\mathrm{\Phi }_{(e)},`$ (23)
$`V_a={\displaystyle \frac{1}{2}}D_a\mathrm{\Phi }_{(e)},\mathrm{\Sigma }=\mathrm{\Delta }\mathrm{\Phi }_{(e)},`$ (24)
where $`ϵ^{ab}`$ is the two-dimensional antisymmetric tensor on $`_1`$. On $`_{ex}`$, we also find that (17)-(19) are reduced to the same form as (21)-(24) with $`=0`$. The exterior solution $`\mathrm{\Phi }_{(o),(e)}^{(ex)}`$ and the interior solution $`\mathrm{\Phi }_{(o),(e)}^{(in)}`$ to (21) are
$`\mathrm{\Phi }_{(o),(e)}^{(ex)}={\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}e^{im\varphi }\left\{AH_\mu ^{(1)}(\beta r)+BH_\mu ^{(2)}(\beta r)\right\},`$ (25)
$`\mathrm{\Phi }_{(o),(e)}^{(in)}={\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}e^{im\varphi }\left\{CP_\nu ^m(x)+DQ_\nu ^m(x)\right\},`$ (26)
where $`H_\mu ^{(1)}(\beta r)`$ and $`H_\mu ^{(2)}(\beta r)`$ are the Hankel function of the first and the second class, $`P_\nu ^m(x)`$ and $`Q_\nu ^m(x)`$ are the associated Legendre function of the first and second class, and $`\mu :=m/(1\alpha )`$, $`\beta :=\kappa /(1\alpha )`$, $`\nu (\nu +1):=\kappa ^2/\widehat{\alpha }^2`$ and $`x:=\sqrt{1\widehat{\alpha }^2r^2}`$. The coefficients $`A`$ and $`B`$ correspond to the amplitude of the outgoing and the incident wave, respectively. The regularity condition at the axis $`r=0`$ on $`\mathrm{\Phi }_{(o),(e)}^{(in)}`$ leads $`D=0`$.
Now, we construct the global solutions to the perturbed Einstein equations in $``$ by matching the exterior and the interior solutions (25) and (26) along the thick string surface $`𝒮`$ ($`r=r_{}`$). The perturbed solutions should satisfy the perturbed junction conditions $`[\delta q_{\mu \nu }]=0`$ and $`[\delta K_\nu ^\mu ]=0`$, where $`\delta q_{\mu \nu \pm }`$ is the perturbed intrinsic metric and $`\delta K_{\nu \pm }^\mu `$ is the perturbed extrinsic curvature of $`𝒮`$. These perturbed quantities are described by the metric perturbations, which are represented by $`\mathrm{\Phi }_{(e),(o)}`$ through (22) and (23). After some calculations, the perturbed junction conditions tell us
$$[\mathrm{\Phi }_{(o,e)}]=0,[D_{}\mathrm{\Phi }_{(o,e)}]=0,$$
(27)
where $`D_{}=n^aD_a`$ and $`n^a=(/\rho )^a`$.
Substituting (25) and (26) into (27), we have
$`A`$ $`=`$ $`UB,`$ (28)
$`U`$ $`=`$ $`{\displaystyle \frac{\kappa r_{}H_{\mu +1}^{(2)}(\beta r_{})P_\nu ^m(x_{})\left(\sqrt{1x_{}^2}P_\nu ^{m+1}(x_{})+\alpha mP_\nu ^m(x_{})\right)H_\mu ^{(2)}(\beta r_{})}{\kappa r_{}H_{\mu +1}^{(1)}(\beta r_{})P_\nu ^m(x_{})\left(\sqrt{1x_{}^2}P_\nu ^{m+1}(x_{})+\alpha mP_\nu ^m(x_{})\right)H_\mu ^{(1)}(\beta r_{})}},`$ (29)
where $`x_{}:=\sqrt{1\widehat{\alpha }^2r_{}^2}=1\alpha .`$ The absolute value of $`U`$ with $`m=1`$ and $`\alpha =0.3`$ is illustrated in Fig.2.
The deformation of $`𝒮`$ is represented by $`V^a`$ on $`𝒮`$. By the appropriate choice of the function $`\zeta ^a`$, we fix the gauge freedom $`X^aX^a\zeta ^a`$ in the neighborhood of $`𝒮`$ so that
$$X_a|_𝒮:=X_{a\pm }=\frac{1}{}V_a|_𝒮=\frac{1}{2}D_a\mathrm{\Phi }_{(e)}|_𝒮.$$
(30)
Since $`iV_a\omega S`$ is a precise momentum perturbation, $`X_a|_𝒮S`$ does represent the deformation of $`𝒮`$.
Now, we consider the thin string case. Physically, a “thin string” means a string whose thickness $`r_{}`$ is sufficiently smaller than the wavelength of gravitational wave. In this paper, we consider the situation $`ϵ:=\beta r_{}1`$ with the finite outside deficit angle $`\alpha `$ and take the leading order of $`ϵ`$ for the thin string case.
Further, we note that only $`m=1`$ mode in (25) and (26) shows the motion of a Nambu-Goto thin string. $`m=0`$ and $`m>1`$ modes are irrelevant to a thin string. Hence, in the thin string case, the displacement $`X_S^a`$ of a Nambu-Goto string by the gravitational wave scattering is given by
$$X_S^a:=X^a|_𝒮(m=1,ϵ1)S=\frac{\kappa B_{m=1}S}{2(1\alpha )\mathrm{\Gamma }(\frac{2\alpha }{1\alpha })}\left(\frac{ϵ}{2}\right)^{\frac{\alpha }{1\alpha }}e^{i\varphi }\left(n^a+i\tau ^a\right),$$
(31)
where $`\tau ^a=(1/r)(/\varphi )^a`$ and $`n^a=(1\alpha )(/r)^a`$. (31) shows that the string is deformed while the incident wave exists on the string. In the same order calculation where $`X_S^a`$ is given by (31), we obtain the trivial scattering data $`U1`$ from (29).
The order of magnitude of $`|X_S^a|:=\sqrt{\gamma _{ab}X_S^aX_S^b}`$ is estimated as follows:
$$|X_S^a|/r_{}\kappa ^2|B_{m=1}|ϵ^{\alpha /(1\alpha )}/(\kappa r_{})\kappa ^2|B_{m=1}|ϵ^{\alpha /(1\alpha )1}.$$
(32)
The amplitude of gravitational wave ($`F`$ or $`F_{ab}`$) is proportional to $`|\kappa ^2B|1`$. If $`ϵ^{1\alpha /(1\alpha )}<\kappa ^2|B|1`$, then $`|X_S^a|>r_{}`$. Therefore, the magnitude of the displacement may become larger than the string thickness within the linear perturbation framework.
Thus, we have obtain the result that there is neither resonance nor phase shift in the scattering problem, nevertheless the string is deformed by the gravitational waves. In particular, the fact that there is no resonance means an infinite thin string does not emit gravitational wave spontaneously by its oscillatory motion. It should be noted that the trivial scattering does not dictate the absence of gravitational lensing effect by the deficit angle $`\alpha `$ on $`_{ex}`$. We will explicitly see the lensing effect by the scattering of wave packet which is suitably constructed by (25) because the mode functions already include the effect of $`\alpha `$.
Using the linearized Einstein equation, we have found that the perturbative string displacement $`X_S^a`$ is represented by the gravitational wave $`\mathrm{\Phi }_{(e)}`$. In this sense, the dynamical degree of freedom of the string displacement is given by that of the gravitational waves on the string surface. Further, (30) and the trivial scattering data show that the string is bent while the incident wave is passing through the string worldsheet. These behaviors are same as that for a self-gravitating spherical Nambu-Goto wall in the first order with respect to its oscillation amplitude . This is our main conclusion.
To obtain the above results, we have first considered the scattering by a thick string. We regard that a “thin string” is not a string with the thickness $`r_{}0`$ but that whose thickness $`r_{}`$ is sufficiently smaller than the wavelength of gravitational wave. For a fixed amplitude of the incident wave, the magnitude of the string displacement $`X_S^a`$ depends on $`r_{}`$. If we take the limit $`r_{}0`$ for fixed the wavelength of gravitational wave, these is no response of string motion for finite incident gravitational wave. This result is consistent with that obtained by Unruh et.al.. Both their and our results show that the straight string cannot bend in the limit $`r_{}0`$. This mathematical limit will be irrelevant for strings formed during phase transition in the early universe, because they have finite thickness. Further, the scattering data (29) has the resonance poles at $`\beta r_{}1`$ or larger as seen in Fig.2. This suggest that a cosmic string oscillates spontaneously and emits the gravitational waves with the frequencies of the order of the string thickness. In this situation, the dynamics of strings is no longer approximated by that of thin Nambu-Goto strings.
One might think that our result depends sensitively on the distribution of $`\sigma `$ in (3) and our model would be too artificial because of the step function distribution of $`\sigma `$. Further, the above analysis does not include $`\kappa =0`$ mode which includes “cosmic string traveling waves” discussed in Ref., and one might think that the dynamical degree of freedom of the string spontaneous oscillation is in this $`\kappa =0`$ mode. However, we obtain the same conclusion in the thin string case even when the background $`\sigma `$ is different from the step function and our conclusion is unchanged in the thin string case even if we include the $`\kappa =0`$ mode into our consideration. These two points will be discussed in a separated paper .
The authors thank Professor Minoru Omote and Professor Akio Hosoya for their continuous encouragement. This work was partially supported by Soryushi Shogakukai and Yukawa Shogakukai (A.I.). |
warning/0001/astro-ph0001466.html | ar5iv | text | # LiBeB Evolution: Three Models
## 1. Introduction
Cosmic-ray driven nucleosynthesis has been known to be important for the origin of the light elements Li, Be and B (LiBeB) for three decades (Reeves, Fowler, & Hoyle 1970). But constraints on the relevant evolutionary models could only be obtained after LiBeB abundances of low metallicity stars started to become available (e.g. Ryan et al. 1990). The principal three models currently considered are: (i) the cosmic-ray interstellar model (hereafter CRI), in which the cosmic-ray source composition at all epochs of Galactic evolution is assumed to be similar to that of the average ISM at that epoch (Vangioni-Flam et al. 1990; Fields & Olive 1999); (ii) the CRI+LECR model, in which metal enriched low energy cosmic rays (LECRs) are superimposed onto the CRI cosmic rays (Cassé, Lehoucq, & Vangioni-Flam 1995; Vangioni-Flam et al. 1996; Ramaty, Kozlovsky, & Lingenfelter 1996); and (iii) the cosmic-ray superbubble model (hereafter CRS), in which the cosmic-ray source composition is taken to be constant, independent of the ISM metallicity (Ramaty et al. 1997;2000; Higdon, Lingenfelter, Ramaty 1998). Both the CRS cosmic rays and the LECRs are thought to be accelerated out of supernova enriched matter in superbubbles.
Because of the excess of the observed Be abundances in low metallicity stars over the predictions of the CRI model, and motivated by reports of the detection of C and O nuclear gamma-ray lines from the Orion star formation region (Bloemen et al. 1994), LECRs, enriched in C and O relative to protons and $`\alpha `$ particles, were superimposed on the CRI cosmic rays, hence the CRI+LECR model . These LECRs, with maximum energies not exceeding about 100 MeV/nucleon, were thought (e.g. Ramaty 1996) to be responsible for the gamma rays reported from Orion. It was suggested that such enriched LECRs might be accelerated out of metal-rich winds of massive stars and supernova ejecta (Bykov & Bloemen 1994; Ramaty et al. 1996; Parizot, Cassé, & Vangioni-Flam 1997) by an ensemble of shocks in superbubbles (Bykov & Fleishman 1992; Parizot Cassé, & Vangioni-Flam 1997). The Orion gamma-ray data, however, have been retracted (Bloemen et al. 1999). Nonetheless, as the possible existence of the postulated LECRs remains, new gamma-ray line data are needed to determine the role of LECRs in LiBeB production.
Recent O abundance data, which suggest that \[O/Fe\] increases with decreasing \[Fe/H\] at low metallicities (Israelian et al. 1998; Boesgaard et al. 1999), led Fields and Olive (1999) to reexamine the viability of the CRI model. More recent measurements (Fulbright & Kraft 1999; Westin et al. 1999) argue against such an \[O/Fe\] increase. But as demonstrated in (Ramaty et al. 2000), and also shown below, this model is inconsistent with cosmic-ray energetics, an \[O/Fe\] increase notwithstanding.
Alternatively, it was suggested (Lingenfelter, Ramaty, & Kozlovsky 1998; Higdon et al. 1998), that the Be evolution can be best understood in the CRS model, in which the cosmic-ray metals at all epochs of Galactic evolution are accelerated predominantly out of supernova ejecta. Lingenfelter et al. (1998) and Lingenfelter & Ramaty (1999) showed that the arguments (e.g. Meyer, Drury, & Ellison 1997) against the supernova ejecta origin of the current epoch cosmic rays can be answered, and Higdon et al. (1998) and Higdon, Lingenfelter, & Ramaty (1999) showed that the most likely scenario is collective acceleration by successive supernova shocks of ejecta-enriched matter in the interiors of superbubbles. This scenario is consistent with the delay between nucleosynthesis and acceleration (time scales $``$10<sup>5</sup> yr), suggested by the <sup>59</sup>Co and <sup>59</sup>Ni observations (Wiedenbeck et al. 1999). In both the CRS and CRI+LECR models, the bulk of the Be in the early Galaxy is produced by accelerated C and O interacting with ambient H and He. That these “inverse reactions” are dominant in the early Galaxy was first suggested by Duncan, Lambert, & Lemke (1992).
In the present paper we present results from a complete set of LiBeB evolutionary calculations for all three models using our production code described in Ramaty et al. (1997) and evolutionary code detailed in Ramaty et al. (2000).
## 2. Analysis
Figure 1 shows the evolution of \[O/Fe\] as a function of \[Fe/H\]. In order to account for the possible rise of \[O/Fe\] with decreasing \[Fe/H\], we introduced mixing delays, i. e. the delayed deposition of the synthesized products into the star forming regions due to differences in transport and mixing. We choose a short mixing time for oxygen because we expect the bulk of the O and other volatiles in the ejecta to mix with the ISM after the remnant slows down to local sound speeds. But we consider longer mixing times for Fe, assuming that the bulk of the ejected Fe is incorporated into high velocity refractory dust grains which continue moving for longer periods of time before they stop and can be incorporated into newly forming stars. The incorporation of a large fraction of the synthesized Fe into dust grains is supported by observations of both supernova 1987A and the Galactic 1.809 MeV gamma-ray line resulting from the decay of <sup>26</sup>Al (for more details and references see Ramaty et al. 2000). We see that with the mixing delays (solid curves in Figure 1) both the WW95 and TS cases (see figure caption) become consistent with the Israelian et al. (1998) and Boesgaard et al. (1999a) data, showing that delayed Fe deposition could indeed be the cause for the rise of \[O/Fe\]. In this connection, it is interesting to note that, unlike \[O/Fe\], the abundance ratios of the $`\alpha `$-nuclei Mg, Si, Ca and Ti relative to Fe do not increase with decreasing \[Fe/H\] below \[Fe/H\] $`=1`$ (Ryan, Norris, & Beers 1996). This may be consistent with the fact that these elements are also refractory, and thus are affected by mixing in the same way as is Fe. A test may be provided by sulfur, which is volatile, and thus should show a rise similar to the rise of \[O/Fe\] vs. decreasing \[Fe/H\] (G. Israelian, private communication, 1999).
Figure 2 shows the employed C and O abundances of the ISM and the cosmic ray source. For the ISM (left panel) we take \[C/H\]$``$\[Fe/H\], because at early times both the C and Fe come primarily from core collapse supernovae of massive stars, while at later times the increased C contribution from the winds of intermediate mass stars is compensated by the Fe contribution from the thermonuclear supernovae of the white dwarf remnants of such stars (see Timmes, Woosley & Weaver 1995). The O abundances follows from the results of evolutionary calculations shown in Figure 1. Unlike in Ramaty et al. (2000), where we used a constant He abundances, here we allow He/H (by number) to vary slowly from 0.08 at very low metallicities to 0.1 at \[Fe/H\]=0. For the cosmic-ray source (right panel), we define the logarithmic ratios \[C/H\] and \[O/H\] in the same way as is done for the corresponding ISM values, including normalization to solar (not current epoch cosmic ray source) abundances. As the CRS cosmic rays are accelerated primarily out of supernova ejecta enriched superbubbles, \[C/H\] and \[O/H\] are constant, set equal to current epoch cosmic-ray values. For the CRI model we scale \[C/H\] and \[O/H\] to the ISM values with enhancement factors of 1.5 and 2, consistent with the mass-to-charge dependent acceleration of volatiles (Ellison, Drury & Meyer 1997), except at \[Fe/H\]=0 where the cosmic ray CRS and CRI values are equal. For details on the rest of the employed cosmic-ray sources abundances see Ramaty et al. (2000). For the LECRs we adopt the CRS abundances. The CRS, CRI and LECR source energy spectra are power laws in momentum with high energy exponential cutoffs (characteristic energy $`E_0`$), which we set to an ultrarelativistic value for the CRS and CRI cosmic rays and to 30 MeV/nucleon for the LECRs.
Figure 3 shows the resultant $`Q/W`$’s, the total number of nuclei $`Q`$ produced by an accelerated particle distribution normalized to the integral cosmic-ray energy $`W`$, for a given source energy spectrum, and cosmic ray and ambient medium compositions as described above. We make the reasonable assumption that the accelerated particle source energy spectrum is independent of \[Fe/H\]. The resultant CRS and CRI $`Q/W`$’s for Be and <sup>6</sup>Li (left panel) are for $`\mathrm{X}_{\mathrm{esc}}=10\mathrm{gcm}^2`$, typical of currently inferred values for “leaky box” cosmic-ray propagation models. While $`\alpha `$$`\alpha `$ dominated $`Q(^6\mathrm{Li})/W`$ is not very different for the CRS and CRI models, $`Q(\mathrm{Be})/W`$ is drastically different for the two models, reflecting the fact that efficient Be production in the early Galaxy can only result from C and O enriched accelerated particles. The different O abundances employed in the calculation of $`Q/W`$ for the two mixing delays cases (see Figure 2) lead to significantly different $`Q(\mathrm{Be})/W`$’s for CRI model. The LECR $`Q/W`$’s (right panel) show that, while removal of the protons and $`\alpha `$ particles (the CRS(metal) composition) significantly increases $`Q/W`$ for Be, it essentially leaves $`Q(^6\mathrm{Li})/W`$ unchanged, because the lack of <sup>6</sup>Li production by $`\alpha `$ particles is compensated by a smaller W due to the absence of the protons and alphas.
Figure 4 shows the Be evolution for the CRS and CRI models. In the calculations, 10<sup>50</sup> erg per supernova are imparted to the cosmic rays, a value in very good agreement with current epoch cosmic-ray energetics. We note that even though the overall slope of log(Be/H) vs. \[Fe/H\] is practically unity, while that of log(Be/H) vs. \[O/H\] is significantly steeper (0.96$`\pm `$0.04 and 1.45$`\pm `$0.04, respectively, Boesgaard et al. 1999b), the CRS model provides a good fit to these evolutionary trends, particularly if n<sub>H</sub> is near 0.1. Such a low value might not be unreasonable for an average halo hydrogen density if 10<sup>10</sup>M are spread over a few kpc<sup>3</sup>. The calculated log(Be/H) vs. \[Fe/H\] is flatter than log(Be/H) vs. \[O/H\] because the delayed deposition of the synthesized Be, caused by the low n<sub>H</sub>, is compensated by the delayed Fe deposition, due to the incorporation of Fe in high velocity dust, but not compensated by the very short delay of the deposition of O, which is mostly volatile. As in Ramaty et al. (2000), we see that the CRI model, normalized to a reasonable energy in cosmic rays per supernova, severely underproduces the measured Be abundances. However, unlike in that paper where we showed the result only for log(Be/H) vs. \[Fe/H\], here we show that the same result also holds for log(Be/H) vs. \[O/H\]. This removes the remaining ambiguity concerning our argument against the result of Fields & Olive (1999), who claimed that the CRI model would be viable if instead of Fe ejecta per supernova based on calculations, which are somewhat uncertain, they used values based on their fit to the increasing \[O/Fe\] with decreasing \[Fe/H\]. As there is no such uncertainty concerning the O ejected masses, our present result unequivocally demonstrates that the CRI model is untenable.
Figure 5 shows the Be evolution for the CRI+LECR model. Here, as before, 10<sup>50</sup> erg per supernova are imparted to the CRI cosmic rays, but in order to achieve a fit to the log(Be/H) vs. \[Fe/H\] data, we had to add more energy (1.5$`\times `$10<sup>50</sup> erg per supernova) in LECRs. The need for this increased energy can be seen in Figure 3, where the $`Q(\mathrm{Be})/W`$ for the LECR model with the CRS composition (right panel) is lower by about a factor of 2 than the corresponding value for the CRS model (left panel). Returning to Figure 5, we see that the CRI+LECR model leads to a log(Be/H) vs. \[O/H\] evolutionary curve with a slope which is only slightly steeper than 1, while, as mentioned above, the data indicate a slope of 1.45$`\pm `$0.04 (Boesgaard et al. 1999b). Indeed, the simplest evolutionary considerations for both the CRS and LECR models would predict that log(Be/H) vs. \[O/H\] should have a slope of 1. But, as shown above, the delay introduced by cosmic-ray transport with n<sub>H</sub>=0.1 will steepen the slope. However, the LECRs slow down much faster than the higher energy CRS cosmic rays. Thus, the calculation of Figure 5 assigns no delay to the Be deposition, leading to a possible inconsistency between the predictions of the CRI+LECR model and the log(Be/H) vs. \[O/H\] data.
Figure 6 shows the boron isotopic ratio (left panel) and log(B/Be) (right panel) vs. Fe/H for the CRS and CRI+LECR models. In the evolutionary calculations, the production ratios $`Q(^{11}\mathrm{B}/^{10}\mathrm{B})=2.4`$ and $`Q(\mathrm{B}/\mathrm{Be})=14`$ for the CRS and CRI cosmic rays, and $`Q(^{11}\mathrm{B}/^{10}\mathrm{B})=3.3`$ and $`Q(\mathrm{B}/\mathrm{Be})=22`$ for the LECRs, were taken as independent of \[Fe/H\]. These numerical values are from Ramaty et al. (1997) and are valid for high energy CRS and CRI cosmic rays and LECRs with $`E_0=30`$ MeV/nucleon. In order to reproduce the meteoritic <sup>11</sup>B/<sup>10</sup>B both the CRS and CRI+LECR models require the addition of $`\nu `$-produced <sup>11</sup>B (Woosley & Weaver 1995). We take into account the metallicity dependence of this <sup>11</sup>B production, and the fact that only the core collapse supernovae produce <sup>11</sup>B by neutrinos. We find that for the CRS model the meteoritic data can be fit with $`f_\nu `$=0.18, a lower value than we found in Ramaty et al. (2000) because of the lower cosmic-ray energy per supernova that we use here (10<sup>50</sup> vs. 1.5$`\times `$10<sup>50</sup> erg that we used in that paper). The rise in <sup>11</sup>B/<sup>10</sup>B and B/Be for n<sub>H</sub> = 0.1 below \[Fe/H\] of about $`2`$ in the CRS model is due to the delayed deposition of the cosmic-ray produced Be and B relative to the $`\nu `$-produced <sup>11</sup>B for which we took a short delay time, the same as for O (1 Myr). For the CRI+LECR model, in which the LECRs employ a larger energy per supernova, $`f_\nu `$ = 0.2 is required to fit the boron isotope data. The rise in <sup>11</sup>B/<sup>10</sup>B and B/Be with decreasing \[Fe/H\] is mostly due to the larger value of the corresponding LECR production ratios. It is evident from Figure 6 that future high precision measurements of <sup>11</sup>B/<sup>10</sup>B and B/Be as functions of \[Fe/H\] will distinguish between the models.
Figure 7 shows the Li evolution. The CRS model underproduces the <sup>6</sup>Li abundance for \[Fe/H\]$`<2`$, suggesting the existence of pregalactic or extragalactic <sup>6</sup>Li sources. With the lower cosmic-ray normalization mentioned above, the CRS model also slightly underproduces the meteoritic <sup>6</sup>Li. With the discovery of solar flare produced <sup>6</sup>Li in the solar wind (via measurements in lunar soil, Chaussidon & Robert 1999), the possibility of some locally produced <sup>6</sup>Li in the solar system must be considered. The CRI+LECR model produces more <sup>6</sup>Li relative to Be, simply because the $`\alpha `$$`\alpha `$ cross section for <sup>6</sup>Li production peaks in the nonrelativistic region.
## 3. Discussion and Conclusions
We have summarized a complete set of O, Fe and LiBeB evolutionary calculation. We have considered the three principal evolutionary models, CRI in which the cosmic-ray source composition at all epochs of Galactic evolution is similar to that of the average ISM at that epoch, CRI+LECR in which metal enriched low energy cosmic rays (LECRs) are superimposed onto the CRI cosmic rays, and CRS in which the cosmic-ray source, accelerated in superbubbles, has a constant composition, independent of the ISM metallicity. By considering the evolutionary trend of log(Be/H) vs. both \[Fe/H\] and \[O/H\], we demonstrated that the CRI model is energetically untenable. Although the CRI+LECR mix considered here is consistent with the Be evolution vs. Fe, a plausible scenario for producing the required mix has yet to be proposed.
For Fe, our code allows for a delay between nucleosynthesis and deposition into star forming regions due to the incorporation of the synthesized Fe into high velocity dust. This delay could provide an explanation for the possible rise of \[O/Fe\] with decreasing \[Fe/H\] indicated by some of the data. A test for this scenario would be the demonstration that the abundances of refractory $`\alpha `$-elements Mg, Si, Ca and Ti relative to Fe do not increase with decreasing \[Fe/H\] below \[Fe/H\]$`=1`$, but that volatile sulfur does rise.
For the LiBeB there is also a delay. Due to the transport of the cosmic rays, LiBeB synthesis lags behind the explosion of the supernova responsible for accelerating the cosmic rays by as much a hundred million years, depending on the average gas density in the halo of the early Galaxy. We show that this delay, combined with the delayed Fe deposition, could provide an explanation for the steeper evolutionary trend of log(Be/H) vs. \[O/H\] than vs. \[Fe/H\], as indicated by the recent data of Boesgaard et al. (1999b).
The trends of <sup>11</sup>B/<sup>10</sup>B and B/Be vs. \[Fe/H\] show structure resulting from the above mentioned delays, as well as from the hybrid nature of the CRI+LECR model. Future observations of these ratios may distinguish between the models. The fact that the <sup>6</sup>Li abundances at \[Fe/H\] below $`2`$ seem inconsistent with all the models suggests the existence of pregalactic or extragalactic <sup>6</sup>Li sources.
## References
Bloemen, H. et al. 1994, A&A, 281, L5
Bloemen, H. et al. 1999, ApJ, 521, L137
Boesgaard, A., King, J., Deliyannis, C., & Vogt, S. 1999a, AJ, 117, 492
Boesgaard, A., Deliyannis, C., King, J., Ryan, S., Vogt, S., & Beers, T. 1999b, AJ, 117, 1548
Bykov, A., & Bloemen, H. 1994, A&A, 283, L1
Bykov, A. M., & Fleishman, G. D. 1992, MNRAS, 15, 269
Cassé, M., Lehoucq, R., & Vangioni-Flam, E. 1995, Nature, 373, 318
Chaussidon, M., & Robert, F. 1995, Nature, 374, 337
Chaussidon, M., & Robert, F. 1999, Nature, 402, 270
Duncan, D. K., Lambert, D. L., & Lemke, M., 1992, ApJ, 401, 584
Duncan, D. K. et al. 1997, ApJ, 488, 338
Ellison, D. C., Drury, L.O’C., & Meyer, J-P. 1997, ApJ, 487, 197
Fields, B. D., & Olive, K. A. 1999, ApJ, 516, 797
Fulbright, J. P, & Kraft, R. P. 1999, AJ, 118, 527
Garcia Lopez, R. J. et al. 1998, ApJ, 500, 241
Grevesse, N., Noels, A., & Sauval, A. J. 1996, in: Cosmic Abundances, eds. S. S. Holt, & G. Sonneborn, ASP Conf. Ser., 99, (San Francisco: ASP), 117
Higdon, H. C., Lingenfelter, R. E., & Ramaty, R., 1998, ApJ, 509, L33
Higdon, H. C., Lingenfelter, R. E., & Ramaty, R., 1999, 26th Internat. Cosmic Ray Conf., eds. D. Kieda et al. (Salt Lake City), 4, 144
Hobbs, L. M. 1999, in: LiBeB, Cosmic Rays, and Related X- and Gamma-Rays, eds. R. Ramaty et al., ASP Conf. Ser., 171, (San Francisco: ASP), 23
Israelian, G., Garcia Lopez, R. J., & Rebolo, R. 1998, ApJ, 507, 805
Lambert, D. L. et al. 1998, ApJ, 494, 614
Lingenfelter, R. E., & Ramaty, R., 1999, 26th Internat. Cosmic Ray Conf., eds. D. Kieda et al. (Salt Lake City), 4, 148
Lingenfelter, R. E., Ramaty, R., & Kozlovsky, B. 1998, ApJ, 500, L153
Meyer, J-P., Drury, L. O’C., & Ellison, D. C. 1997, ApJ, 487, 182
Molaro, P. 1999, in: LiBeB, Cosmic Rays, and Related X- and Gamma-Rays, eds. R. Ramaty et al. , ASP Conf. Ser., 171, (San Francisco: ASP), 6
Parizot, E. M. G., Cassé, M., & Vangioni-Flam, E. 1997, A&A, 328,107
Ramaty, R. 1996, A&A(Suppl.), 120, C373
Ramaty, R., Kozlovsky, B., & Lingenfelter, R. E. 1996, ApJ, 456, 525
Ramaty, R., Kozlovsky, B., Lingenfelter, R., & Reeves, H. 1997, ApJ, 488, 730
Ramaty, R., Scully, S., Lingenfelter, R., & Kozlovsky, B. 2000, ApJ, in press
Reeves, H., Fowler, W. A., & Hoyle, F. 1970, Nature, 226, 727
Ryan, S. G., Norris, J. E., & Beers, T. C. 1996, ApJ, 471, 254
Ryan, S., Bessell, M., Sutherland, R., & Norris, J. 1990, ApJ, 348, L57
Timmes, F. X., Woosley, S. E., & Weaver, T. A., 1995. ApJ(suppl), 98, 617
Tsujimoto, T., & Shigeyama, T. 1998, ApJ, 508, L151
Vangioni-Flam E., Cassé, M., Audouze, J., & Oberto, Y. 1990, ApJ, 364, 568
Vangioni-Flam, E., Cassé, M., Fields, B., & Olive, K. 1996, A&A, 468, 199
Westin, J. et al. 1999, ApJ, in press (astro-ph/9910376)
Wiedenbeck, M. E. et al. 1999, ApJ(Letters), 523, L61
Woosley, S. E. & Weaver, T. A. 1995, ApJS, 101, 181 |
warning/0001/hep-th0001132.html | ar5iv | text | # D-branes at Singular Curves of Calabi-Yau Compactifications
## 1 Introduction
During the last years our understanding of non-perturbative aspects of string- and field theories has improved tremendously. In this process, the analysis of D-branes played a prominent role. While many of these investigations have been devoted to string theories in a flat background, it has been tried to extend this knowledge to curved backgrounds preserving fewer supersymmetries using various approaches. In this context, Calabi-Yau (CY) compactifications are of particular interest. There exist essentially two ways to approach $`D`$-branes in CY-manifolds: in the large radius limit we can describe them geometrically as branes wrapping around holonomy cycles . It has been found that D-branes in a geometrical phase are naturally described in terms of K-theory classes . On the other hand, these geometrically intuitive concepts are not available in the stringy regime. Here, we can employ methods of boundary conformal field theory (CFT) instead to study D-branes at Gepner points . The comparison of these two approaches has been initiated in (see also for more recent work in this direction).
The present work focuses mainly on extending the CFT constructions of D-branes at the Gepner points of CY compactifications. While a large class of such D-branes was obtained in , a closer investigation of their RR-charges shows that none of them corresponds to branes wrapping the exceptional cycles which appear from the resolution of possible singularities. Since there are only a few non-singular cases (of which the quintic is the most prominent), one would certainly hope that CFT-techniques can provide additional D-branes for almost all Gepner models. This is the problem we are going to address below. In particular, we will construct so-called boundary states that are associated with branes wrapping the additional three-cycles induced by the resolution of a $`_2`$-singularity over a curve in the CY-space.
D-branes in a flat space with orbifold singularities have been studied extensively in the string theory literature, starting with the work of Douglas and Moore . In particular, it is known how to describe D-branes wrapping the collapsed cycles at the orbifold point. The main idea is to attach Chan-Paton factors to the ends of the open strings and define an orbifold action on them. Boundary states corresponding to branes away from the orbifold fixed points are obtained by summing over the brane’s pre-images in the covering space. At the fixed points, however, the expressions for boundary states can involve contributions from twisted sectors of the theory, leading to a charge under RR-potentials coming from the twisted sector. We shall see these concepts reappearing in the non-geometrical CFT analysis.
Our analysis of D-branes in Gepner models will be formulated in the framework of simple current orbifolds. For closed string theories, the necessary techniques were developed long ago in . Open string descendants of simple current orbifolds were investigated systematically by Fuchs and Schweigert and more recently in . Some non-trivial examples have also been studied previously in . The general results show that the boundary states of can be further ‘resolved’, if the action of an appropriate set $`\mathrm{\Gamma }`$ of simple currents on conformal families of the theory possesses short orbits, i.e. orbits of length being less than the order $`|\mathrm{\Gamma }|`$ of $`\mathrm{\Gamma }`$.
Some general aspects of simple current orbifolds will be discussed in the next section. In particular, we shall explain how one can obtain a large set of D-brane states in the orbifold theory by an appropriate projection. For reasons to become clear later, the branes that result from this construction will be called untwisted D-branes. In Section 3 we illustrate the whole procedure at the example of A-type boundary states in Gepner models, thereby recovering precisely the A-type D-brane states listed in . Geometrically, these correspond to branes wrapping middle dimensional cycles on the CY-manifold . As we argue in Section 4, the untwisted D-branes do not wrap exceptional cycles. This motivates our search for additional D-brane states in Section 5. There we will show that untwisted D-branes at a $`_2`$ singularity over some curve $`C`$ can be further resolved. Explicit formulas for the associated boundary states and the open string partition functions are provided. We finally conclude with a number of remarks on possible extensions. These include the analysis of B-type boundary states and of branes wrapping $`_N`$-singularities. While the present techniques do not suffice to resolve branes at $`_{N2}`$-singularities, they can be used to study B-type boundary states and the comparison with the results of provides new evidence for mirror symmetry in the open string sector.
## 2 Orbifolds and untwisted D-branes
The aim of this section is to review some results on simple current orbifolds and a general method for the construction of D-branes in the untwisted sector of the orbifold theory. This will be applied to Gepner models below.
### 2.1 Simple current orbifolds - the bulk theory
The simple current techniques developed in allow to construct new modular invariants from existing ones. A well-known class of examples is provided by Gepner models , where the GSO-projected partition function is obtained using the spectral flow operator as a simple current. Other applications in string theory include the construction of $`(0,2)`$ models which lead to $`𝒩=1`$ space-time supersymmetric theories in four uncompactified dimensions .
In this subsection, we briefly summarize the main results of . Consider some given bulk theory with a bosonic chiral algebra $`𝒲`$. We label classes of irreducible representations of $`𝒲`$ by labels $`i,j,k`$ taken from an index set $`𝒥`$. Within $`𝒥`$ we may find some non-trivial classes $`g𝒥`$ such that the fusion product of $`g`$ with any other $`j𝒥`$ gives again a single class $`gj=gj𝒥`$. Such classes $`g`$ are called simple currents and the set $`𝒞`$ of all these simple currents forms an abelian subgroup $`𝒞𝒥`$. The product in $`𝒞`$ is inherited from the fusion product of representations. From now on, let us fix some subgroup $`\mathrm{\Gamma }𝒞`$.
Through the fusion of representations, the index set $`𝒥`$ comes equipped with an action $`\mathrm{\Gamma }\times 𝒥𝒥`$ of the group $`\mathrm{\Gamma }`$ on labels $`j𝒥`$. Under this action, $`𝒥`$ may be decomposed into orbits. The space of these orbits will be denoted by $`𝒥/\mathrm{\Gamma }`$ and we use the symbol $`[j]_\mathrm{\Gamma }`$ to denote the orbit represented by $`j𝒥`$. By construction, the length of the orbit of the identity $`1𝒥`$ is given by the order $`|\mathrm{\Gamma }|`$ of the group $`\mathrm{\Gamma }`$. Other orbits may be shorter since there can be fixed points, i.e. labels $`j𝒥`$ for which
$$gj=j\text{ for some }g\mathrm{\Gamma }.$$
(1)
The subgroup of all simple currents leaving some $`j𝒥`$ fixed is called the stabilizer of $`j`$:
$$𝒮_j=\{g\mathrm{\Gamma }gj=j\}$$
(2)
Two labels $`j_1,j_2`$ possess isomorphic stabilizers, if they are on the same orbit, i.e. if there exists an element $`g\mathrm{\Gamma }`$ such that $`gj_1=j_2`$.
For all simple currents $`g\mathrm{\Gamma }`$ and all labels $`j𝒥`$ we define the monodromy charge $`Q_g(j)𝐑/𝐙S^1`$ of $`j`$ with respect to $`g`$ by
$$Q_g(j):=h_g+h_jh_{gj}\text{mod}1.$$
(3)
Here, $`h_l`$ denotes the non-integer part of the conformal dimension of the conformal primary that is associated with $`l𝒥`$. The meaning of $`Q_g(j)`$ can be easily understood once we choose two fields $`J(z)`$ and $`\psi (w)`$ from the conformal families $`g\mathrm{\Gamma }𝒥`$ and $`j𝒥`$, respectively. Their operator product will give fields $`\psi ^{}(w)`$ within a single conformal family $`gj`$, i.e.
$$J(z)\psi (w)(zw)^{h_J+h_\psi h_\psi ^{}}\psi ^{}(w)+\mathrm{}.$$
(4)
If we move $`z`$ once around $`w`$ we pick up some phase factor $`\mathrm{exp}(2\pi iQ_g(j))`$ which is given by the monodromy charge defined above. Let us finally note that the map $`\mathrm{exp}(2\pi iQ_.(j)):\mathrm{\Gamma }S^1`$ defined by $`g\mathrm{exp}(2\pi iQ_g(j))`$ gives rise to a 1-dimensional representation of the group $`\mathrm{\Gamma }`$.
In the case that the simple currents have integer conformal weight, we can include them into the chiral algebra of the theory to form an extended chiral algebra $`𝒲(\mathrm{\Gamma })`$. Note that in this case the monodromy charge $`Q_g(j)`$ depends only on the equivalence class $`[j]=[j]_\mathrm{\Gamma }`$ of $`jJ`$ in the space $`𝒥/\mathrm{\Gamma }`$ of orbits. An orbit $`[j]`$ is said to be invariant, if $`Q_g([j])=Q_g(j)=0`$ for all $`g\mathrm{\Gamma }`$. We can now write down the partition function for the new orbifold theory. Under certain conditions that are satisfied for all cases to be discussed, conformal families of the extended algebra $`𝒲(\mathrm{\Gamma })`$ are labeled by pairs $`([j]_\mathrm{\Gamma },\tau )`$ of invariant orbits along with irreducible representations $`\tau `$ of the associated stabilizer subgroup, i.e. by $`[j]_\mathrm{\Gamma }`$ with $`Q_\mathrm{\Gamma }([j])=0`$ and $`\tau :𝒮_jU(1)`$. The new orbifold theory possesses a diagonal modular invariant partition function with respect to the extended algebra $`W(\mathrm{\Gamma })`$, i.e.
$$Z^{\text{orb}}=\underset{[j],Q_\mathrm{\Gamma }([j])=0}{}|𝒮_{[j]}||\underset{g\mathrm{\Gamma }/𝒮_j}{}\chi _{gj}|^2$$
(5)
Characters appearing with a multiplicity given by the order of the stabilizer are associated with inequivalent representations of the larger algebra $`𝒲(\mathrm{\Gamma })`$. Upon restriction to the original chiral algebra $`𝒲`$, these representations become equivalent. The new partition function is obviously non-diagonal with respect to the smaller algebra $`𝒲`$.
### 2.2 D-branes in the untwisted sector
In this section we shall explain how to construct D-branes in the untwisted sector of simple current orbifolds. Starting from D-branes in the original theory we can obtain certain boundary states for the orbifold theory by a projection onto invariant sectors. This projection method was invented in to get GSO invariant boundary states in Gepner models. The GSO projection can be understood through simple current extension or orbifolds. The methods of (for A-type boundary states) can easily be extended to open descendants of other simple current modular invariants. The idea is to construct boundary states of the theory with a non-diagonal orbifold partition function from boundary states of the original diagonal theory.
So let us suppose we are given our original theory with a partition function
$$Z=\underset{j𝒥}{}|\chi _j|^2.$$
We are looking for associated theories on the half plane such that chiral fields obey the gluing condition $`W(z)=(\mathrm{\Omega }\overline{W})(\overline{z})`$ all along the boundary $`z=\overline{z}`$ and for all chiral fields $`W`$ of the theory. For one special choice of the gluing automorphism $`\mathrm{\Omega }`$ one can apply Cardy’s prescription to find a complete set of elementary boundary states of the original theory (see e.g. for more details). These are given by the formula
$$|I=\underset{j}{}\frac{S_{Ij}}{\sqrt{S_{0j}}}|j.$$
(6)
Here, the element $`|j`$ is the generalized coherent (‘Ishibashi’-) state associated with the sector $`_j\overline{}_j`$ of the theory. $`S`$ is the modular S-matrix of the rational theory under consideration.
Let us now consider the orbifold of the original theory obtained from the action $`g\mathrm{exp}(2\pi iQ_g)`$ of our abelian group $`\mathrm{\Gamma }`$ on the state space $``$. Here we think of the monodromy charge $`Q_g`$ as an operator on the space $`=_j_j\overline{}_j`$ that acts by multiplication with $`Q_g(j)`$ upon restriction to the subspaces $`_j\overline{}_j`$. As usual, the untwisted sector of the orbifold theory contains the states of the original theory that are invariant under the action of $`\mathrm{\Gamma }`$. We can map the states from $``$ to the untwisted sector of the orbifold theory with the help of the projector
$$P^0=\frac{1}{|\mathrm{\Gamma }|}\underset{g\mathrm{\Gamma }}{}\mathrm{exp}(2\pi iQ_g).$$
The idea in was to take the (appropriately rescaled) image of the Cardy boundary states (6) under $`P^0`$ as boundary states of the orbifold theory, i.e.
$$|I_{\text{proj}}=\kappa P^0|I=\frac{\kappa }{|\mathrm{\Gamma }|}\underset{g\mathrm{\Gamma }}{}e^{2\pi iQ_g}\underset{j}{}\frac{S_{Ij}}{\sqrt{S_{0j}}}|j.$$
(7)
By construction, these states are linear combinations of generalized coherent states in the orbifold theory. Hence, it only remains to check that they give rise to consistent open string spectra. To this end one may exploit the following property of modular S-matrices
$$S_{gIj}=e^{2\pi iQ_g(j)}S_{Ij}.$$
(8)
Insertion into our previous expression for the projected boundary states $`|I_{\text{proj}}`$ yields the interesting result
$$|I_{\text{proj}}=\frac{\kappa }{|\mathrm{\Gamma }|}\underset{g\mathrm{\Gamma }}{}|gI.$$
(9)
This means that in a general model, the charge projected boundary states are orbits of boundary states in the theory before orbifolding. The formula should be regarded as an algebraic analogue of the geometric prescription to sum over pre-images of a brane in the covering space. Hence, the open string spectra associated with these the boundary states $`|I_{\text{proj}}`$ are
$`Z_{IJ}^{\text{proj}}(q)`$ $`:=`$ $`{}_{\text{proj}}{}^{}\theta I|\stackrel{~}{q}^{\frac{1}{2}(L_0+\overline{L}_0)\frac{c}{24}}|J_{\text{proj}}^{}`$ (10)
$`=`$ $`{\displaystyle \frac{\kappa ^2}{|\mathrm{\Gamma }|^2}}{\displaystyle \underset{g,h\mathrm{\Gamma },k}{}}N_k^{gI^{}hJ}\chi _k(q)={\displaystyle \frac{\kappa ^2}{|\mathrm{\Gamma }|}}{\displaystyle \underset{g,k}{}}N_k^{gI^{}J}\chi _k(q).`$
Here, $`\theta `$ denotes the CPT operator of the bulk theory and $`I^{}`$ is the label conjugate to $`I`$ (see for details). If we choose
$$\kappa =\sqrt{|\mathrm{\Gamma }|}$$
(11)
then the coefficients of the characters on the right hand side are guaranteed to be integer. But they may still possess a common divisor. This happens whenever $`I`$ or $`J`$ is fixed under the action of some element $`g\mathrm{\Gamma }`$. If one of the labels is on such a short orbit, then the sum in (9) consists of several equal terms. For $`I=J`$, the number of equal terms is given by the order $`|𝒮_I|`$ of the stabilizer subgroup $`𝒮_I\mathrm{\Gamma }`$ of $`I`$. We shall see below that in the orbifold theory there exist $`|𝒮_I|`$ boundary states associated with the label $`I`$ which differ by a representation of the stabilizer subgroup $`𝒮_I`$. The projection analysis provides one linear combination of these $`|𝒮_I|`$ states which contains only contributions from the untwisted sector.
## 3 Gepner models and untwisted D-branes
The aim now is to apply the general theory outlined above to an important class of examples, namely to Gepner models (see also for a review). These are exactly solvable CFTs which are used to study strings moving on a Calabi-Yau manifold at small radius . Their construction employs an orbifold-like projection starting from a tensor products of $`r`$ $`𝒩=2`$ minimal models. In our presentation we shall assume that there are $`d=1`$ complex, transverse, external dimensions in light cone gauge and that the number $`r`$ of minimal models equals $`r=5`$.
### 3.1 The building blocks of Gepner models
Our main building blocks are $`𝒩=2`$ minimal models at level $`k`$. These are SCFTs with central charge $`c=\frac{3k}{k+2}<3`$ . One can label primaries of the bosonic subalgebra by 3 integers, $`(l,m,s)`$ taking values in the range
$$0lk;0|ms|l;s\{1,0,1\};l+m+s=0\mathrm{mod}2.$$
(12)
In many respects, $`l`$ and $`m`$ behave like the familiar labels in an $`SU(2)_k`$ WZW model. The third label $`s`$ determines the spin structure. States with $`s=0`$ are in the NS sector while $`s=\pm 1`$ correspond to the two chiralities in the $`R`$ sector. The conformal weights and $`U(1)`$ charges of these primary fields can be computed by means of the formulas
$`h_{m,s}^l`$ $`=`$ $`{\displaystyle \frac{l(l+2)m^2}{4(k+2)}}+{\displaystyle \frac{s^2}{8}}\text{mod}1,`$ (13)
$`q_{m,s}^l`$ $`=`$ $`{\displaystyle \frac{m}{k+2}}{\displaystyle \frac{s}{2}}\text{mod}2.`$ (14)
There are some distinguished labels $`(l,\pm l,0)𝒥`$ in the NS sector which are associated with the $`𝒩=2`$ (anti-) chiral primaries. They play a special role in the Landau-Ginzburg description of minimal models where they are identified with powers $`X^l(l,l,0)`$ of the Landau-Ginzburg field $`X`$.
Our set $`𝒥`$ of conformal families contains triples $`(l,m,s)`$ from the standard range. It will be rather convenient below to consider an extended set $`\stackrel{~}{𝒥}`$ of labels $`(l,m,s)\{0,\mathrm{},k\}\times _{2k+4}\times _4`$ with the only additional constraint $`l+m+s=`$ even. For each label $`(l,m,s)𝒥`$ there exist exactly two labels in $`\stackrel{~}{𝒥}`$, namely $`(l,m,s)`$ and $`(kl,m+k+2,s+2)`$. In other words, the extended set $`\stackrel{~}{𝒥}`$ carries an action of $`_2`$ that maps $`(l,m,s)`$ to $`(kl,m+k+2,s+2)`$ such that our original set $`𝒥`$ is simply the quotient $`\stackrel{~}{𝒥}/_2`$. Passing from $`\stackrel{~}{𝒥}`$ to $`𝒥`$ is known as field identification.
The simple currents of an $`𝒩=2`$ minimal model can be determined from the fusion rules. For the label $`l`$ these are given by the usual $`SU(2)`$ fusion rules while both other labels add like representations of the abelian groups $`_{2k+4}`$ and $`_4`$, respectively. This implies that e.g. $`(0,1,1)`$ and $`(0,0,2)`$ are both simple currents. They are of special interest in the context of Gepner models and will be used to generate our simple current group $`\mathrm{\Gamma }`$. $`(0,1,1)`$ is the spectral flow by $`1/2`$ unit and $`(0,0,2)`$ the world-sheet supersymmetry generator. $`(0,0,2)`$ is a simple current of order $`2`$ and can be used to group the world-sheet fields into supermultiplets. The order of the simple current $`(0,1,1)`$ is model dependent. To see this, we apply the current $`2k+4`$ times to the identity. This will lead us back to the identity whenever the level $`k`$ is even. Since $`(0,0,2)`$ is nowhere on this orbit, $`(0,0,2)`$ and $`(0,1,1)`$ togther generate the simple current group $`\mathrm{\Gamma }=\mathrm{\Gamma }_k=_{2k+4}\times _2`$ for even level $`k`$. When $`k`$ is odd, however, we reach the field $`(0,0,2)`$ after $`2k+4`$ applications of $`(0,1,1)`$ and hence we have to apply the simple current $`(0,1,1)`$ another $`2k+4`$ times. In this case, the orbit contains the label $`(0,0,2)`$ and hence our orbifold group is $`\mathrm{\Gamma }=\mathrm{\Gamma }_k=_{4k+8}`$ for odd $`k`$.
Let us briefly describe the orbits for the action of $`\mathrm{\Gamma }_k`$ on the set $`𝒥`$. Again, we need to treat the cases of even and odd $`k`$ separately. If $`k`$ is odd, the orbifold group $`\mathrm{\Gamma }_k`$ acts freely so that all orbits have length $`4k+8`$. For even level $`k`$, however, we generate short orbits of length $`2k+4`$ whenever we start from a field $`(l,m,s)`$ with $`l=k/2`$ because the label $`l=k/2`$ is invariant under field identification. The stabilizer for these short orbits is a subgroup $`_2\mathrm{\Gamma }_k`$.
From formula (13) and the definition relation (3) we can compute the monodromy charge:
$$Q_{(0,\mu ,\sigma )}^k(l,m,s):=\frac{m\mu }{2k+4}\frac{s\sigma }{4}\text{mod}1\text{ for all }(0,\mu ,\sigma )\mathrm{\Gamma }_k.$$
Obviously, this expression defines a representation of our orbifold group $`\mathrm{\Gamma }_k`$ for any choice of $`(l,m,s)𝒥`$.
$`𝒩=2`$ characters and their modular properties are described e.g. in . The characters $`\chi _{l,m,s}`$ of the representations listed above can be indexed by arbitrary labels $`(l,m,s)`$ in the extended set $`\stackrel{~}{𝒥}`$. In fact, we can use the field identification to define $`\chi _{l,m,s}`$ for $`(l,m,s)`$ outside the standard range (12) by $`\chi _{m,s}^l=\chi _{m+k+2,s+2}^{kl}`$. For the modular S-matrix one has the explicit formula
$$S_{(l,m,s),(l^{},m^{},s^{})}^k=\frac{1}{\sqrt{2}(k+2)}\mathrm{sin}\pi \frac{(l+1)(l^{}+1)}{k+2}e^{i\pi mm^{}/(k+2)}e^{i\pi ss^{}/2}$$
(15)
It is easy to check that $`S`$ satisfies the relation (8) for all $`g\mathrm{\Gamma }_k`$.
In addition to the minimal models, Gepner’s constructions involve fermions from the external space-time sector of the theory. Their bosonic subalgebra is an $`SO(2)_1`$ current algebra. Its representations are labeled by $`s=0,\pm 1,2_4`$ and they possess the obvious abelian fusion rules. All sectors in this theory are simple currents and their monodromy charge is given by
$$Q_\sigma ^𝖿(s):=\frac{s\sigma }{4}\text{mod}1\text{ for }\sigma =0,\pm 1,2.$$
The property (8) is obeyed by the modular S-matrix $`S_{s,s^{}}^𝖿=(1/2)\mathrm{exp}(i\pi ss^{}/2)`$ of the $`SO(2)_1`$ current algebra.
### 3.2 Gepner models in the bulk
Let us now consider a tensor product of $`r=5`$ minimal models with levels $`k_i,i=1,\mathrm{},r`$ such that their central charges adds up to $`c=9`$. To get a full string theory, one needs to add a sector containing ghosts and a level $`k=1`$ current algebra that comes with the space-time sector. These tensor products do not give consistent string backgrounds with 4d spacetime SUSY. But there exists some orbifold theory of this tensor product theory that satisfies all the requirements of a consistent string background. For its description we need further notations. Let us introduce the following vectors
$$\lambda :=(l_1,\mathrm{},l_r)\mathrm{and}\mu :=(s_0;m_1,\mathrm{},m_r;s_1,\mathrm{},s_r)$$
to label the tensor product of representations $`(l_j,m_j,s_j)`$ of the individual minimal models and of the representations $`s_0=0,2,\pm 1`$ that come with the level $`k=1`$ current algebra. The associated product of characters $`\chi _{m_i,s_i}^{l_i}`$ and $`\chi _{s_0}`$ is denoted by $`\chi _\mu ^\lambda (q)`$.
Next, we introduce the special $`(2r+1)`$-dimensional vectors $`\beta _0`$ with all entries equal to 1, and $`\beta _j`$, $`j=1,\mathrm{},r`$, having zeroes everywhere except for the 1st and the $`(r+1+j)`$th entry which are equal to 2. These vectors stand for particular elements in the group $`_4\times _i\mathrm{\Gamma }_{k_i}`$. It is easily seen that they generate a subgroup $`\mathrm{\Gamma }=_K\times _2^r`$ where $`K:=\mathrm{lcm}(2k_j+4)`$. Elements of this subgroup will be denoted by $`𝝂=(\nu ,\nu _1,\mathrm{},\nu _r)`$. The monodromy charge of a pair $`(\lambda ,\mu )`$ is
$`Q_𝝂(\lambda ,\mu )`$ $`=`$ $`\nu \beta _0\mu +{\displaystyle \underset{i=1}{\overset{r}{}}}\nu _i\beta _i\mu \text{mod}1`$ (16)
$`\text{where}\beta _0\mu `$ $`:=`$ $`{\displaystyle \frac{s_0}{4}}{\displaystyle \underset{j=1}{\overset{r}{}}}{\displaystyle \frac{s_j}{4}}+{\displaystyle \underset{j=1}{\overset{r}{}}}{\displaystyle \frac{m_j}{2k_j+4}},`$ (17)
$`\beta _j\mu `$ $`:=`$ $`{\displaystyle \frac{s_0}{2}}{\displaystyle \frac{s_j}{2}}.`$ (18)
The orbifold group $`\mathrm{\Gamma }`$ acts on the labels $`\lambda `$ and $`\mu `$ in the obvious way. There appear orbits of maximal length $`K2^r`$ and short orbits of length $`K2^{r1}`$. The latter are characterized by the property that $`\lambda =(l_1,\mathrm{},l_r)`$ satisfy $`l_i=k_i/2`$ for all $`i`$ such that $`2k_i+4`$ is not a factor in $`K/2`$.
As we mentioned before, the complete construction requires to include ghosts. Since the ghost sector will not play an important role in the later sections, we shall constrain ourselves to some brief remarks that are necessary in understanding Gepner models from an orbifold point of view.
When we consider the full theory, the field generating the $`_K`$-symmetry contains a factor from the ghost sector and the space-time part of the spin field $`S^\alpha `$. In the $`\pm 1/2`$ picture, the operator can be represented as
$$U(z)=e^{\pm i\frac{\varphi }{2}}e^{\pm i\frac{1}{2}\sqrt{\frac{c}{3}}X}S_\alpha $$
(19)
Here, $`\varphi `$ denotes the bosonized super-ghost and we introduced the bosonic field $`X`$ whose derivative gives the U(1) current $`J(z)=_{i=0}^rJ_i(z)`$ in the tensor product of the $`SO(2)_1`$ theory with the minimal models. The precise relation is $`J(z)=i\sqrt{\frac{c}{3}}X`$.
The operator (19) is a simple current and its internal part agrees with the simple current $`\beta _0`$ in the tensor product considered above. Since the operator (19) has total weight one, it can be added to the chiral algebra and we can use the formula (5) to determine the partition function of the orbifold theory. The formula requires to determine invariant orbits, i.e. orbits with vanishing monodromy charge. Taking the OPE of the spectral flow (19) with a vertex operator that represent space-time scalars of the theory gives the monodromy charge
$$\stackrel{~}{Q}_𝝂(\lambda ,\mu )=\nu \left(\frac{\beta _0\mu }{2}+\frac{1}{2}\right)+\underset{i=1}{\overset{r}{}}\nu _i\beta _i\mu \text{mod}1$$
by comparison with the general formula (4). $`\stackrel{~}{Q}`$ effectively replaces the monodromy charge $`Q`$ introduced in eq. (16). The orbits of vanishing monodromy charge are those of odd integer $`U(1)`$ charge.
To write a partition function of physical states one has to extract the physical degrees of freedom. This can be done by a projection onto light-cone variables which removes, in particular, the ghost sector. In practice, the light-cone degrees of freedom may be read off directly in the canonical ghost pictures.
An important reason to include ghosts is that the fields of the theory aquire the right commutation properties. One would like to incorporate this feature in the physical theory. Therefore, in the partition function, the fields are counted with a ghost-charge dependent phase factor $`\mathrm{exp}(2\pi iq_{ghost})`$. This means that the states with half-integer ghost charge, i.e. the RR-states, contribute with a negative sign.
Having discussed all aspects of the ghosts, which are relevant to the simple current construction of the Gepner partition function, we will discard them from now on and consider only the physical (light-cone) degrees of freedom. We are prepared to write down the partition function for a Gepner model describing a superstring compactification to $`4`$ dimensions. It is given by
$$Z_G^{(r)}(\tau ,\overline{\tau })=\frac{1}{2}\frac{(\mathrm{Im}\tau )^2}{|\eta (q)|^2}\underset{\lambda ,\mu ;\stackrel{~}{Q}(\lambda ,\mu )=0}{}\underset{\nu ,\nu _j}{}(1)^\nu \chi _\mu ^\lambda (q)\chi _{\mu +\nu \beta _0+\nu _1\beta _1+\mathrm{}+\nu _r\beta _r}^\lambda (\overline{q})$$
The sign is the usual one occurring in (space-time) fermion one-loop diagrams. The $`\tau `$-dependent factor in front of the sum accounts for the free bosons associated to the $`2`$ physical transversal dimensions of flat external space-time, while the $`1/2`$ is simply due to the field identification mentioned above. Except for these modifications, the formula for $`Z_G`$ is the same as eq. (5). Elements $`g=\nu \beta _0+\mathrm{}\nu _r\beta _r`$ of the orbifold group $`\mathrm{\Gamma }`$ are labeled by $`\nu ,\nu _i`$ so that the second sum is over the full group $`\mathrm{\Gamma }`$. Short orbits appear twice in the summation and give rise to an extra factor of $`2`$ which is the order of the corresponding stabilizer subgroup. Since our orbifold group $`\mathrm{\Gamma }`$ is abelian, we used additive notation for the action of elements $`g\mathrm{\Gamma }`$ on the labels $`\lambda ,\mu `$.
### 3.3 Untwisted D-branes in Gepner models
Using the general formalism outlined in Section 2.2 one can find a large set of boundary states which respect the $`𝒩=2`$ world-sheet algebras of each minimal model factor of the Gepner model separately. To this end we start with Cardy boundary states of the tensor product theory. They are given by the expression (6) along with the formula (15) for the modular S-matrices of minimal models and the simple expression for the S-matrix of $`SO(2)_1`$ that we spelled out before. The generalized coherent states $`|j`$ are now parametrized by pairs $`j=(\lambda ,\mu )`$. Cardy’s boundary states belong to some gluing condition $`W(z)=\mathrm{\Omega }\overline{W}(\overline{z}),z=\overline{z},`$ which becomes $`J_i(z)=\overline{J}_i(\overline{z})`$ on the U(1)-currents of the individual theories. This means that they are A-type boundary conditions in the sense of .
The boundary states $`|I=:|\mathrm{\Lambda },\mathrm{\Xi }`$ we have just described depend on a spin vector $`\mathrm{\Lambda }=(L_1,\mathrm{},L_r)`$ and a charge vector $`\mathrm{\Xi }=(S_0;M_1,\mathrm{},M_r;S_1,\mathrm{},S_r)`$. From these states in the tensor product theory we can pass to boundary states of the Gepner model using the general strategy explained in Section 2.2. The projected boundary states in the orbifold theory are given by (see eqs. (9,11))
$$|\mathrm{\Lambda },\mathrm{\Xi }_{\text{proj}}=\frac{1}{\sqrt{K2^r}}\underset{\nu ,\nu _i}{}(1)^\nu (1)^{\frac{\widehat{s}_0^2}{2}}|\mathrm{\Lambda },\mathrm{\Xi }+\nu \beta _0+\nu _1\beta _1+\mathrm{}+\nu _r\beta _r.$$
Here, the element $`\widehat{s}_0`$ is an operator acting on closed string states which measures the value $`s_0`$. The whole factor $`(1)^{\widehat{s}_0^2/2}`$ is needed to guarantee that in the open string partition function (similar to the closed string partition function) fields are counted with a phase factor referring to their ghost charge. If we insert the formula (6) with the appropriate modular S-matrix on the right hand side, we obtain the expressions established in ,
$$|\alpha :=|\mathrm{\Lambda },\mathrm{\Xi }_{\text{proj}}=\underset{\lambda ,\mu ;\stackrel{~}{Q}(\lambda ,\mu )=0}{}(1)^{\frac{s_0^2}{2}}B_\alpha ^{\lambda ,\mu }|\lambda ,\mu .$$
(20)
with the coefficients:
$$B_\alpha ^{\lambda ,\mu }=\frac{\sqrt{K2^r}}{2}e^{i\pi \frac{s_0S_0}{2}}\underset{j=1}{\overset{r}{}}\frac{1}{\sqrt{\sqrt{2}(k_j+2)}}\frac{\mathrm{sin}(l_j,L_j)_{k_j}}{\sqrt{\mathrm{sin}(l_j,0)_{k_j}}}e^{i\pi \frac{m_jM_j}{k_j+2}}e^{i\pi \frac{s_jS_j}{2}}.$$
(21)
Here $`(l,l^{})_k=\pi (l+1)(l^{}+1)/(k+2)`$. For these A-type boundary states the Ishibashi states are built on diagonal primary states, i.e. states in the untwisted sector, in accordance with our general theory in Section 2. The associated partition functions (10) aquire the following form (see also ):
$`Z_{\stackrel{~}{\alpha }\alpha }^A(q)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\lambda ^{},\mu ^{}}{}}{\displaystyle \underset{\nu =0}{\overset{K1}{}}}{\displaystyle \underset{\nu _i=0,1}{}}(1)^{s_0^{}+S_0\stackrel{~}{S}_0}\delta _{s_0^{}\stackrel{~}{S}_0+S_0+\nu +2{\scriptscriptstyle \nu _i}2}^{(4)}`$ (22)
$`\times {\displaystyle \underset{j=1}{\overset{r}{}}}N_{L_j,\stackrel{~}{L}_j}^{l_j^{}}\delta _{\nu +M_j\stackrel{~}{M}_j+m_j^{}}^{(2k_j+4)}\delta _{s_j^{}\stackrel{~}{S}_j+S_j+\nu +2\nu _j}^{(4)}\chi _\mu ^{}^\lambda ^{}(q).`$
The factor $`1/2`$ in front of the right hand side accounts for the fact that field identification causes each character to appear twice when we sum over $`\lambda ^{},\mu ^{}`$ taken from the extended range.
If the two boundary conditions $`\alpha ,\stackrel{~}{\alpha }`$ appearing in eq. (22) are both labeled by monodromy invariant orbits, they give rise to a monodromy invariant open string spectrum, i.e. to a spectrum that contains only odd-integer charges. One should recall, however, that non-invariant orbits of $`\mathrm{\Gamma }`$ are also admissible as labels for boundary conditions. The condition for a supersymmetric open string spectrum consisting of monodromy invariant states is that the U(1) charge of the two orbits labeling the boundary conditions $`\alpha `$ and $`\stackrel{~}{a}`$ differs by an even integer.
## 4 Singular curves on Calabi-Yau 3-folds
The boundary states (20) we have constructed so far are only charged under $`(c,c)`$-fields in the untwisted sector. Typically, there exist additional $`(c,c)`$-fields in the twisted sectors of Gepner models. Their appearance is related to singularities of the associated Calabi-Yau spaces. The aim of this section is to explain this relation in some more detail. Once this is understood, it motivates the search for additional boundary states that are charged under fields in the twisted sectors. This will be addressed in the next section. Our presentation here will be rather sketchy. The interested reader can find more explanations and details e.g. in .
Suppose we are given a Gepner model composed from $`r=5`$ minimal models with levels $`(k_1,\mathrm{},k_5)`$ chosen such that their central charges add up to $`c=9`$. Let us define integers $`\omega _i=K/(2k_i+4)`$ where $`K=\mathrm{lcm}(2k_i+4)`$ as before. Assuming that the Gepner model has an A-type modular invariant partition function, as we did before, we associate with it the following space $`M=P_{(\omega _1,\mathrm{},\omega _5)}[\frac{K}{2}]`$ which is defined by the equation
$$M:z_1^{k_1+2}+\mathrm{}z_5^{k_5+2}=0$$
evaluated in a weighted projective space where $`(z_1,\mathrm{},z_5)(\lambda ^{\omega _1}z_1,\mathrm{},\lambda ^{\omega _5}z_5)`$ for $`\lambda ^{}`$. It is well known that the $`(c,c)`$-fields with total left- and right-moving charges $`q`$, $`\overline{q}`$ satisfying $`q=1=\overline{q}`$ correspond to harmonic (2,1)-forms on $`M`$. We denote the space of these forms by $`H^{2,1}`$.
Typically, the space $`M`$ possesses singularities which can be either singular points or singular curves. In string theory these singularities are resolved. Desingularization of singular curves provides a non-vanishing contribution to the Hodge number $`h^{21}=\mathrm{dim}H^{2,1}`$. We will make a more quantitative statement momentarily after a brief description of the possible singularities.
A singular curve on $`M`$ exists whenever three of the numbers $`\omega _i`$ have some non-trivial factor in common. To be specific, we suppose that these are $`\omega _3,\omega _4,\omega _5`$ and we denote their largest common divisor by $`N`$. Now our weighted projective space carries an action of $`_N`$ defined by $`(z_1,\mathrm{},z_5)(\eta ^{\omega _1}z_1,\mathrm{},\eta ^{\omega _5}z_5)`$ where $`\eta `$ is some $`N^{th}`$ root of unity. By our choice of $`N`$, $`\eta ^{\omega _i}=1`$ for $`i=3,4,5`$ and hence, there is a two complex dimensional subspace in the weighted projective space that remains fixed under the action of $`_N`$. This subspace is parametrized by classes of $`(0,0,z_3,z_4,z_5)`$. The surface $`M`$ intersects with this singularity along the curve given by the equation
$$C:z_3^{k_3+2}+z_4^{k_4+2}+z_5^{k_5+2}=0.$$
Putting all this together we see that $`M`$ possesses a $`_N`$ singularity over each point of the curve $`C`$. At a generic point, the space transverse to the singular curve looks like the quotient $`^2/_N`$ where $`_N`$ acts on points $`(w_1,w_2)^2`$ by $`(w_1,w_2)(\eta w_1,\eta ^1w_2)`$.
The singularity at the origin of $`^2/_N`$ is known as a rational double point of type $`A_{N1}`$. It is resolved by gluing in a chain of $`N1`$ projective spaces $`\mathrm{IP}^1`$ which intersect pairwise transversely in one point and have self-intersection numbers $`2`$. In other words, the intersection matrix equals the negative Cartan matrix for an $`A_{N1}`$ Dynkin diagram. When we resolve the $`_N`$ singularity along the whole fixed curve $`C`$, we obtain locally a product of the curve $`C`$ and the chain of spheres we have just described. The resolution of a single fixed curve $`C`$ shifts the Hodge number $`h^{21}`$ by $`(N1)g`$ with $`g`$ being the genus of the curve $`C`$ . This can be understood by sending a one-cycle of the curve $`C`$ to the three-cycle swept out by the spheres fibered over the one-cycle . (See also for a discussion of singular curves in the context of gauge symmetry enhancement in type II theories.)
Let us turn back to the Gepner models and compare what we have learned about the resolution of singularities with the appearance of $`(c,c)`$-fields in the twisted sectors. As we have pointed out before, the Landau-Ginzburg description of minimal models involves an identification of the coordinate $`z_i`$ with the chiral primary field $`\mathrm{\Phi }_{m,s}^{i;l}=\mathrm{\Phi }_{1,0}^{i;1}`$ in the $`i^{th}`$ minimal model. The subgroup $`_K\mathrm{\Gamma }`$ acts on the latter by multiplication with the phase $`\mathrm{exp}(2\pi i/(2k_i+4))`$. Note that the number $`N`$ must be contained as a factor in $`K`$ since each of the numbers $`\omega _3,\omega _4,\omega _5`$ is contained in $`K`$. Hence, we obtain an action of $`_N_K`$ by
$$\mathrm{\Phi }_{1,0}^{i;1}e^{\frac{2\pi i}{2k_i+4}\frac{K}{N}}\mathrm{\Phi }_{1,0}^{i;1}=e^{\frac{2\pi i}{N}\omega _i}\mathrm{\Phi }_{1,0}^{i;1}=\eta ^{\omega _i}\mathrm{\Phi }_{1,0}^{i;1},$$
where $`\eta =\mathrm{exp}(2\pi i/N)`$. This is precisely the transformation law of the coordinate $`z_i`$ under the action of the $`_N`$ that is responsible for the singularity in the weighted projective space. Hence, we may identify the subgroup $`_N_K\mathrm{\Gamma }`$ of the orbifold group with the geometrically acting group on our Calabi-Yau manifold $`M`$.
Combining the last paragraph with the previous discussion we can now formulate the relation between the singularities and $`(c,c)`$-fields more precisely . In fact, we expect to find $`g`$ of the $`(c,c)`$-fields with $`q=1=\overline{q}`$ in each of the $`\nu =0`$ mod $`K/N`$-twisted sectors. As there exist $`N1`$ such sectors (the untwisted sector $`\nu =0`$ is not included), the number of new $`(c,c)`$-fields matches the shift of the Hodge number that we attributed to the desingularization of a single $`_N`$-singular curve $`C`$.
While these new $`(c,c)`$-fields possess the same total right- and left moving U(1)-charge, i.e. $`q=\overline{q}`$ by construction, the U(1)-charges in the individual minimal models can be different. This may be seen from the obvious property $`q^{(i)}\overline{q}^{(i)}=\nu /(k_i+2)`$ mod $`1`$ of fields in the $`\nu `$-twisted sector (it holds for $`\nu _i=0`$). The right hand side of these equations is non-zero unless $`2\nu `$ can be devided by all the $`2k_i+4`$. But this requires $`2\nu =K`$ and hence we conclude that only the new $`(c,c)`$-fields in the $`\nu =K/2`$-twisted sector have equal left- and right-moving charges $`q^{(i)}`$ = $`\overline{q}^{(i)}`$ in all the individual minimal models. It is easy to see that all such fields lie on short orbits of our orbifold group $`\mathrm{\Gamma }`$.
The boundary states we discuss in this work satisfy A-type boundary conditions in each minimal model rather than A-type boundary conditions for the diagonally embedded superconformal algebra only. This means that generalized coherent states in our boundary states are necessarily based on primary fields satisfying the condition $`q^{(i)}`$ =$`\overline{q}^{(i)}`$. Hence, all we can hope for in the following is to find boundary states which are charged under the new $`(c,c)`$-fields in the $`\nu =K/2`$-twisted sector, i.e. under the $`(c,c)`$-fields that arise from resolving a $`_2`$-singularity along a fixed curve $`C`$.
Example: The degree eight Fermat CY space $`\mathrm{IP}_{1,1,2,2,2}[8]`$ corresponds to the choice $`(k=6)^2(k=2)^3`$ for the minimal models used in Gepner’s construction. The surface has only one $`_2`$ singularity over a curve $`C`$ of genus $`g=3`$. Resolving this singularity shifts the Hodge number $`h^{21}`$ by 3. It is not difficult to list all the $`(c,c)`$-fields with $`q=1=\overline{q}`$ in this case. Except from those in the untwisted sector, there appear 3 additional such fields in the $`\nu =8`$-twisted sector with labels for their left-movers being given by
$$(3,3,0)\times (3,3,0)\times (1,1,0)\times (0,0,0)\times (0,0,0)$$
and by similar expressions obtained through a permutation of the last three factors. Obviously, these fields are all on short orbits of the group $`\mathrm{\Gamma }=_{16}\times _2^5`$.
## 5 Twisted D-branes in Gepner models
We have argued in the previous section that a typical Gepner model has $`(c,c)`$-fields in the twisted sectors. Geometrically, this situation corresponds to the resolution of singularities inherited from the embedding projective space. D-branes can wrap the extra homology cycles on the resolved manifold. Therefore, we expect that there are boundary states containing Ishibashi states built on the $`(c,c)`$ fields in twisted sectors of the models. This is necessary in order to have D-branes which are charged under all the possible massless RR potentials in the theory. Since locally the singularities we shall be concerned with look like $`_2`$ orbifolds of flat space, we will first review the construction of boundary states for fractional branes in flat space. This will help to motivate the procedure that allows to find twisted sector boundary states in Gepner models.
### 5.1 Twisted boundary states in flat space
D0 branes at orbifold singularities of flat space have been discussed from the open string point of view in . There are different types of branes depending on the representation of the orbifold group on the Chan-Paton factors. The group $`_2`$ has two inequivalent irreducible one-dimensional representations which are distinguished by a sign. These two representations lead to two types of fractional branes. In fact, to each of the two representations one can associate a boundary state (see also for a related construction in a different context) for branes sitting at the origin $`z=(z_1,z_2)=(0,0)`$. They are given by
$$|D(z=0);\pm =\frac{1}{\sqrt{2}}|D(z=0)\pm \sqrt{2}|D(z=0)^{\mathrm{tw}}.$$
Here, $`|D(z=0)^{\mathrm{tw}}`$ is an Ishibashi state obtained from the twisted sector of the $`_2`$-orbifold theory. After modular transformation, the transition amplitude between two such twisted boundary states acquires the form
$`Z_{++}=Z_{}(q)`$ $`=`$ $`Z_{|D0}(q)+Z^{\mathrm{tw}}(q)=\mathrm{tr}\left({\displaystyle \frac{1+g}{2}}q^{H_{op}}\right),`$
$`Z_+=Z_+(q)`$ $`=`$ $`Z_{|D0}(q)Z^{\mathrm{tw}}(q)=\mathrm{tr}\left({\displaystyle \frac{1g}{2}}q^{H_{op}}\right),`$
where $`g`$ is the non-trivial element of $`_2`$ acting on both Chan-Paton labels and oscillators. The twisted part $`|D(z=0)^{\mathrm{tw}}`$ of the boundary state is responsible for the term $`Z^{\mathrm{tw}}(q)`$ in the open string partition function in which the group element $`g`$ is inserted.
If we add the two boundary states $`|D;+`$ and $`|D;`$ we obtain a boundary state without contributions from twisted sectors. This state comes with the regular representation of the orbifold group on Chan-Paton labels and it can be obtained directly by the projection method applied to the boundary state $`|D(x=0)`$ of the theory on the two-fold cover of the orbifold space. Turning the argument around, we see that boundary states at fixed points obtained with the projection method can be further decomposed by taking into account twisted sectors. This is the strategy we will now follow in our discussion of twisted boundary states for Gepner models.
### 5.2 Twisted boundary states in the Gepner model
In this section we will construct additional boundary states for the Gepner model, which are charged under the twisted sector states in short orbits of the orbifold group $`\mathrm{\Gamma }`$. The method will be similar to the one used in flat space. The role of the state $`|D(x=0)^{\mathrm{tw}}`$ is played Ishibashi states built from states in $`K/2`$-twisted sector of the Gepner models. Our aim is to improve the boundary states associated with short orbits of $`\mathrm{\Gamma }`$ by adding twisted sector Ishibashi states.
There are a few general statements one can make about this procedure whenever $`_2`$ appears as the maximal stabilizer subgroup of the theory. In this case, the short orbits have length $`K2^{r1}`$. The application of the projection (9) leads to two equal pieces in the projected sum. Each of the two pieces can be interpreted as a sum over elements of the orbifold group modulo the stabilizer. It means that the boundary states we started with are not elementary but rather have to be replaced by a linear combination of the projected state and Ishibashi states from the $`\nu =K/2`$ twisted sector, i.e.
$$|\mathrm{\Lambda }_s,\mathrm{\Xi }_s_{\text{proj}}\beta |\mathrm{\Lambda }_s,\mathrm{\Xi }_s_{\text{proj}}\pm \underset{\lambda _s,\mu _s}{}\alpha _{\lambda _s,\mu _s}|\lambda _s,\mu _s^{\mathrm{tw}},$$
(23)
where and $`|\lambda _s,\mu _s`$ is an Ishibashi state built on a short orbit primary labeled by $`\lambda _s`$ and $`\mu _s`$. The coefficients $`\alpha `$ and $`\beta `$ in the boundary state have to be chosen in such a way that the boundary state leads to a consistent open string spectrum. We will see that we can easily get information on $`\beta `$ from general considerations, whereas $`\alpha `$ is model dependent.
Let us start the discussion by computing the open string spectrum in the case that one boundary state is given by the improved short orbit boundary state and the other one is given by a projected long orbit boundary state. In this case, we can get conditions on $`\beta `$ only, since the twist fields are orthogonal to all fields of the untwisted sector. The open string sector is determined by the fusion rules of the elements of the short orbit with those of the long orbit (cf. eq. (10) above). According to our general discussion following eq. (10), the overall multiplicity of the fields propagating in the open string sector is $`2\beta `$. Therefore, minimal normalization suggests to pick
$$\beta =\frac{1}{2}.$$
Of course one has to check that this choice gives consistent open string spectra when tested against other short orbit states of the form (23). We will do this in the explicit examples below.
To prepare for the discussion of Gepner models, we determine the improved boundary states in a single $`𝒩=2`$ minimal model, where we mod out by the group generated by the current $`(0,1,1)`$. We pick a minimal model with $`k=0`$ mod $`4`$. In this case, the length of a generic orbit of the current $`(0,1,1)`$ is $`2k+4`$ and there is a short orbit for $`l=k/2`$. As we outlined above, our plan is to add Ishibashi states $`|k/2,m+\nu ,s+\nu ^{\mathrm{tw}}`$ built on twisted sector fields to the short orbit boundary states. Our ansatz for the improved boundary states is
$`|{\displaystyle \frac{k}{2}},M,S`$ $`=`$ $`{\displaystyle \frac{1}{2}}|{\displaystyle \frac{k}{2}},M,S_{\text{proj}}+`$
$`+\stackrel{~}{\alpha }{\displaystyle \frac{1}{2^{1/4}(k+2)}}{\displaystyle \underset{\nu =0}{\overset{k+1}{}}}{\displaystyle \underset{m,s}{\overset{ev}{}}}(1)^\nu e^{i\pi m\frac{M+\nu }{k+2}}e^{i\pi s\frac{S+\nu }{2}}|{\displaystyle \frac{k}{2}},m,s^{\mathrm{tw}},`$
where the sum over $`m`$ and $`s`$ is constrained to $`m+s=even`$. We allow for an overall factor $`\stackrel{~}{\alpha }`$ in front of the twisted part of the boundary state. It will be adjusted later so that we get a consistent open string spectrum. To this end, we have to compute the transition amplitude between the twisted parts of the boundary state and to perform a modular transformation. This results in the following expression for the contribution $`Z^{\mathrm{tw}}`$ of the twisted sector to the full open string partition function :
$$Z^{\mathrm{tw}}=\stackrel{~}{\alpha }^2\frac{2}{k+2}\underset{\nu =0}{\overset{k+1}{}}\underset{l^{},m^{},s^{}}{\overset{ev}{}}(1)^\nu \mathrm{sin}\pi \frac{l^{}+1}{2}\delta _{\nu +m^{}+\frac{k+2}{2}(\nu +s^{})}^{(2k+4)}\delta _{\nu +s^{}}^{(2)}\chi _{l^{},m^{}s^{}}(q).$$
(25)
From our previous discussion on the minimal normalization of $`\beta `$ we know that the untwisted part of the full partition function computed from the state (5.2) has a factor $`1/2`$ standing in front of each character. We pick $`\stackrel{~}{\alpha }`$ in such a way that there is a factor of $`1/2`$ in front of the twisted contribution as well. This means that we take $`\stackrel{~}{\alpha }`$ to be
$$\stackrel{~}{\alpha }^2=\frac{k+2}{4}.$$
(26)
For this choice to be consistent, all characters in the full partition function should come with integer coefficients. The twisted sector gives negative contributions for $`l^{}=2`$ mod $`4`$, positive contributions for $`l^{}=0`$ mod $`4`$ and no contributions for $`l^{}`$ odd. With our choice of $`\stackrel{~}{\alpha }`$, all contributing characters come with a factor $`\pm 1/2`$ in front.
From the untwisted part of the partition function we obtain characters with even $`l^{}`$. As we mentioned before, their coefficient is $`1/2`$. Characters with odd $`l^{}`$ are forbidden by the fusion rules. Putting everything together, we have shown that the sum of the twisted and the untwisted part of the partition function contains the characters with $`l^{}=0`$ mod $`4`$ with an integer coefficient. The characters with $`l^{}=2`$ mod $`4`$ get subtracted, so that they are not propagating in the open string sector.
Equation (26) for $`\stackrel{~}{\alpha }`$ admits the choice of a sign similar to the flat space case. In our general discussion, this corresponds to a character of the stabilizer subgroup. Since our orbifold group is $`_2`$, the latter reduces to the choice of a sign. The open string partition function has an opposite sign in front of the twisted part in the case that the two boundary states have opposite signs. In this case, the characters with $`l^{}=0`$ mod $`4`$ get removed from the partition function and those with $`l^{}=2`$ mod $`4`$ survive. This can be summarized in the following formulas
$`Z_{++}=Z_{}`$ $`=`$ $`{\displaystyle \underset{\nu =0}{\overset{k+1}{}}}{\displaystyle \underset{l^{},m^{},s^{}}{}}(1)^\nu \delta _l^{}^{(4)}\delta _{\nu +m^{}+\frac{k+2}{2}(\nu +s^{})}\delta _{\nu +s^{}}^{(2)}\chi _{l^{},m^{},s^{}}(q)`$
$`Z_+=Z_+`$ $`=`$ $`{\displaystyle \underset{\nu =0}{\overset{k+1}{}}}{\displaystyle \underset{l^{},m^{},s^{}}{}}(1)^\nu \delta _{l^{}+2}^{(4)}\delta _{\nu +m^{}+\frac{k+2}{2}(\nu +s^{})}\delta _{\nu +s^{}}^{(2)}\chi _{l^{},m^{},s^{}}(q).`$
Let us now discuss short orbit states in the full Gepner model. The length of a short orbit is $`K/2`$. In the individual minimal models, this means that $`l_i=k_i/2`$ in the case that $`2k_i+4`$, the orbit length of the single minimal model, does not divide $`K/2`$. In the case that $`K/2`$ is a multiple of the orbit length of a minimal model, an arbitrary primary can be chosen in that particular model. If the number of minimal models, for which we need $`l_i=k_i/2`$, is $`r^{}`$, the factors in the minimal models can be ordered in such a way that these minimal models are the factors $`1,\mathrm{},r^{}`$. This means that the labels of the short orbit boundary state are
$$\mathrm{\Lambda }_s=(\frac{k_1}{2},\mathrm{},\frac{k_r^{}}{2},L_{r^{}+1},\mathrm{},L_r)$$
(27)
with arbitrary labels $`\mathrm{\Xi }`$. Our general discussion of short orbit states and the discussion of the single minimal model in this section motivates the following ansatz for a short orbit boundary state in the full Gepner model
$`|\mathrm{\Lambda }_s,\mathrm{\Xi }_s`$ $`=`$ $`{\displaystyle \frac{1}{2}}|\mathrm{\Lambda }_s,\mathrm{\Xi }_s_{\text{proj}}+`$
$`+\alpha cN{\displaystyle \underset{\nu =0}{\overset{\frac{K}{2}1}{}}}{\displaystyle \underset{\nu _j=0,1}{}}{\displaystyle \underset{\mu }{}}{\displaystyle \underset{l_{r^{}+1},\mathrm{},l_r}{}}(1)^\nu (1)^{\frac{s_0^2}{2}}{\displaystyle \underset{j=r^{}+1}{\overset{r}{}}}{\displaystyle \frac{\mathrm{sin}(l_j,L_j)_{k_j}}{\sqrt{\mathrm{sin}(l_j,0)_{k_j}}}}\times `$
$`\times {\displaystyle \underset{j=1}{\overset{r}{}}}e^{i\pi \frac{M_j+\nu }{k_j+2}m_j}e^{i\pi \frac{s_j}{2}(S_j+\nu +2\nu _j)}e^{i\pi \frac{s_0}{2}(S_j+\nu 2{\scriptscriptstyle \nu _j})}|{\displaystyle \frac{k_1}{2}},\mathrm{},{\displaystyle \frac{k_r^{}}{2}},l_{r^{}+1},\mathrm{},l_r;\mu ^{\mathrm{tw}},`$
$`\text{xxxxxx}\text{ where}c=\left(\sqrt{2}{\displaystyle \underset{j}{}}\sqrt{\sqrt{2}(k_j+2)}\right)^1`$
and $`N=1/\sqrt{K2^r}`$. To determine the prefactor $`\alpha `$ of the twisted part of the boundary state, we compute the transition amplitude between the boundary state and itself and modular transform to the open string sector. This leads to the following result for the twisted part of the partition function
$`Z^{\mathrm{tw}}`$ $`=`$ $`\alpha ^2{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{K2^r}}{\displaystyle \underset{\nu }{}}{\displaystyle \underset{\nu _j}{}}{\displaystyle \underset{\lambda ^{},\mu ^{}}{}}(1)^\nu \delta _{\nu 2{\scriptscriptstyle \nu _j}+s_0^{}2}^{(4)}`$
$`{\displaystyle \underset{j=1}{\overset{r^{}}{}}}{\displaystyle \frac{2}{k_j+2}}\mathrm{sin}\pi {\displaystyle \frac{l_j^{}+1}{2}}\delta _{\nu +m_j^{}+\frac{(k_j+2)}{2}(\nu +2\nu _j+s_j^{})}^{(2k_j+4)}\delta _{\nu +2\nu _j+s_j^{}}^{(2)}`$
$`{\displaystyle \underset{j=r^{}+1}{\overset{r}{}}}N_{L_jL_j}^{l_j^{}}\delta _{\nu +m_j^{}}^{(2k_j+4)}\delta _{2\nu _j+\nu +s_j^{}}^{(4)}\chi _{\lambda ^{},\mu ^{}}(q).`$
The computation for the factors $`1,\mathrm{},r^{}`$ is like that for the short orbits in a single minimal model and the computation for the other factors works exactly as the computation for the projected boundary states of long orbits.
Given that we have characters with a prefactor of $`1/2`$ in the untwisted part of the partition function, the consistent choice for $`\alpha `$ is
$$\alpha =\pm \sqrt{K2^{rr^{}1}\underset{j=1}{\overset{r^{}}{}}(k_j+2)}.$$
(28)
Again, there are two different boundary states $`|\lambda _s,\mathrm{\Xi }_s;\pm `$ depending on the sign chosen for $`\alpha `$. Adding the twisted and untwisted part of the partition function, we see that characters with an odd $`l_j^{}`$ do not appear in both the untwisted and twisted part. Those with an even number of $`l_j^{}=2`$ mod $`4`$ add up and appear with multiplicity one in the total partition function. Those with an odd number of $`l_j^{}=2`$ mod $`4`$ appear with a negative sign in the twisted part of the partition function and get removed from the spectrum in the total partition function.
In the case of the Gepner model, there is usually more than one short orbit, and to check the consistency of the boundary states constructed above we have to compute the partition functions for all combinations of boundary states. The short orbits differ in the choice of $`L_j`$ in the minimal model factors which are not restricted to be on short orbits and in the $`\mathrm{\Xi }_s`$. The computation of the partition function for two different boundary states is not much different from that above. In the $`SU(2)`$ fusion coefficients, one of the $`L_j`$ is replaced by some other label $`\stackrel{~}{L}_j`$ for all $`j=r^{}+1,\mathrm{},r`$. Furthermore, one has to substitute the sum $`m_j^{}+\nu `$ with $`m_j^{}+\nu +M_j\stackrel{~}{M}_j`$ and similarly $`s_j^{}+\nu `$ with $`s_j^{}+\nu +S_j\stackrel{~}{S}_j`$.
## 6 Conclusions and Outlook
In this paper, we have constructed A-type boundary states in the Gepner model, which cannot be obtained by the projection method. These boundary states contain Ishibashi states from twisted sectors. In particular, they carry RR-charge under the RR-fields in the twisted sector of the Gepner model. Geometrically, this means that the branes wrap exceptional cycles, i.e. cycles coming from the resolution of a singularity. It has been shown in examples that the rank of the intersection matrix computed from projected boundary states equals the number of homology three-cycles that do not come from the resolution. This clearly demonstrates that the projection method is missing interesting boundary states whenever the corresponding CY-manifold develops singular curves. Some of these extra boundary states are now provided by our construction.
In this paper, we considered only A-type boundary states. B-type boundary states in Gepner models have been obtained in . It is possible to reinterpret the formulas given there in terms of simple currents orbifolds.
For B-type boundary states, the RR-charges associated with $`(a,c)`$-fields become important and replace the $`(c,c)`$-fields in our discussion above. It is well known that mirror symmetry relates these two different types of bulk fields, i.e. the space of $`(a,c)`$-fields on the original 3-fold $`M`$ gets mapped into the $`(c,c)`$-fields on the mirror $`W`$ and vice versa. Let us recall that mirror symmetry of closed string theories can be understood as an orbifold construction with the group
$$𝒢=\left(\underset{i=1}{\overset{r}{}}𝒢_i\right)/\mathrm{\Gamma }^{}\left(\underset{i=1}{\overset{r}{}}_{k_i+2}\right)/_{\frac{K}{2}}$$
where $`𝒢_i`$ is the group generated by the element $`(2;0,\mathrm{},m_i=2,\mathrm{},0;0,\mathrm{},s_i=2,\mathrm{},0)`$ and $`\mathrm{\Gamma }^{}`$ is the intersection between the orbifold group $`\mathrm{\Gamma }`$ and the product of groups $`𝒢_i`$ in the nominator. Most importantly, this orbifold construction allows to compute the bulk partition function for the theory on the mirror $`W`$.
Following our general discussion above, one can extend these constructions of the mirror theory to obtain A-type boundary states on $`W`$ through an orbifold procedure with the group $`\mathrm{\Gamma }\times 𝒢`$. Obviously, the resulting spectra in the open string sector will be organized in terms of orbits of this orbifold group. On the other hand, orbits of $`\mathrm{\Gamma }\times 𝒢`$ do show up in the explicit formulas for open string spectra of B-type boundary states in the string theory on the original 3-fold $`M`$ . The fact that the partition functions of $`\mathrm{\Gamma }\times 𝒢`$-projected (A-type) D-branes on $`W`$ coincide with those found in for B-type branes on $`M`$ confirms that the orbifold description of mirror symmetry extends to open strings (cp. also ).
In the detailed analysis of the $`\mathrm{\Gamma }\times 𝒢`$-orbifold theories one finds a rich pattern of orbits with short orbits of various lengths appearing for a typical model. This signals the existence of additional boundary states which are not obtained by projection from the tensor product theory. In principal, these states can be constructed along the lines of Section 5 above. Only technically this becomes more involved due to the possibly complicated orbit structure. The resulting B-type boundary states on $`M`$ are charged under fields which arise from desingularization of singular points and of singular curves.
In B-type boundary states have been compared with vector bundles on Calabi-Yau manifolds, which are elliptic or K3 fibrations. For short orbit states it was found that there are a number of vacua propagating. The interpretation of was that these boundary states do not correspond to elementary branes. This is indeed supported by our analysis. The elementary short orbit states should contain a twisted part.
Let us remark that our techniques can also be applied to boundary states in bulk theories with a D-type modular invariant. Boundary states for such theories have been discussed recently in . In the context of simple current constructions, D-even modular invariants can be constructed by modding out an additional $`_2`$ current in a minimal model. The open string sector is then organized in orbits of this current.
The CFT-construction we have used in this work allow to obtain a large number of boundary states extending the set of states which were known before, if the underlying Calabi-Yau 3-fold possesses $`_2`$-singularities along curves. It is interesting to point out once more that the geometric scenario of Section 4 suggests the existence of many other important D-brane states that are charged under $`(c,c)`$-fields from the $`K/N`$-twisted sectors when $`N2`$. The latter come with the desingularization of a $`_{N2}`$-singularity over a curve $`C`$ on $`M`$. While boundary states charged under such $`K/N`$-twisted $`(c,c)`$-fields satisfy A-type gluing conditions for the diagonally embedded superconformal algebra, it would be inconsistent with the U(1) charges of the corresponding $`(c,c)`$-fields to require such A-type gluing conditions for each minimal model separately. But so far, all treatments of D-branes in Gepner models assumed the strong version of the gluing condition which preserves all the individual superconformal algebras. It remains an interesting open problem to relax this requirement and, in particular, to construct $`D`$-branes associated with the resolution of $`_{N2}`$ singularities over curves $`C`$ in $`M`$.
Acknowledgements: We would like to thank M. Bianchi, R. Blumenhagen, M. Douglas, H. Liu, A. Recknagel, C. Römelsberger, R. Schimmrigk, Y. Stanev and S. Theisen for very useful discussions. V.S. is grateful to the Rutgers string group for their hospitality. I.B. thanks the II. Institut für Theoretische Physik for supporting her stay at Hamburg. We also want to thank A. Recknagel for his helpful remarks on the manuscript. |
warning/0001/cond-mat0001011.html | ar5iv | text | # Spin-Dependent Josephson Current through Double Quantum Dots and Measurement of Entangled Electron States
\[
## Abstract
We study a double quantum dot each dot of which is tunnel-coupled to superconducting leads. In the Coulomb blockade regime, a spin-dependent Josephson coupling between two superconductors is induced, as well as an antiferromagnetic Heisenberg exchange coupling between the spins on the double dot which can be tuned by the superconducting phase difference. We show that the correlated spin states—singlet or triplets—on the double dot can be probed via the Josephson current in a dc-SQUID setup.
PACS numbers: 73.23.-b, 73.23.Hk, 74.50.+r
\]
In recent years, electronic transport through strongly interacting mesoscopic systems has been the focus of many investigations . In particular, a single quantum dot coupled via tunnel junctions to two non-interacting leads has provided a prototype model to study Coulomb blockade effects and resonant tunneling in such systems. These studies have been extended to an Anderson impurity or a quantum dot coupled to superconductors . In a number of experimental and theoretical papers, the spectroscopic properties of a quantum dot coupled to two superconductors have been studied. Further, an effective dc Josephson effect through strongly interacting regions between superconducting leads has been analyzed . More recently, on the other hand, research on the possibility to control and detect the spin of electrons through their charges has started. In particular in semiconducting nanostructures, it was found that the direct coupling of two quantum dots by a tunnel junction can be used to create entanglement between spins , and that such spin correlations can be observed
in charge transport experiments .
Motivated by these studies we propose in the present work a new scenario for inducing and detecting spin correlations, viz., coupling a double quantum dot (DD) to superconducting leads by tunnel junctions as shown in Fig. 1. It turns out that this connection via a superconductor induces a Heisenberg exchange coupling between the two spins on the DD. Moreover, if the DD is arranged between two superconductors (see Fig. 1), we obtain a Josephson junction (S-DD-S). The resulting Josephson current depends on the spin state of the DD and can be used to probe the spin correlations on the DD.
Model— The double-dot (DD) system we propose is sketched in Fig. 1: Two quantum dots ($`a`$,$`b`$), each of which contains one (excess) electron and is connected to two superconducting leads ($`L,R`$) by tunnel junctions (indicated by dashed lines) . There is no direct coupling between the two dots. The Hamiltonian describing this system consists of three parts, $`H_S+H_D+H_TH_0+H_T`$. The leads are assumed to be conventional singlet superconductors that are described by the BCS Hamiltonian
$`H_S`$ $`=`$ $`{\displaystyle \underset{j=L,R}{}}{\displaystyle _{\mathrm{\Omega }_j}}{\displaystyle \frac{d𝐫}{\mathrm{\Omega }_j}}\{{\displaystyle \underset{\sigma =,}{}}\psi _\sigma ^{}(𝐫)h(𝐫)\psi _\sigma (𝐫)`$ (2)
$`+\mathrm{\Delta }_j(𝐫)\psi _{}^{}(𝐫)\psi _{}^{}(𝐫)+h.c.\},`$
where $`\mathrm{\Omega }_j`$ is the volume of lead $`j`$, $`h(𝐫)=(i\mathrm{}+\frac{e}{c}𝐀)^2/2m\mu `$, and $`\mathrm{\Delta }_j(𝐫)=\mathrm{\Delta }_je^{i\varphi _j(𝐫)}`$ is the pair potential. For simplicity, we assume identical leads with same chemical potential $`\mu `$, and $`\mathrm{\Delta }_L=\mathrm{\Delta }_R=\mathrm{\Delta }`$. The two quantum dots are modelled as two localized levels $`ϵ_a`$ and $`ϵ_b`$ with strong on-site Coulomb repulsion $`U`$, described by the Hamiltonian
$$H_D=\underset{n=a,b}{}\left[ϵ\underset{\sigma }{}d_{n\sigma }^{}d_{n\sigma }+Ud_n^{}d_nd_n^{}d_n\right],$$
(3)
where we put $`ϵ_a=ϵ_b=ϵ`$ ($`ϵ>0`$) for simplicity. $`U`$ is typically given by the charging energy of the dots, and we have assumed that the level spacing of the dots is $`U`$ (which is the case for small GaAs dots), so that we need to retain only one energy level in $`H_D`$. Finally, the DD is coupled in parallel (see Fig. 1) to the superconducting leads, described by the tunneling Hamiltonian
$$H_T=\underset{j,n,\sigma }{}[te^{i\frac{\pi }{\mathrm{\Phi }_0}_{𝐫_n}^{𝐫_{j,n}}𝑑𝐥𝐀}\psi _\sigma ^{}(𝐫_{j,n})d_{n\sigma }+h.c.],$$
(4)
where $`𝐫_{j,n}`$ is the point on the lead $`j`$ closest to the dot $`n`$. Here, $`\mathrm{\Phi }_0=hc/2e`$ is the superconducting flux quantum. Unless mentioned otherwise, it will be assumed that $`𝐫_{L,a}=𝐫_{L,b}=𝐫_L`$ and $`𝐫_{R,a}=𝐫_{R,b}=𝐫_R`$.
Since the low-energy states of the whole system are well separated by the superconducting gap $`\mathrm{\Delta }`$ as well as the strong Coulomb repulsion $`U`$ ($`\mathrm{\Delta },ϵUϵ`$), it is sufficient to consider an effective Hamiltonian on the reduced Hilbert space consisting of singly occupied levels of the dots and the BCS ground states on the leads. To lowest order in $`H_T`$, the effective Hamiltonian is
$$H_{\mathrm{𝑒𝑓𝑓}}=PH_T\left[(E_0H_0)^1(1P)H_T\right]^3P,$$
(5)
where $`P`$ is the projection operator onto the subspace and $`E_0`$ is the ground-state energy of the unperturbed Hamiltonian $`H_0`$. (The 2nd order contribution leads to an irrelevant constant.) The lowest-order expansion (5) is valid in the limit $`\mathrm{\Gamma }\mathrm{\Delta },ϵ`$ where $`\mathrm{\Gamma }=\pi t^2N(0)`$ and $`N(0)`$ is the normal-state density of states per spin of the
leads at the Fermi energy. Thus, we assume that $`\mathrm{\Gamma }\mathrm{\Delta },ϵUϵ`$, and temperatures which are less than $`ϵ`$ (but larger than the Kondo temperature).
Effective Hamiltonian — There are a number of virtual hopping processes that contribute to the effective Hamiltonian (5), see Fig. 2 for a partial listing of them. Collecting these various processes, one can get the effective Hamiltonian in terms of the gauge-invariant phase differences $`\varphi `$ and $`\phi `$ between the superconducting leads and the spin operators $`𝐒_a`$ and $`𝐒_b`$ of the dots (up to a constant and with $`\mathrm{}=1`$)
$`H_{\mathrm{𝑒𝑓𝑓}}`$ $`=`$ $`J_0\mathrm{cos}(\pi f_{AB})\mathrm{cos}(\varphi \pi f_{AB})`$ (6)
$`+`$ $`\left[(2J_0+J)(1+\mathrm{cos}\phi )+2J_1(1+\mathrm{cos}\pi f_{AB})\right]`$ (8)
$`\times \left[𝐒_a𝐒_b1/4\right].`$
Here $`f_{AB}=\mathrm{\Phi }_{AB}/\mathrm{\Phi }_0`$ and $`\mathrm{\Phi }_{AB}`$ is the Aharonov-Bohm (AB) flux threading through the closed loop indicated by the dashed lines in Fig. 1. One should be careful to define gauge-invariant phase differences $`\varphi `$ and $`\phi `$ in (6). The phase difference $`\varphi `$ is defined as usual by
$$\varphi =\varphi _L(𝐫_L)\varphi _R(𝐫_R)\frac{2\pi }{\mathrm{\Phi }_0}_{𝐫_R}^{𝐫_L}𝑑\mathrm{}_a𝐀,$$
(9)
where the integration from $`𝐫_R`$ to $`𝐫_L`$ runs via dot $`a`$ (see Fig. 1). The second phase difference, $`\phi `$, is defined by
$$\phi =\varphi _L(𝐫_L)\varphi _R(𝐫_R)\frac{\pi }{\mathrm{\Phi }_0}_{𝐫_R}^{𝐫_L}(d\mathrm{}_a+d\mathrm{}_b)𝐀.$$
(10)
The distinction between $`\varphi `$ and $`\phi `$, however, is not significant unless one is interested in the effects of an AB flux through the closed loop in Fig. 1 (see Ref. for an example of such effects). The coupling constants appearing in (6) are defined by
$`J`$ $`=`$ $`{\displaystyle \frac{2\mathrm{\Gamma }^2}{ϵ}}\left[{\displaystyle \frac{1}{\pi }}{\displaystyle \frac{dx}{f(x)g(x)}}\right]^2`$ (11)
$`J_0`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }^2}{\mathrm{\Delta }}}{\displaystyle \frac{dxdy}{\pi ^2}\frac{1}{f(x)f(y)[f(x)+f(y)]g(x)g(y)}}`$ (12)
$`J_1`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }^2}{\mathrm{\Delta }}}{\displaystyle \frac{dxdy}{\pi ^2}\frac{g(x)[f(x)+f(y)]2\zeta g(y)}{g(x)^2g(y)[g(x)+g(y)][f(x)+f(y)]}},`$ (13)
where $`\zeta =ϵ/\mathrm{\Delta }`$, $`f(x)=\sqrt{1+x^2}`$, and $`g(x)=\sqrt{1+x^2}+\zeta `$.
Eq. (6) is one of our main results. A remarkable feature of it is that a Heisenberg exchange coupling between the spin on dot $`a`$ and on dot $`b`$ is induced by the superconductor. This coupling is antiferromagnetic (all $`J`$’s are positive) and thus favors a singlet ground state of spin $`a`$ and $`b`$. This in turn is a direct consequence of the assumed singlet nature of the Cooper pairs in the superconductor . As discussed below, an immediate observable consequence of $`H_{\mathrm{𝑒𝑓𝑓}}`$ is a spin-dependent Josephson current from the left to right superconducting lead (see Fig. 1) which probes the correlated spin state on the DD.
The various terms in (6) have different magnitudes. In particular, the processes leading to the $`J_1`$ term involve quasiparticles only as can be seen from its AB-flux dependence which has period $`2\mathrm{\Phi }_0`$. In the limits we will consider below, this $`J_1`$ term is small and can be neglected.
In the limit $`\zeta 1`$, the main contributions come from processes of the type depicted in Fig. 2 (a) and (b), making $`J_00.1(\mathrm{\Gamma }^2/\zeta ϵ)\mathrm{log}\zeta `$ dominant over $`J`$ and $`J_1`$. Thus, (6) can be reduced to
$`H_{\mathrm{𝑒𝑓𝑓}}`$ $``$ $`J_0\mathrm{cos}(\pi f_{AB})\mathrm{cos}(\varphi \pi f_{AB})`$ (14)
$`+`$ $`2J_0(1+\mathrm{cos}\phi )\left[𝐒_a𝐒_b{\displaystyle \frac{1}{4}}\right],`$ (15)
up to order $`(\mathrm{log}\zeta )/\zeta `$. As can be seen in Fig. 2 (a), the first term in (14) has the same origin as that in the single-dot case : Each dot separately constitutes an effective Josephson junction with coupling energy $`J_0/2`$ (i.e. $`\pi `$-junction) between the two superconductors. The two resulting junctions form a dc SQUID, leading to the total Josephson coupling in the first term of (14). The Josephson coupling in the second term in (14), corresponding to processes of type Fig. 2 (b), depends on the correlated spin states on the double dot: For the singlet state, it gives an ordinary Josephson junction with coupling $`2J_0`$ and competes with the first term, whereas it vanishes for the triplet states. Although the limit $`\mathrm{\Delta }ϵUϵ`$ is not easy to achieve with present-day technology, such a regime is relevant, say, for two atomic impurities embedded between the grains of a granular superconductor.
More interesting and experimentally feasible is the case $`\zeta 1`$. In this regime, the effective Hamiltonian (6) is dominated by a single term (up to terms of order $`\zeta `$),
$$H_{\mathrm{𝑒𝑓𝑓}}J(1+\mathrm{cos}\phi )\left[𝐒_a𝐒_b\frac{1}{4}\right],$$
(16)
with $`J2\mathrm{\Gamma }^2/ϵ`$. The processes of type Fig. 2 (b) and (c) give rise to (16). Below we will propose an experimental setup based on (16).
Before proceeding, we digress briefly on the dependence of $`J`$ on the contact points. Unlike the processes of type Fig. 2 (a), those of types Fig. 2 (b), (c), and (d) depend on $`\delta r_L=|𝐫_{L,a}𝐫_{L,b}|`$ and $`\delta r_R=|𝐫_{R,a}𝐫_{R,b}|`$, see the remark below Eq. (4). For the tunneling Hamiltonian (4), one gets (putting $`\delta r=\delta r_L=\delta r_R`$)
$$J(\delta r)=\frac{8t^4}{ϵ}\left|_0^{\mathrm{}}\frac{d\omega }{2\pi }\frac{F^R(\delta r,\omega )F^A(\delta r,\omega )}{\omega +ϵ}\right|^2,$$
(17)
where $`F^{R/A}(𝐫,\omega )`$ is the Fourier transform of the Green’s function in the superconductors, $`F^{R/A}(𝐫,t)=i\mathrm{\Theta }(\pm t)\{\psi _{}(𝐫,t),\psi _{}(0,0)\}`$ . E.g., in the limit $`\epsilon \mathrm{\Delta }\mu `$, we find $`J(\delta r)J(0)e^{2\delta r/\xi }\mathrm{sin}^2(k_F\delta r)/(k_F\delta r)^2`$ up to order $`1/k_F\xi `$, with $`k_F`$ the Fermi wave vector in the leads. Thus, to have $`J(\delta r)`$ non-zero, $`\delta r`$ should not exceed the superconducting coherence length $`\xi `$.
Probing spins with a dc-SQUID — We now propose a possible experimental setup to probe the correlations (entanglement) of the spins on the dots, based on the
effective model (16). According to (16) the S-DD-S structure can be regarded as a spin-dependent Josephson junction. Moreover, this structure can be connected with an ordinary Josephson junction to form a dc-SQUID-like geometry, see Fig. 3. The Hamiltonian of the entire system is then given by
$`H`$ $`=`$ $`J[1+\mathrm{cos}(\theta 2\pi f)]\left(𝐒_a𝐒_b{\displaystyle \frac{1}{4}}\right)`$ (19)
$`+\alpha J(1\mathrm{cos}\theta ),`$
where $`f=\mathrm{\Phi }/\mathrm{\Phi }_0`$, $`\mathrm{\Phi }`$ is the flux threading the SQUID loop, $`\theta `$ is the gauge-invariant phase difference across the auxiliary junction ($`J^{}`$), and $`\alpha =J^{}/J`$ with $`J^{}`$ being the Josephson coupling energy of the auxiliary junction. One immediate consequence of (19) is that at zero temperature, we can effectively turn on and off the spin exchange interaction: For half-integer flux ($`f=1/2`$), singlet and triplet states are degenerate at $`\theta =0`$. Even at finite temperatures, where $`\theta `$ is subject to thermal fluctuations, singlet and triplet states are almost degenerate around $`\theta =0`$. On the other hand, for integer flux ($`f=0`$), the energy of the singlet is lower by $`J`$ than that of the triplets.
This observation allows us to probe directly the spin state on the double dot via a Josephson current across the dc-SQUID-like structure in Fig. 3. The supercurrent through the SQUID-ring is defined as $`I_S=(2\pi c/\mathrm{\Phi }_0)H/\theta `$, where the brackets refer to a spin expectation value on the DD. Thus, depending on the spin state on the DD we find
$$I_S/I_J=\{\begin{array}{cc}\mathrm{sin}(\theta 2\pi f)+\alpha \mathrm{sin}\theta \hfill & \text{(singlet)}\hfill \\ \alpha \mathrm{sin}\theta \hfill & \text{(triplets)}\hfill \end{array},$$
(20)
where $`I_J=2eJ/\mathrm{}`$. When the system is biased by a dc current $`I`$ larger than the spin- and flux-dependent critical current, given by $`\mathrm{max}_\theta \{|I_S|\}`$, a finite voltage $`V`$ appears. Then one possible experimental procedure might be as follows (see Fig. 4). Apply a dc bias current such that $`\alpha I_J<I<(\alpha +1)I_J`$. Here, $`\alpha I_J`$ is the critical current of the triplet states, and $`(\alpha +1)I_J`$ the critical current of the singlet state at $`f=0`$, see (20). Initially prepare the system in an equal mixture of singlet and triplet states by tuning the flux around $`f=1/2`$. (With electron $`g`$-factors $`g0.5\text{}20`$ the Zeeman splitting on the dots is usually small compared with $`k_BT`$ and can thus be ignored.) The dc voltage measured in this mixture will be given by $`(V_0+3V_1)/4`$, where $`V_0(V_1)2\mathrm{\Delta }/e`$ is the (current-dependent) voltage drop associated with the singlet (triplet) states. At a later time $`t=0`$, the flux is switched off (i.e. $`f=0`$), with $`I`$ being kept fixed. The ensuing time evolution of the system is characterized by three time scales: the time $`\tau _{\mathrm{𝑐𝑜ℎ}}\mathrm{max}\{1/\mathrm{\Delta },1/\mathrm{\Gamma }\}1/\mathrm{\Gamma }`$ it takes to establish coherence in the S-DD-S junction, the spin relaxation time $`\tau _{\mathrm{𝑠𝑝𝑖𝑛}}`$ on the dot, and the switching time $`\tau _{\mathrm{𝑠𝑤}}`$ to reach $`f=0`$. We will assume $`\tau _{\mathrm{𝑐𝑜ℎ}}\tau _{\mathrm{𝑠𝑝𝑖𝑛}},\tau _{\mathrm{𝑠𝑤}}`$, which is not unrealistic in view of measured spin decoherence times in GaAs exceeding $`100`$ ns . If $`\tau _{\mathrm{𝑠𝑤}}<\tau _{\mathrm{𝑠𝑝𝑖𝑛}}`$, the voltage is given by $`3V_1/4`$ for times less than $`\tau _{\mathrm{𝑠𝑝𝑖𝑛}}`$, i.e. the singlet no longer contributes to the voltage. For $`t>\tau _{\mathrm{𝑠𝑝𝑖𝑛}}`$ the spins have relaxed to their ground (singlet) state, and the voltage vanishes. One therefore expects steps in the voltage versus time (solid curve in Fig. 4). If $`\tau _{\mathrm{𝑠𝑝𝑖𝑛}}<\tau _{\mathrm{𝑠𝑤}}`$, a broad transition region of the voltage from the initial value to $`0`$ will occur (dashed line in Fig. 4) .
To our knowledge, there are no experimental reports on quantum dots coupled to superconductors. However, hybrid systems consisting of superconductors (e.g., Al or Nb) and 2DES (InAs and GaAs) have been investigated by a number of groups . Taking the parameters of those materials, a rough estimate leads to a coupling energy $`J`$ in (16) or (19) of about $`J\text{0.05–0.5K}`$. This corresponds to a critical current scale of $`I_J\text{5–50}\text{nA}`$.
In conclusion, we have investigated double quantum dots each dot of which is coupled to two superconductors. We have found that in the Coulomb blockade regime the Josephson current from one superconducting lead to the other is different for singlet or triplet states on the double dot. This leads to the possibility to probe the spin states of the dot electrons by measuring a Josephson current.
Acknowledgment — We would like to thank G. Burkard, C. Strunk, and E. Sukhorukov for discussions and the Swiss NSF for support. |
warning/0001/hep-ph0001043.html | ar5iv | text | # I Introduction
## I Introduction
Chiral gauge theories, in which at least part of the matter field content is in complex representations of the gauge group, play an important role in efforts to extend the standard model. These include grand unified theories, dynamical breaking of symmetries, and theories of quark and lepton substructure. An important distinction from vector-like theories such as QCD is that since at least some of the chiral symmetries are gauged, mass terms that would explicitly break these chiral symmetries are forbidden in the Lagrangian. Another key feature is that the fermion content is subject to a constraint not present in vectorial gauge theories, the cancellation of gauge and gravitational anomalies.
Chiral theories received much attention in the 1980’s , focusing on their strong coupling behavior in the infrared. One possibility is confinement with the gauge symmetry as well as global symmetries unbroken, realized by the formation of gauge singlet, massless composite fermions. Another is confinement with intact gauge symmetry but with some of the global symmetries broken spontaneously, leading to the formation of gauge-singlet Goldstone bosons. It is also possible for these theories to exist in the Higgs phase, dynamically breaking their own gauge symmetries . Depending on particle content, they might even remain weakly coupled. This will happen if the theory has an interacting but weak infrared fixed point. The symmetries will then remain unbroken, and the infrared and underlying degrees of freedom will be the same.
Supersymmetric (SUSY) chiral theories have also received considerable attention over the years, since most of the known examples of dynamical supersymmetry breaking involve these kinds of theories .
Studies of chiral gauge theories have typically made use of the ’t Hooft global anomaly matching conditions along with $`1/N`$ expansion, and not-so-reliable most attractive channel (MAC) analysis and instanton computations. Direct approaches using strong coupling lattice methods are still difficult. Another indirect approach developed recently takes the form of an inequality limiting the number of massless degrees of freedom in the infrared description of a field theory to be no larger than the number of ultraviolet degrees of freedom. It was conjectured to apply to all asymptotically free theories whose infrared behavior is also governed by a fixed point, not necessarily free.
The inequality is formulated using finite temperature as a device to probe all energy scales, with the degree-of-freedom count defined using the free energy of the field theory. The zero-temperature theory of interest is characterized using the quantity $`f_{IR}`$, related to the free energy by
$$f_{IR}\underset{T0}{lim}\frac{(T)}{T^4}\frac{90}{\pi ^2},$$
(1)
where $`T`$ is the temperature and $``$ is the conventionally defined free energy per unit volume. The limit is well defined if the theory has an infrared fixed point. For the special case of an infrared-free theory, $`f_{IR}`$ is simply the number of massless bosons plus $`7/4`$ times the number of 2-component massless Weyl fermions. The corresponding expression in the large $`T`$ limit is
$$f_{UV}\underset{T\mathrm{}}{lim}\frac{(T)}{T^4}\frac{90}{\pi ^2}.$$
(2)
This limit will be well defined if the theory has an ultraviolet fixed point. For an asymptotically free theory $`f_{UV}`$ counts the underlying, ultraviolet degrees of freedom in a similar way.
In terms of these quantities, the conjectured inequality for any asymptotically-free theory is
$$f_{IR}f_{UV}.$$
(3)
This inequality has not been proven, but in Ref. it was shown to agree with all known results and then used to derive new results for several strongly coupled, vector-like gauge theories. It was applied to chiral theories in Ref. . The principal focus there was on the possibility of preserving the global symmetries through the formation of massless composite fermions.
In this paper, we examine further the two non-supersymmetric chiral theories of Ref. , both rich in possible phase structure. One is the Bars-Yankielowicz (BY) model involving fermions in the two-index symmetric tensor representation, and the other is a generalized Georgi-Glashow (GGG) model involving fermions in the two-index antisymmetric tensor representation. In each case, in addition to fermions in complex representations, a set of $`p`$ anti fundamental-fundamental pairs are included and the allowed phases are considered as a function of $`p`$. Several possible phases emerge, consistent with global anomaly matching and the above inequality.
Graphs for the various $`f_{IR}`$’s vs $`p`$ are plotted for each model. They lead us to the suggestion that each of these theories will choose from among the allowed phases the one that minimizes $`f_{IR}`$ at each value of $`p`$. This idea is useful only to distinguish among phases for which the $`f_{IR}`$’s are computable, those that are infrared free or governed by a weak infrared fixed point. It may be rephrased in terms of the entropy per unit volume $`S(T)`$ of the system (the derivative of the pressure per unit volume $`P(T)=(T)`$). Provided only that the limit in Eq. (1) exists, $`S(T)`$ near freeze out will be given by $`S(T)=(2\pi ^2/45)T^3f_{IR}`$ plus higher order terms in the low temperature expansion. Thus the minimum number of degrees of freedom corresponds to the minimum entropy at approach to freeze out, consistent with global anomaly matching. Now of course at any finite T, the preferred phase is chosen from among all the states entering the partition function by minimizing the free energy density, which becomes the energy density at $`T=0`$. Comparing these quantities for different states when the theory is strongly interacting is, however, generally a strong coupling problem. The conjecture for these theories is the following. Assuming that the candidate phases revealed by anomaly matching are states entering the partition function, the state with the relatively lower energy density is approached at a lower rate proportional to the lower $`f_{IR}`$. The latter quantity is computable in terms of only the effective low energy theory which is infrared free for the phases being considered here.
We focus almost completely on symmetry breaking patterns corresponding to the formation of bilinear condensates. We suggest that in general, the phase corresponding to confinement with all symmetries unbroken, where all the global anomalies are matched by massless composite fermions, is not preferred. Instead, the global symmetries associated with fermions in real representations break spontaneously via bilinear condensate formation as in QCD. With respect to the fermions in complex representations, however, the formation of bilinear condensates is suggested to be disfavored relative to confinement and the preservation of the global symmetries via massless composite fermion formation.
It is interesting to note that in the real world case of two-flavor QCD (a vector-like theory with all fermions in a real representation) nature prefers to minimize $`f_{IR}`$. Neglecting the small bare quark masses, global anomaly matching admits two possible low energy phases, broken chiral symmetry through the formation of the bilinear $`<\overline{\psi }\psi >`$ condensate, or unbroken chiral symmetry through the formation of confined massless baryons. Both effective low energy theories are infrared free. The three Goldstone bosons of the former (chosen by nature) lead to $`f_{IR}=3`$, and the two massless composite Dirac fermions of the latter lead to $`f_{IR}=7`$.
Bilinear condensate formation is of course not the only possibility in a strongly coupled gauge field theory. We have so far extended our discussion to include general condensate formation for one simple example, the $`SU(5)`$ Georgi-Glashow model, which has fermions in only complex representations and has only a $`U(1)`$ global symmetry. This symmetry can be broken via only higher dimensional condensates. Interestingly, this breaking pattern, with confinement and unbroken gauge symmetry, leads to the minimum value of $`f_{IR}`$. This highlights the important general question of the pattern of symmetry breaking in chiral theories when arbitrary condensate formation is considered. Higher dimensional condensates might play an important role, for example, in the dynamical breaking of symmetries in extensions of the standard model .
The Bars-Yankielowicz model is discussed in Section II and the generalized Georgi-Glashow model is discussed in Section III. In Section IV, we briefly describe two supersymmetric chiral theories: the supersymmetric version of the $`SU(5)`$ Georgi-Glashow model and the closely related $`32`$ model. In Section V we summarize and conclude.
## II The Bars Yankielowicz (BY) Model
This model is based on the single gauge group $`SU(N3)`$ and includes fermions transforming as a symmetric tensor representation, $`S=\psi _L^{\{ab\}}`$, $`a,b=1,\mathrm{},N`$; $`N+4+p`$ conjugate fundamental representations: $`\overline{F}_{a,i}=\psi _{a,iL}^c`$, where $`i=1,\mathrm{},N+4+p`$; and $`p`$ fundamental representations, $`F^{a,i}=\psi _L^{a,i},i=1,\mathrm{},p`$. The $`p=0`$ theory is the basic chiral theory, free of gauge anomalies by virtue of cancellation between the antisymmetric tensor and the $`N+4`$ conjugate fundamentals. The additional $`p`$ pairs of fundamentals and conjugate fundamentals, in a real representation of the gauge group, lead to no gauge anomalies.
The global symmetry group is
$$G_f=SU(N+4+p)\times SU(p)\times U_1(1)\times U_2(1).$$
(4)
Two $`U(1)`$’s are the linear combination of the original $`U(1)`$’s generated by $`Se^{i\theta _S}S`$ , $`\overline{F}e^{i\theta _{\overline{F}}}\overline{F}`$ and $`Fe^{i\theta _F}F`$ that are left invariant by instantons, namely that for which $`_jN_{R_j}T(R_j)Q_{R_j}=0`$, where $`Q_{R_j}`$ is the $`U(1)`$ charge of $`R_j`$ and $`N_{R_j}`$ denotes the number of copies of $`R_j`$.
Thus the fermionic content of the theory is
(5)
where the first $`SU(N)`$ is the gauge group, indicated by the square brackets.
For all the models considered in this paper, the beta function is generically written as
$$\beta =\mu \frac{d\alpha }{d\mu }=\beta _1(\frac{\alpha ^2}{2\pi })\beta _2(\frac{\alpha ^3}{4\pi ^2})+O(\alpha ^4),$$
(6)
where the terms of order $`\alpha ^4`$ and higher are scheme-dependent. For the present model, we have $`\beta _1=3N2(2/3)p`$ and $`\beta _2=(1/4)\{13N^230N+1+12/N2p((13/3)N1/N)\}`$. Thus the theory is asymptotically free for
$$p<(9/2)N3.$$
(7)
We shall restrict $`p`$ so that this condition is satisfied.
Because of asymptotic freedom, the thermodynamic free-energy may be computed in the $`T\mathrm{}`$ limit. An enumeration of the degrees of freedom leads to
$$f_{UV}=2(N^21)+\frac{7}{4}[\frac{N(N+1)}{2}+(N+4)N+2pN].$$
(8)
The infrared realization of this theory will vary depending on the number $`p`$ of conjugate fundamental-fundamental pairs. We begin by discussing the $`p=0`$ theory and then map out the phase structure as function of $`p`$.
### A The $`p=0`$ Case
For $`p=0`$, the fermions are in complex representations of the $`SU(N)`$ gauge group and the global symmetry group is $`G_f=SU(N+4)\times U_1(1)`$. The theory is strongly coupled at low energies, so it is expected either to confine or to break some of the symmetries, consistent with global anomaly matching.
The possibility that the $`p=0`$ theory confines with the full global symmetry group $`G_f`$ unbroken has been considered previously in the literature . All the global anomalies of the underlying theory may be matched at low energies providing that the massless spectrum is composed of gauge singlet composite fermions transforming according to the antisymmetric second-rank tensor representation of $`SU(N+4)`$. They are described by the composite operators $`\overline{F}_{[i}S\overline{F}_{j]}`$ and have charge $`N`$ under the $`U_1(1)`$ global symmetry.
With only these massless composites in the low energy spectrum, there are no dimension-four interactions, so the composites are noninteracting in the infrared. Therefore the thermodynamic free energy may be computed in the limit $`T0`$. Enumerating the degrees of freedom gives
$$f_{IR}^{sym}(p=0)=\frac{7}{4}\frac{(N+4)(N+3)}{2},$$
(9)
where the superscript indicates that the full global symmetry is intact. Clearly $`f_{IR}^{sym}(p=0)<f_{UV}(p=0)`$, satisfying the inequality of Eq. (3.
While the formation of confined massless composite fermions and the preservation of $`G_f`$ is consistent with anomaly matching and the thermal inequality, the same can be seen to be true of broken symmetry channels. We consider first the Higgs phase corresponding to the maximally attractive channel. It is
$`\begin{array}{cc}& \end{array}\times \overline{\begin{array}{c}\end{array}}\begin{array}{c}\end{array},`$
leading to the formation of the $`S\overline{F}`$ condensate
$$\epsilon ^{\gamma \delta }S_\gamma ^{ai}\overline{F}_{\{a,i\},\delta },$$
(10)
where $`\gamma ,\delta =1,2`$ are spin indices and $`a,i=1,\mathrm{},N`$ are gauge and flavor indices. This condensate breaks $`U_1(1)`$ and all the gauge symmetries, and it breaks $`SU(N+4)`$ to $`SU(4)`$. But the $`SU(N)`$ subgroup of $`SU(N+4)`$ combines with the gauge group, leading to a new global symmetry $`SU^{}(N)`$. For this group, $`\overline{F}_{a,iN}`$ is reducible, to the symmetric $`\overline{F}^S=\overline{F}_{\{a,i\}}`$ and the anti-symmetric $`\overline{F}^A=\overline{F}_{[a,i]}`$ representations.
The broken $`SU(N+4)`$ generator
$`Q_{(N+4)}=\left(\begin{array}{c}\text{4}\text{}\text{}\text{}\text{}\text{}\\ & \mathrm{}& & & & \\ & & 4& & & \\ & & & N& & \\ & & & & \mathrm{}& \\ & & & & & N\end{array}\right),`$
combines with $`Q_1`$ giving a residual global symmetry $`U_1^{}(1)=`$ $`\frac{1}{N+4}(2Q_1Q_{(N+4)})`$. The breakdown pattern thus is
$$\left[SU(N)\right]\times SU(N+4)\times U_1(1)SU^{}(N)\times SU(4)\times U_1^{}(1).$$
(11)
The gauge bosons have become massive as have some fermions. The fermionic spectrum, with respect to the residual global symmetry is
| | | $`SU^{}(N)`$ | $`SU(4)`$ | $`U_1^{}(1)`$ |
| --- | --- | --- | --- | --- |
| | $`S`$ | | $`1`$ | $`2`$ |
| massive | | | | |
| | $`\overline{F}^S`$ | $`\overline{\begin{array}{cc}& \end{array}}`$ | $`1`$ | $`2`$ |
| | $`\overline{F}^A`$ | $`\overline{\begin{array}{c}\\ \end{array}}`$ | $`1`$ | $`2`$ |
| massless | | | | |
| | $`\overline{F}_{i>N}`$ | $`\overline{\begin{array}{c}\end{array}}`$ | $`\overline{\begin{array}{c}\end{array}}`$ | $`1`$ |
(12)
This breaking pattern gives $`N^2+8N`$ Goldstone bosons, $`N^21`$ of which are eaten by the gauge bosons. So only $`8N+1`$ remain as part of the massless spectrum along with the massless fermions. The global anomalies are again matched by this spectrum. Those associated with the unbroken group $`SU^{}(N)\times SU(4)\times U_1^{}(1)`$ are matched by the massless fermions, while those associated with the broken global generators are matched by the Goldstone bosons. Since the Goldstone bosons do not couple singly to the massless fermions (no dimension-four operators), the effective zero-mass theory is free at low energies.
It follows that the thermodynamic free energy may be computed at $`T0`$ by counting the degrees of freedom. The result is
$$f_{IR}^{Higgs}(p=0)=(8N+1)+\frac{7}{4}[\frac{1}{2}N(N1)+4N],$$
(13)
where the superscript indicates that the gauge symmetry is (partially) broken. Just as in the case of the symmetric phase, the inequality Eq. (3) is satisfied: $`f_{IR}^{Higgs}(p=0)<f_{UV}(p=0)`$.
As an aside, we note that according to the idea of complementarity this low energy phase may be thought of as having arisen from confining gauge forces rather than the Higgs mechanism . Confinement then would partially break the global symmetry to the above group forming the necessary Goldstone bosons. It would also produce gauge singlet massless composite fermions to replace precisely the massless elementary fermions in the above table.
We have identified two possible phases of this theory consistent with global anomaly matching and the inequality Eq. (3). One confines and breaks no symmetries. The other breaks the chiral symmetry according to Eq. (11). For any finite value of $`N`$, $`f_{IR}^{sym}(p=0)<f_{IR}^{Higgs}(p=0)`$. The symmetric phase is thus favored if the number of degrees of freedom, or the entropy of the system near freeze-out, is minimized. In the limit $`N\mathrm{}`$, the Goldstone bosons do not contribute to leading order, and $`f_{IR}^{sym}(p=0)f_{IR}^{Higgs}(p=0)`$. We return to a discussion of the infinite $`N`$ limit after describing the general ($`p>0`$) model.
What about other symmetry breaking phases of the $`p=0`$ theory corresponding to bilinear condensate formation? In addition to $`S\overline{F}`$ condensates, there are also $`SS`$ and $`\overline{F}\overline{F}`$ possibilities. Several of these correspond to attractive channels, although not maximally attractive, due to gluon exchange. We have considered all of them for the case $`N=3`$, and have shown that the effective low energy theory is infrared free and that the number of low energy degrees of freedom is larger than the symmetric phase. Whether this is true for symmetry breaking patterns corresponding to bilinear condensate formation for general $`N`$ remains to be seen. Higher dimensional condensate formation is yet to be studied for any of these choices.
### B The General Case
We next consider the full range of $`p`$ allowed by asymptotic freedom: $`0<p<(9/2)N3`$. For $`p`$ near $`(9/2)N3`$, an infrared stable fixed point exists, determined by the first two terms in the $`\beta `$ function. This can be arranged by taking both $`N`$ and $`p`$ to infinity with the difference $`(9/2)Np`$ fixed, or at finite $`N`$ by continuing to nonintegral $`p`$. The infrared coupling is then weak and the theory neither confines nor breaks symmetries. The fixed point leads to an approximate, long-range conformal symmetry. As $`p`$ is reduced, the screening of the long range force decreases, the coupling increases, and confinement and/or symmetry breaking set in. We consider three strong-coupling possibilities, each consistent with global anomaly matching.
#### 1 Confinement with no symmetry Breaking
It was observed by Bars and Yankielowicz that confinement without chiral symmetry breaking is consistent with global anomaly matching provided that the spectrum of the theory consists of massless composite fermions transforming under the global symmetry group as follows:
(14)
The effective low-energy theory is free. In Ref. , the thermodynamic free energy for this phase was computed, giving
$$f_{IR}^{sym}=\frac{7}{4}[\frac{1}{2}(N+4+p)(N+3+p)+p(N+4+p)+\frac{1}{2}p(p+1)].$$
(15)
The inequality $`f_{IR}^{sym}<f_{UV}`$ was then invoked to argue that this phase is possible only if $`p`$ is less than a certain value (less than the asymptotic freedom bound). For large $`N`$, the condition is $`p<(15/14)^{1/2}N`$.
#### 2 Chiral symmetry breaking
Since this theory is vector-like with respect to the $`p`$ $`F`$-$`\overline{F}`$ pairs, it may be anticipated that these pairs condense according to
$$\overline{\begin{array}{c}\end{array}}\times \begin{array}{c}\end{array}1,$$
(16)
leading to a partial breaking of the chiral symmetries. The gauge-singlet bilinear condensate (fermion mass) is of the form
$$\epsilon ^{\gamma \delta }F_\gamma ^{a,i}\overline{F}_{a,N+4+i,\delta },$$
(17)
where $`i=1,\mathrm{},p`$.
This leads to the symmetry breaking pattern $`SU(N+4+p)\times SU(p)\times U_1(1)\times U_2(1)SU(N+4)\times SU_V(p)\times U_1^{}(1)\times U_2^{}(1)`$, producing $`2pN+p^2+8p`$ gauge singlet Goldstone bosons. The $`U^{}(1)^{}s`$ are combinations of the $`U(1)^{}s`$ and the broken generator of $`SU(N+4+p)`$
Q(N+4+p)=(
-p
pN+4N+4).subscript𝑄𝑁4𝑝fragments
-p
missing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionfragmentspmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionfragmentsN4missing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionfragmentsN4Q_{(N+4+p)}=\left(\begin{tabular}[]{ccc|ccc}$-p$&&&&&\\
&$\ddots$&&&&\\
&&$-p$&&&\\
\hline\cr&&&$N+4$&&\\
&&&&$\ddots$&\\
&&&&&$N+4$\end{tabular}\right). (18)
At this stage, the remaining massless theory is the $`p=0`$ theory described above, together with the $`2pN+p^2+8p`$ gauge-singlet Goldstone bosons. Since the Goldstone bosons are associated with the broken symmetry, there will be no dimension-four interactions between them and the $`p=0`$ theory. This theory may therefore be analyzed at low energies by itself, leading to the possible phases described above. Two possible phases of the $`p=0`$ theory were discussed in detail. One corresponds to confinement and massless composite fermion formation with no chiral symmetry breaking. For the general theory, this corresponds to
* Partial chiral symmetry breaking but no gauge symmetry breaking. The vector-like $`p`$ pairs of $`F`$ and $`\overline{F}`$ condense, and others form composite fermions.
The massless spectrum consists of the $`2pN+p^2+8p`$ Goldstone bosons together with the $`(N+4)(N+3)/2`$ composite fermions of the $`p=0`$ sector. All are confined. The final global symmetry is $`SU(N+4)\times SU_V(p)\times U_1^{}(1)\times U_2^{}(1)`$. Global anomalies are matched partially by the massless composites and partially by the Goldstone bosons. Since both theories are infrared free, the free energy may be computed in the $`T0`$ limit to give
$$f_{IR}^{brk+sym}=(2pN+p^2+8p)+\frac{7}{4}[\frac{1}{2}(N+4)(N+3)].$$
(19)
The inequality Eq. (3) thus allows this phase for $`p`$ less than a certain value below the asymptotic freedom bound but above the value at which the symmetric phase becomes possible. For large N, the limit is $`p/N`$ less than $`2.83`$.
The other phase of the $`p=0`$ theory considered above, corresponds to the MAC for symmetry breaking and the Higgsing of the gauge group with a further breaking of the chiral symmetry. For the general theory ($`p>0`$), it leads to
* Further chiral symmetry breaking and gauge symmetry breaking.
The final global symmetry is $`SU^{}(N)\times SU(4)\times SU_V(p)\times U_1^{}(1)\times U_2^{}(1)`$. The massless spectrum consists of the $`2pN+p^2+8p`$ Goldstone bosons associated with the $`p`$ $`F`$-$`\overline{F}`$ pairs, together with the $`8N+1`$ Goldstone bosons and $`N(N+1)/2+4N`$ massless elementary fermions of the $`p=0`$ sector. Global anomalies are matched partially by Goldstone bosons and partially by the remaining massless fermions. The effective low energy theories are infrared free, and we have
$`f_{IR}^{brk+Higgs}`$ $`=`$ $`(2pN+p^2+8p)+(8N+1)`$ (21)
$`+{\displaystyle \frac{7}{4}}[{\displaystyle \frac{1}{2}}N(N1)+4N].`$
Three possible phases of the general Bars-Yankielowicz model have now been identified. In Fig. 1, we summarize the computation of $`f_{IR}/N^2`$ for each phase and compare with $`f_{UV}`$ for the choice $`N=3`$. Other choices are qualitatively the same. Each phase satisfies the inequality Eq. (3) for $`p/N`$ small enough. As $`p`$ is reduced, the first phase allowed by the inequality corresponds to confinement with condensation of the $`p`$ fermions in the real representation of the gauge group and the breaking of the associated chiral symmetry, along with unbroken chiral symmetry and massless composite fermion formation in the $`p=0`$ sector. The degree of freedom count is denoted by $`f_{IR}^{brk+sym}`$.
The two other phases are also allowed by the inequality as $`p`$ is reduced further. But for any finite value of $`N`$ and for any value of $`p>0`$, the curve for $`f_{IR}^{brk+sym}`$ is the lowest of the $`f_{IR}`$ curves. Thus the lowest infrared degree-of-freedom count corresponds to a complete breaking of the chiral symmetry associated with the $`p`$ $`F`$-$`\overline{F}`$ pairs (the vector- like part of the theory with the fermions in a real representation of the gauge group), and no breaking of the chiral symmetry associated with the $`p=0`$ sector (the part of the theory with the fermions in complex representations).
It is instructive to examine this model in the infinite $`N`$ limit. If the limit is taken with $`p/N`$ fixed, the curves for $`f_{IR}^{brk+sym}/N^2`$ and $`f_{IR}^{brk+Higgs}/N^2`$ become degenerate for all values of $`p/N`$, and are below the curve for $`f_{IR}^{sym}`$. If the limit $`N\mathrm{}`$ is taken with $`p`$ fixed, all the curves become degenerate, and the phases are not distinguished by the number of degrees of freedom. The authors of Ref. analyzed the model in the $`N\mathrm{}`$ limit with confinement assumed and noted that the $`U_1(1)`$ symmetry cannot break because no appropriate order parameter can form in this limit. This is consistent with the above discussion since each of the phases preserves the $`U_1(1)`$ for any $`N`$.
As with the $`p=0`$ theory, there are other possible symmetry breaking phases corresponding to bilinear condensate formation. Some of these are attractive channels, although not maximally attractive, due to gluon exchange. We have considered several possibilities. Each leads to an effective low energy theory that is infrared free, and each gives a larger value of $`f_{IR}`$ than the phase corresponding to the lowest curve in Fig. 1: complete breaking of the chiral symmetry associated with $`p`$ additional $`F`$-$`\overline{F}`$ pairs and no breaking of the chiral symmetry associated with the sector of the theory with the fermions in complex representations.
Symmetry breaking by higher dimensional condensate formation is yet to be considered, and we have nothing to say about possible strongly coupled infrared phases such as a strongly coupled nonabelian Coulomb phase.
## III The Generalized Georgi-Glashow (GGG) Model
This model is similar to the BY model just considered. It is an $`SU(N5)`$ gauge theory, but with fermions in the anti-symmetric, rather than symmetric, tensor representation. The complete fermion content is $`A=\psi _L^{[ab]},`$ $`a,b=1,\mathrm{},N`$; an additional $`N4+p`$ fermions in the conjugate fundamental representations: $`\overline{F}_{a,i}=\psi _{a,iL}^c,`$ $`i=1,\mathrm{},N4+p`$; and $`p`$ fermions in the fundamental representations, $`F^{a,i}=\psi _L^{a,i}`$, $`i=1,\mathrm{},p`$.
The global symmetry is
$$G_f=SU(N4+p)\times SU(p)\times U_1(1)\times U_2(1).$$
(22)
where the two $`U(1)`$’s are anomaly free. With respect to this symmetry, the fermion content is
(23)
For the $`\beta `$-function, we have $`\beta _1=3N+2(2/3)p`$ and $`\beta _2=(1/4)\{13N^2+30N+1+12/N2p((13/3)N1/N)\}`$. Thus the theory is asymptotically free if
$$p<(9/2)N+3.$$
(24)
We restrict $`p`$ so that this condition is satisfied. Because of asymptotic freedom, the thermodynamic free-energy may be computed in the $`T\mathrm{}`$ limit. We have
$$f_{UV}=2(N^21)+\frac{7}{4}[\frac{N(N1)}{2}+(N4)N+2pN].$$
(25)
As with the BY model, we first discuss the $`p=0`$ theory and then consider the general case.
### A The $`p=0`$ Case
The global symmetry group is $`G_f=SU(N4)\times U_1(1)`$. The theory is strongly coupled at low energies, so it is expected either to confine or to break some of the symmetries, consistent with global anomaly matching .
In the case of complete confinement and unbroken symmetry, to satisfy global anomaly matching the massless spectrum consists of gauge singlet composite fermions $`\overline{F}_{\{i}A\overline{F}_{j\}}`$ transforming according to the symmetric second-rank tensor representation of $`SU(N4)`$ with charge $`N`$ under the $`U_1(1)`$ global symmetry . The composites are noninteracting in the infrared. Therefore the thermodynamic free energy may be computed in the limit $`T0`$. Enumerating the degrees of freedom gives
$$f_{IR}^{sym}(p=0)=\frac{7}{4}\frac{(N4)(N3)}{2}.$$
(26)
Clearly $`f_{IR}^{sym}(p=0)<f_{UV}(p=0)`$, satisfying the inequality Eq. (3).
We next consider symmetry breaking due to bilinear condensate formation by first examining the maximally attractive channel :
$$\begin{array}{c}\\ \end{array}\times \overline{\begin{array}{c}\end{array}}\begin{array}{c}\end{array},$$
(27)
leading to the formation of the $`A\overline{F}`$ condensate
$$\epsilon ^{\gamma \delta }A_\gamma ^{ai}\overline{F}_{a,i,\delta },$$
(28)
where $`\gamma ,\delta =1,2`$ are spin indices, $`a=1,\mathrm{},N`$, is a gauge index and $`i=1,\mathrm{},N4`$ is a flavor index. This condensate breaks the $`U_1(1)`$ symmetry and breaks the gauge symmetry $`SU(N)`$ to $`SU(4)`$. The broken gauge subgroup $`SU(N4)`$ combines with the flavor group, leading to a new global symmetry $`SU^{}(N4)`$, while the broken gauge $`SU(N)`$ generator
$`Q_{(N)}=\left(\begin{array}{c}\text{4}\text{}\text{}\text{}\text{}\text{}\\ & \mathrm{}& & & & \\ & & 4& & & \\ & & & 4N& & \\ & & & & \mathrm{}& \\ & & & & & 4N\end{array}\right),`$
combines with $`U_1(1)`$ to form a residual global symmetry $`U^{}(1)`$. The remaining symmetry is thus $`[SU(4)]\times SU^{}(N4)\times U_1^{}(1)`$. All Goldstone bosons are eaten by gauge bosons.
We have
$$\overline{F}_{a,i}=\left(\begin{array}{c}\overline{F}_{j,i}\overline{F}_{[j,i]}+\overline{F}_{\{j,i\}}\\ \overline{F}_{c,i}\end{array}\right)$$
(29)
and
$$A^{ab}=\left(\begin{array}{cc}A^{ij}& A^{ic}\\ & A^{cd}\end{array}\right),$$
(30)
where $`a,b=1,\mathrm{},N`$, $`i,j=1,\mathrm{},N4`$, and $`c,d=N3,\mathrm{},N`$. The $`A\overline{F}`$ condensate pairs $`\overline{F}_{[j,i]}`$ with $`A^{ij}`$ and $`\overline{F}_{c,i}`$ with $`A^{ic}`$. This leaves only $`A^{cd}`$, which is neutral under $`U^{}(1)`$, as the fermion content of the $`SU(4)`$ gauge theory.
This $`SU(4)`$ theory is also strongly coupled in the infrared and we expect it to confine. The most attractive channel for condensate formation, for example, is
$$\begin{array}{c}\\ \end{array}\times \begin{array}{c}\\ \end{array}\begin{array}{c}\text{}\\ \\ \\ \end{array}=1,$$
(31)
leading to the bilinear condensate
$$\epsilon ^{\gamma \delta }A_\gamma ^{ab}A_\delta ^{cd}\epsilon _{1\mathrm{}(N4)abcd},$$
(32)
a singlet under the gauge group. Thus, in the infrared, the only massless fermions are the $`\overline{F}_{\{j,i\}}^{}s`$ in the symmetric two-index tensor representation of $`SU^{}(N4)`$. Interestingly, the massless fermion content and the low energy global symmetry are precisely the same for the symmetric and Higgs phases. Therefore,
$$f_{IR}^{higgs}(p=0)=f_{IR}^{sym}(p=0)=\frac{7}{4}[\frac{1}{2}(N4)(N3)].$$
(33)
The fermions are composite in the first case and elementary in the second. This is another example of the complementarity idea . While the two phases are not distinguished by the low energy considerations used here, they are different phases. However, other ideas involving energies on the order of the confinement and/or breaking scales will have to be employed to distinguish them.
A general study of the phases of chiral gauge theories should include higher dimensional as well as bilinear condensate formation. We have done this for one case, the $`p=0`$ $`SU(N=5)`$ model, which possesses only a $`U(1)`$ global symmetry. Among the various phases that may be considered is one that confines but breaks the global $`U(1)`$. This corresponds to the formation of gauge invariant higher dimensional condensates, for example of the type $`\left(\overline{F}A\overline{F}\right)^2`$. There is no bilinear condensate for this breaking pattern. Global anomaly matching is satisfied by the appearance of a single massless Goldstone boson and no other massless degrees of freedom. This phase clearly minimizes the degree of freedom count (the entropy near freeze-out), among the phases described by infrared free effective theories. The unbroken phase, by contrast, must include a massless composite fermion for anomaly matching, and therefore gives a larger $`f_{IR}`$. This suggests that higher dimensional condensate formation may indeed be preferred in this model. It will be interesting to study this possibility in more detail and to see whether higher dimensional condensate formation plays an important role in the larger class of chiral theories considered here and in other theories.
### B The General Case
The full range of $`p`$ allowed by asymptotic freedom may be considered just as it was for the BY model. For $`p`$ near $`(9/2)N+3`$, an infrared stable fixed point exists, determined by the first two terms in the $`\beta `$ function. The infrared coupling is then weak and the theory neither confines nor breaks symmetries. As $`p`$ decreases, the coupling strengthens, and confinement and/or symmetry breaking set in. We consider two possibilities consistent with global anomaly matching.
#### 1 Confinement with no symmetry breaking
It is known that confinement without chiral symmetry breaking is consistent with global anomaly matching provided that the spectrum of the theory consists of gauge singlet massless composite fermions transforming under the global symmetry group as follows:
(34)
The effective low energy is free. Thus the thermodynamic free energy may be computed in the limit $`T0`$ to give
$$f_{IR}^{sym}=\frac{7}{4}[\frac{1}{2}(N4+p)(N3+p)+p(N4+p)+\frac{1}{2}p(p1)].$$
(35)
The inequality Eq. (3) allows this phase when $`p/N`$ is less than $`2.83`$, for large $`N`$.
#### 2 Chiral symmetry breaking
As in the BY model it may be expected that the fermions in a real representation of the gauge group (the $`p`$ $`F`$-$`\overline{F}`$ pairs) will condense in the pattern
$$\overline{\begin{array}{c}\end{array}}\times \begin{array}{c}\end{array}1.$$
(36)
The gauge-singlet bilinear condensate (fermion mass) is of the form
$$\epsilon ^{\gamma \delta }F_\gamma ^{a,i}\overline{F}_{a,N4+i,\delta },$$
(37)
where $`i=1,\mathrm{},p`$, leading to the symmetry breaking pattern
| $`SU(N4+p)\times SU(p)\times U_1(1)\times U_2(1)`$ |
| --- |
| $`SU(N4)\times SU_V(p)\times U_1^{}(1)\times U_2^{}(1),`$ |
(38)
and producing $`2pN+p^28p`$ gauge singlet Goldstone bosons.
The $`U^{}(1)^{}s`$ are combinations of the $`U(1)^{}s`$ and the broken generator of $`SU(N4+p)`$
Q(N4+p)=(
-p
pN4N4).subscript𝑄𝑁4𝑝fragments
-p
missing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionfragmentspmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionfragmentsN4missing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionfragmentsN4Q_{(N-4+p)}=\left(\begin{tabular}[]{ccc|ccc}$-p$&&&&&\\
&$\ddots$&&&&\\
&&$-p$&&&\\
\hline\cr&&&$N-4$&&\\
&&&&$\ddots$&\\
&&&&&$N-4$\end{tabular}\right)\ . (39)
The remaining massless theory is the $`p=0`$ theory described above, together with the $`2pN+p^28p`$ gauge-singlet Goldstone bosons. Since the Goldstone bosons are associated with the broken symmetry, there will be no dimension-four (Yukawa) interactions between them and the $`p=0`$ fields. The $`p=0`$ theory may therefore be analyzed by itself, leading to the possible phases described above. Two phases were considered, one symmetric and the other broken by the maximally attractive bilinear condensate, and they were seen to lead to identical low energy theories.
Thus, in either case, the degree-of-freedom count for the general theory, corresponding to the breaking of the chiral symmetry associated with the $`p`$ $`F`$-$`\overline{F}`$ pairs, gives
$$f_{IR}^{brk}=(2pN+p^28p)+\frac{7}{4}[\frac{1}{2}(N4)(N3)].$$
(40)
To summarize, two possible phases of the general GGG model have been considered. In Fig. 2, we plot the two computations of $`f_{IR}/N^2`$ for the choice $`N=6`$ and compare them with $`f_{UV}`$. Each phase satisfies the inequality Eq. (3) for $`p`$ below some value. As $`p`$ is reduced, the first phase allowed by the inequality corresponds to partial chiral symmetry breaking. For any $`p`$, $`f_{IR}^{brk}`$ is the lower of the $`f_{IR}`$ curves. Thus the lower infrared degree-of-freedom count corresponds to a complete breaking of the chiral symmetry associated with $`p`$ additional $`F`$-$`\overline{F}`$ pairs (the vector like part of the theory) and no breaking of the chiral symmetry associated with the fermions in a complex representation of the gauge group. Whether the latter behavior is due to confinement or the Higgsing of the gauge group has not been determined. These conclusions remain valid in the infinite $`N`$ limit with $`p/N`$ fixed. If the limit $`N\mathrm{}`$, is taken with $`p`$ fixed, the two curves become degenerate.
As we have already noted, for the $`p=0`$ $`SU(5)`$ theory there is a still lower degree of freedom count when higher dimensional condensates are considered. This will therefore also be true of the general-$`p`$ case for $`SU(5)`$. This even lower count corresponds to complete confinement along with breaking of the chiral symmetry associated with the $`p`$ $`F`$-$`\overline{F}`$ pairs and breaking of the remaining global $`U(1)`$ symmetry. It will be interesting to see whether a preference for this phase can be confirmed by a dynamical study of this model and whether similar higher dimensional condensate formation is favored in a more general class of models.
## IV Two Chiral SUSY Models
Although this paper is devoted principally to non-SUSY chiral models, we briefly describe two chiral SUSY models: the supersymmetric generalization of the one generation $`SU(5)`$ Georgi-Glashow model and the related ($`32`$) model (see for a review of this model and relevant references).
The $`SU(5)`$ model contains a single antisymmetric tensor chiral superfield $`A`$ and an antifundamental chiral superfield $`\overline{F}`$. The vector superfield $`W_\alpha `$ includes the standard vector boson and the associated gluino in the adjoint representation of $`SU(5)`$. The global symmetry is the anomaly-free $`U_R(1)\times U_A(1)`$, and the charge assignments are:
| | $`[SU(5)]`$ | $`U_R(1)`$ | $`U_A(1)`$ |
| --- | --- | --- | --- |
| $`A`$ | | $`1`$ | $`1`$ |
| $`\overline{F}`$ | $`\overline{\begin{array}{c}\end{array}}`$ | $`+9`$ | $`+3`$ |
| $`W_\alpha `$ | $`Adj`$ | $`1`$ | $`0`$ |
(41)
A special feature of this model is that the classical vacuum is unique. The absence of flat directions is due to the fact that there exists no holomorphic gauge invariant polynomial constructed out of the supersymmetric fields. This feature guarantees that when comparing phases through their degree-of-freedom count, we know that we are considering a single underlying theory. By contrast, in SUSY gauge theories with flat directions, non-zero condensates associated with the breaking of global symmetries correspond to different points in moduli space and therefore to different theories.
This model was studied long ago and various possible phases were seen to be consistent with global anomaly matching. One preserves supersymmetry along with the global symmetries. This requires composite massless fermions to saturate the global anomalies. It was shown that there are several, rather complicated, solutions, with at least five Weyl fermions (which for supersymmetry to hold must be cast in five chiral superfields). The charge assignments for one of them is :
$$(5,26),(5,20),(5,24),(0,1),(0,9),$$
(42)
where the first entry is the $`U_A(1)`$ charge and the second is the $`U_R(1)`$ charge of each chiral superfield.
Other possibilities are that SUSY breaks with the global symmetries unbroken or that one or both of the global symmetries together with supersymmetry break spontaneously. It is expected that in a supersymmetric theory without classical flat directions, the spontaneous breaking of global symmetries also signals spontaneous supersymmetry breaking. In these cases, the only massless fields will be the Goldstone boson(s) associated with the broken global symmetries and/or some massless fermions transforming under the unbroken chiral symmetries, together with the Goldstone Weyl fermion associated with the spontaneous supersymmetry breaking.
In Reference it was suggested on esthetic grounds that the supersymmetric solution seems less plausible. Additional arguments that supersymmetry is broken are based on investigating correlators in an instanton background . However a firm solution to this question is not yet available.
Since all the above phases are non interacting in the infrared we may reliably compute $`f_{IR}`$ and note that the phase that minimizes the degree-of-freedom count is the one that breaks supersymmetry and both of the global symmetries. This phase consists of two $`U(1)`$ Goldstone bosons and a single Weyl Goldstino associated with the breaking of SUSY. Thus $`f_{IR}=15/4`$. SUSY preserving phases and those that leave one or both of the $`U(1)^{}s`$ unbroken lead to more degrees of freedom. It will be interesting to see whether further dynamical studies confirm that the maximally broken phase is indeed preferred
This phase is similar to the minimal-$`f_{IR}`$ phase in the nonsupersymmetric $`SU(5)`$ model in that both correspond to higher dimensional condensate formation. In the SUSY case, one can construct two independent order parameters. The one for $`U_R(1)`$ is the gluino condensate (scalar component of the chiral superfield $`W^\alpha W_\alpha `$) while the one for $`U_A(1)`$ can be taken to be the scalar component of the chiral superfield $`\overline{F}_a\overline{F}_bA^{ac}(W^\alpha W_\alpha )_c^b`$.
Finally we comment on a well known and related chiral model for dynamical supersymmetry breaking: the ($`32`$) model. Unlike the models considered so far, this model involves multiple couplings, i.e. two gauge couplings and a Yukawa one. Without the Yukawa interaction the theory posses a run-away vacuum. The model has an $`SU(3)\times SU(2)`$ gauge symmetry and a $`U_Y(1)\times U_R(1)`$ anomaly free global symmetry. As above, the low energy phase that minimizes the number of degrees of freedom is the one that breaks supersymmetry along with both of the global symmetries. The massless spectrum is the same as in the parent chiral $`SU(5)`$ case. In the ($`32`$) model, however, the low energy spectrum has been computed in a self-consistent weak-Yukawa coupling approximation, where it was noted that the $`U_R(1)`$ breaks along with supersymmetry, leaving intact the $`U_Y(1)`$. The spectrum consists of two massless fermions (a Goldstino and the fermion associated with the unbroken $`U_Y(1)`$) and the $`U_R(1)`$ Goldstone boson. If this is indeed the ground state, then the number of infrared degrees of freedom is not minimized in this weak coupling case.
## V Conclusions
We have considered the low energy structure of two chiral gauge theories, the Bars-Yankielowicz (BY) model and the generalized Georgi-Glashow (GGG) model. Each contains a core of fermions in complex representation of the gauge group, along with a set of $`p`$ additional fundamental-anti-fundamental pairs. In each case, for $`p`$ near but not above the value for which asymptotic freedom is lost, the model will have a weak infrared fixed point and exist in the non-abelian Coulomb phase.
As $`p`$ drops, the infrared coupling strengthens and one or more phase transitions to strongly coupled phases are expected. Several possible phases have been identified that are consistent with global anomaly matching, and that satisfy the inequality Eq. (3) for low enough $`p`$. One is confinement with the gauge symmetry and additional global symmetries unbroken. Another is confinement with the global symmetry broken to that of the $`p=0`$ theory. Still another is a Higgs phase, with both gauge and chiral symmetries broken. Both symmetry breaking phases correspond to bilinear condensate formation. The infrared degree of freedom count $`f_{IR}`$ for each of these phases is shown in Figs. 1 and 2, along with the corresponding ultraviolet count $`f_{UV}`$.
We have suggested that at each value of $`p`$, these theories will choose the phase that minimizes the degree of freedom count as defined by $`f_{IR}`$, or equivalently the phase that minimizes the entropy near freeze-out ($`S(T)(2\pi ^2/45)T^3f_{IR}`$). As may be seen from Figs. 1 and 2, this idea leads to the following picture. As $`p`$ drops below some critical value, the $`p`$ fundamental-anti-fundamental pairs condense at some scale $`\mathrm{\Lambda }`$, breaking the full global symmetry to the symmetry of the $`p=0`$ theory and producing the associated Goldstone bosons. For the remaining theory with fermions in only complex representations, the phase with the global symmetry unbroken and the global anomalies matched by massless fermions is preferred to phases with further global symmetry breaking via bilinear condensates. We have not yet shown that this is true relative to all bilinear condensate formation. Also, this does not exclude the possibility that some strongly coupled infrared phase (such as a strong non abelian coulomb phase) leads to the smallest value for $`f_{IR}`$ and is still consistent with global anomaly matching.
We extended our discussion to include general condensate formation for one simple example, the $`SU(5)`$ Georgi-Glashow model with fermions in only complex representations and a single $`U(1)`$ global symmetry. This symmetry can be broken via only a higher dimensional condensate. For this model, interestingly, we noted that the breaking of the $`U(1)`$ with confinement and unbroken gauge symmetry leads to the minimum value of $`f_{IR}`$ among phases that are infrared free. This highlights the important question of the pattern of symmetry breaking in general chiral theories (or any theories for that matter) when arbitrary condensate formation is considered. Higher dimensional condensates could play an important role in the dynamical breaking of symmetries in extensions of the standard model . The enumeration of degrees of freedom in the effective infrared theory is a potentially useful guide to discriminate among the possibilities.
Finally, we commented on two supersymmetric chiral models: the supersymmetric version of the $`SU(5)`$ Georgi-Glashow model and the closely related ($`32`$) model. Both have a $`U_R(1)\times U_Y(1)`$ global symmetry. In each case, the phase that minimizes the number of massless degrees of freedom corresponds to the breaking of SUSY and both of its global symmetries. In the case of the ($`32`$) model, however, an analysis in the case of a weak Yukawa coupling (see for a discussion and relevant references) leads to the conclusion that the $`U_Y(1)`$ is not broken. If this truly represents the ground state in the case of weak coupling, then the degree of freedom count is not the minimum among possible phases that respect global anomaly matching.
To summarize, for the nonsupersymmetric chiral gauge theories discussed here, we have identified a variety of possible zero-temperature phases and conjectured that the theories will choose from among them the one that minimizes the infrared degree of freedom count. Whether this can be proven and whether the idea plays a role in a wider class of theories remains to be seen.
Acknowledgments
We thank Andrew Cohen, Erich Poppitz, Nicholas Read and Robert Shrock for helpful discussions. The work of T.A., Z.D. and F.S. has been partially supported by the US DOE under contract DE-FG-02-92ER-40704. |
warning/0001/hep-th0001106.html | ar5iv | text | # References
hep-th/0001106
AEI-2000-002
On the supersymmetric effective action of Matrix theory
Hermann Nicolai and Jan Plefka
Max-Planck-Institut für Gravitationsphysik
Albert-Einstein-Institut
Am Mühlenberg 1, D-14476 Golm, Germany
Abstract
We present a simple derivation of the supersymmetric one-loop effective action of SU(2) Matrix theory by expressing it in a compact exponential form whose invariance under supersymmetry transformations is obvious. This result clarifies the one-loop exactness of the leading $`(v^2)^2`$ interactions and the absence of non-perturbative corrections.
January 2000
Recently maximally supersymmetric $`SU(N)`$ gauge quantum mechanics in $`d=9`$ has gained prominence due to its relation to the low-energy dynamics of zerobranes in type IIA string theory , the close relation between its $`N\mathrm{}`$ limit and the eleven-dimensional supermembrane , as well as the M theory proposal of . A key feature of this model is the existence of flat directions in the Cartan sector on which scattering states localize. To date almost all investigations of scattering amplitudes in Matrix theory make use of the perturbative construction of an effective Lagrangian for the Cartan valley degrees of freedom at finite $`N`$, which is based on a loopwise expansion around the solution of the classical equations of motions, $`\ddot{x}_I^i=0`$ and $`\dot{\theta }_I=0`$, where $`I,J,\mathrm{}=1,\mathrm{},N1`$. Although this approach simply ignores contributions from bound states all tree level amplitudes computed to date in matrix theory agree with the results obtained from eleven dimensional supergravity . As soon as one goes beyond the tree level regime, however, this correspondence breaks down<sup>1</sup><sup>1</sup>1In fact, it can be shown that there is no Lorentz invariant combination of $`R^4`$ terms that reproduces the matrix theory result.. As argued in the agreement of the $`(v^2)^2`$ and $`(v^2)^3`$ terms in the effective action with tree level supergravity could be solely due to the high amount of supersymmetry in the problem. In particular in it was shown that the leading corrections to the $`SU(2)`$ effective action of order $`(v^2)^2`$ are completely determined by supersymmetry, a claim thereafter made explicit by . In this note we wish to demonstrate how this rather involved analysis may be condensed to a two line argument, yielding the complete form of the supersymmetric $`SU(2)`$ effective action at order $`(v^2)^2`$.
The basic variables for the $`SU(2)`$ theory in the Cartan subalgebra are
$$X^i(t),v^i(t):=\dot{X}^i(t)\mathrm{and}\theta _\alpha (t)$$
(1)
where $`i,j,\mathrm{}=1,\mathrm{},9`$ and $`\alpha ,\beta =1,\mathrm{},16`$. These variables correspond to the diagonal degrees of freedom in the matrix theory; for $`SU(N)`$ we would have $`X_I^i`$ with $`I=1,\mathrm{},(N1)`$.
The supersymmetry variations are given by:
$$\delta X^i=iϵ\gamma ^i\theta +ϵN^i\theta \delta \theta =v^j\gamma _jϵ+Mϵ$$
(2)
where $`N^i`$ and $`M`$ correspond to higher order modifications. For $`N^i=M=0`$, these variations leave invariant the free action
$$S^{(0)}=𝑑t\left[\frac{1}{2}v^2+\frac{i}{2}\theta \dot{\theta }\right]$$
(3)
Besides these terms, the effective action will contain an infinite string of higher order corrections. Since the algebra closes only on-shell, the supersymmetry variations must be modified accordingly such that $`N^i`$ and $`M`$ will no longer vanish. In considering such corrections, one must also take into account nonlinear field redefinitions
$$X^iX^i=X^i(X,v,\theta )\theta _\alpha \theta _\alpha ^{}=\theta _\alpha ^{}(X,v,\theta )$$
(4)
Modifications of the supersymmetry variations induced by such redefinitions do preserve the algebra, but should be discarded as they do not correspond to genuine deformations of the original variations.
Remarkably, even in this simple quantum mechanical context, no nontrivial modifications with $`N^i,M0`$ have so far been explicitly exhibited in the literature, although in evidence for non-trivial $`N^i\theta ^4`$ and $`M\theta ^6`$ modifications was presented. A full treatment is difficult because a complete analysis of the superalgebra and its closure will presumably require the consideration of infinitely many corrections. The problem is aggravated by the fact that for the maximally extended models no off-shell formulation is known<sup>2</sup><sup>2</sup>2Besides, it is doubtful whether an off-shell formalism would be of much use here, as the “rules of the game” are no longer clear: the elimination of auxiliary fields via their equations of motion and via the path integral yield inequivalent results unless the auxiliary fields appear at most quadratically in the Lagrangian..
To simplify matters, one makes the assumption that
$$\frac{dv^i}{dt}=0,\frac{d\theta }{dt}=0$$
(5)
This assumption, which implicitly also underlies all previous work, allows us to drop all derivatives other than those of $`X^i`$ in the variations, and greatly simplifies the calculation; for instance, we can then consistently set $`\delta v^i=0`$ for the unmodified variations. Effectively, the above condition amounts to a reduction of a quantum mechanical system to a “zero-dimensional” system. The freedom of making field redefinitions is reduced accordingly: the only redefinitions compatible with the above reduction are the ones preserving the linear dependence of $`X^i`$ on $`t`$ and the constancy of $`v^i`$ and $`\theta `$.
In the full supersymmetric one-loop effective action was shown to be
$`S^{(1)}`$ $`=`$ $`{\displaystyle }dt[(v^2)^2f(X)+{\displaystyle \frac{i}{2}}v^2v^j_jf(\theta \gamma ^{ij}\theta ){\displaystyle \frac{1}{8}}v^iv^j_k_lf(X)(\theta \gamma ^{ik}\theta )(\theta \gamma ^{jl}\theta )`$ (6)
$`{\displaystyle \frac{i}{144}}v^i_j_k_lf(X)(\theta \gamma ^{ij}\theta )(\theta \gamma ^{km}\theta )(\theta \gamma ^{lm}\theta )`$
$`+{\displaystyle \frac{1}{8064}}_i_j_k_lf(X)(\theta \gamma ^{im}\theta )(\theta \gamma ^{jm}\theta )(\theta \gamma ^{kn}\theta )(\theta \gamma ^{ln}\theta )]`$
Here $`f=f(X)`$ is a harmonic function, i.e. $`f_j_jf(X)=0`$ with the unique rotationally invariant solution $`f=r^7`$ (where $`r:=\sqrt{X^iX^i}`$). Provided one assumes constancy of $`v^i`$ and $`\theta `$ the action $`S^{(1)}`$ must be invariant under the unmodified supersymmetry variations above, as the modified variation of the free action $`S^{(0)}`$ under constant $`v^i`$ and $`\theta `$
$$\delta S^{(0)}=𝑑t_t(v^iϵN^i\theta )$$
(7)
vanishes for asymptotically suppressed corrections, $`lim_{t\pm \mathrm{}}N^i=0`$. Possible leading modifications of the supersymmetry transformations were discussed in , but clearly these do not contribute under the assumption (5).
We will now show that the complicated action $`S^{(1)}`$ can be cast into a much simpler form, whose supersymmetry invariance is very easy to check. Namely, we have
$`S^{(1)}`$ $`=`$ $`{\displaystyle 𝑑t(v^2)^2\mathrm{exp}\left[\frac{i}{2v^2}\theta \gamma ^{ij}\theta v_i_j\right]f(X)}`$ (8)
$`=`$ $`{\displaystyle 𝑑t(v^2)^2f\left(X\frac{i}{2v^2}\theta \gamma ^{ij}\theta v_j\right)}`$
To prove that this action indeed coincides with (6), we first observe that in the above action we can neglect all terms containing the Laplacian (which annihilates $`f`$) as well as terms containing $`v^j_j`$, because with constant $`v^i`$ and $`\theta `$, this term can be pulled out, yielding a total time derivative. To proceed, we show that
$$(\theta \gamma ^{ij}\theta v_i_j)(\theta \gamma ^{kl}\theta _l)(\theta \gamma ^{km}\theta _m)\frac{3}{v^2}(\theta \gamma ^{ij}\theta v_i_j)^3$$
(9)
and
$$\left((\theta \gamma ^{kl}\theta _l)(\theta \gamma ^{km}\theta _m)\right)^2\frac{21}{(v^2)^2}\left(\theta \gamma ^{ij}\theta v_i_j\right)^4$$
(10)
from which the equivalence follows up to fourth order. The symbol $``$ here and below means equality modulo contributions containing $`v^i_i`$ or $``$. To verify the above relations we start out from the Fierz identity (see e.g. for a comprehensive list of Fierz identites)
$$(\theta \gamma ^{ijk}\theta v_i_j)^25(\theta \gamma ^{ij}\theta v_i_j)^2+v^2(\theta \gamma ^{ij}\theta _j)(\theta \gamma ^{ik}\theta _k)$$
(11)
Thereafter one multiplies (11) with $`(\theta \gamma ^{ij}\theta v_i_j)`$ so that its left hand side may be rewritten as
$$(\theta \gamma ^{ij}\theta v_i_j)(\theta \gamma ^{klm}\theta v_k_l)^2\frac{1}{6}v^2(\theta \gamma ^{ij}\theta _i)(\theta \gamma ^{klj}\theta _k)(\theta \gamma ^{mnl}\theta v_m_n).$$
(12)
upon using yet another Fierz identity. Now once more perform a Fierz rearrangement on the last two terms of the above expression to obtain
$$(\theta \gamma ^{ij}\theta v_i_j)(\theta \gamma ^{klm}\theta v_k_l)^2\frac{2}{3}v^2(\theta \gamma ^{ij}\theta v_i_j)(\theta \gamma ^{kl}\theta _k)^2$$
(13)
This is to be contrasted with the right hand side of (11) multiplied with $`(\theta \gamma ^{ij}\theta v_i_j)`$:
$$(\theta \gamma ^{ij}\theta v_i_j)(\theta \gamma ^{klm}\theta v_k_l)^25(\theta \gamma ^{ij}\theta v_i_j)^3+v^2(\theta \gamma ^{ij}\theta v_i_j)(\theta \gamma ^{kl}\theta _k)^2$$
(14)
From (13) and (14) the relation (9) immediately follows. Relation (10) is then shown in a similar manner.
Next we note that the exponential series of (8) terminates already after the fourth order term because
$$(\theta \gamma ^{ij}\theta v_i_j)^50,$$
(15)
which is an immediate consequence of (9) and (10):
$`{\displaystyle \frac{21}{(v^2)^2}}\left(\theta \gamma ^{ij}\theta v_i_j\right)^5`$ $``$ $`(\theta \gamma ^{ij}\theta v_i_j)\left((\theta \gamma ^{kl}\theta _l)(\theta \gamma ^{km}\theta _m)\right)^2`$ (16)
$``$ $`{\displaystyle \frac{3}{v^2}}(\theta \gamma ^{ij}\theta v_i_j)^3(\theta \gamma ^{kl}\theta _l)(\theta \gamma ^{km}\theta _m)`$
$``$ $`{\displaystyle \frac{9}{(v^2)^2}}(\theta \gamma ^{ij}\theta v_i_j)^5`$
where we used (10) in the first and (9) in the second and third lines. Hence there is no need to truncate the series (it would anyhow terminate at order $`\theta ^{16}`$ by the nilpotency of the Grassmann variables).
The supersymmetry of this action with the above assumptions (5) (and $`N^i=M=0`$) can now be proven in two lines:
$`\delta S^{(1)}`$ $`=`$ $`{\displaystyle 𝑑t(v^2)^2\mathrm{exp}\left[\frac{i}{2v^2}\theta \gamma ^{ij}\theta v_i_j\right]\left(\frac{i}{v^2}\delta \theta \gamma ^{ij}\theta v_i_jf(X)+\delta X^i_if(X)\right)}`$ (17)
$`=`$ $`{\displaystyle 𝑑t(v^2)^2\mathrm{exp}\left[\frac{i}{2v^2}\theta \gamma ^{ij}\theta v_i_j\right]\frac{i}{v^2}ϵv/\theta v^i_if(X)}=0.`$
Remarkably, this simple argument works for any action of the form
$$𝑑t\mathrm{exp}\left[\frac{i}{2v^2}\theta \gamma ^{ij}\theta v_i_j\right]g(X,v)$$
(18)
and in particular yields supersymmetric completions of
$$g(r,v)=\frac{(v^2)^m}{r^n}$$
(19)
Uniqueness and therefore a “non-renormalization theorem” holds only for actions with low powers of $`v^2`$, and only if one insists on the absence of terms singular in $`v^2`$. For $`g(r,v)v^2`$, the only way to avoid such singular terms is to require $`_ig=0`$, in which case one is left with the free action only. For $`g(r,v)(v^2)^2`$, a singularity could arise at order $`\theta ^6`$, and is eliminated by means of the requirement $`g=0`$ (and the above Fierz identities implying (15)). Unfortunately for $`g(r,v)(v^2)^3`$ our above arguments fail, as the modified supersymmetry transformations of $`S^{(1)}`$ now do produce non-vanishing terms under the constancy assumption of $`v^i`$ and $`\theta `$. Hence even in this reduced sector (18) cannot be the full answer. At order $`g(r,v)(v^2)^4`$ and beyond, no singular terms can arise, and the choice of $`g(r,v)`$ is not restricted in any way. This is meant by the statement that, at this order there is no “non-renormalization theorem”.
We have thus seen that the one-loop effective action (6) resp. (8) laboriously computed in is indeed completely fixed by supersymmetry. Therefore the agreement of the resulting spin dependent Matrix theory scattering amplitudes with tree level supergravity does not test any dynamical aspects of the M theory proposal, but is solely due to the right amount of supersymmetry in the problem.
Acknowledgement
We wish to thank A. Waldron, who had also guessed (8) independently through an explicit S-Matrix computation, M. Staudacher and A. Tseytlin for useful comments. |
warning/0001/astro-ph0001379.html | ar5iv | text | # One-Point Probability Distribution Functions of Supersonic Turbulent Flows in Self-Gravitating Media
## 1 Introduction
Turbulence is an important ingredient for understanding the properties and characteristics of molecular clouds and star-forming regions. Turbulent gas motions are highly supersonic as indicated by the superthermal line widths ubiquitously observed throughout molecular clouds (Williams, Blitz & McKee 2000). These motions carry enough energy to halt global collapse and act as stabilizing agent for the entire cloud. However, it can be shown that interstellar turbulence decays quite rapidly on time scales of the order of the free-fall time of the system (Mac Low et al. 1998, Stone, Ostriker & Gammie 1998, Padoan & Nordlund 1999). To explain the observed long life times, turbulence in molecular clouds must be constantly driven (Gammie & Ostriker 1996, Mac Low 1999). The interplay between self-gravity on the one hand (leading to local collapse and star formation) and turbulent gas motion on the other hand (trying to prevent this process) plays a key role in determining the structure of molecular clouds. Altogether, understanding the characteristics of compressible, supersonic, and constantly replenished turbulence in self-gravitating media is an important ingredient for an adequate description of molecular clouds dynamics. And vice versa, from analyzing the spatial and dynamical structure of molecular clouds we can gain insight into the phenomenon of turbulence (for an overview over interstellar turbulence see Franco & Carraminana 1999).
Unfortunately, a complete and comprehensive theory of turbulence does not exist. Due to the enormous complexity of the problem, progress has been slow since Kolmogorov’s pioneering work in 1941, where he derived simple scaling laws for incompressible, stationary, and homogeneous turbulence by postulating a self-similar energy cascade downwards from the driving scale to the dissipation range. Most effort has since been put in finding an adequate closure procedure, i.e. in finding a way to express the highest-order correlation in the hierarchy of equations governing turbulent motion (for an excellent overview see Lesieur 1997; also Boratav, Eden & Erzan 1997). However, a satisfying description of turbulence has yet to be found.
Correlation and distribution functions of dynamical variables are frequently deployed for characterizing the kinematical properties of turbulent molecular clouds. Besides using 2-point statistics (e.g. Scalo 1984, Kleiner & Dickman 1987, Kitamura et al. 1993, Miesch & Bally 1994, LaRosa, Shore & Magnani 1999), many studies have hereby concentrated on 1-point statistics, namely on analyzing the probability distribution function (pdf) of the (column) density and of dynamical observables, e.g. of the centroid velocities of molecular lines and their increments. The density pdf has been used to characterize numerical simulations of the interstellar medium by Vázquez-Semadeni (1994), Padoan, Nordlund, & Jones (1997), Passot, & Vázquez-Semadeni (1998) and Scalo et al. (1998). Velocity pdf’s for several star-forming molecular clouds have been determined by Miesch & Scalo (1995) and Miesch, Scalo & Bally (1998). Lis et al. (1996, 1998) analyzed snapshots of a numerical simulation of mildly supersonic, decaying turbulence (without self-gravity) by Porter, Pouquet, & Woodward (1994) and applied the method to observations of the $`\rho `$-Ophiuchus cloud. Altogether, the observed pdf’s exhibit strong non-Gaussian features, they are often nearly exponential with possible evidence for power-law tails in the outer parts. This disagrees with the nearly Gaussian behavior typically found in experimental measurements and numerical models of incompressible turbulence. The observed centroid velocity increment pdf’s are more strongly peaked and show stronger deviations from Gaussianity than numerical models of incompressible turbulence predict. Furthermore, the spatial distribution of the largest centroid velocity differences (determining the tail of the distribution) appears ‘spotty’ across the face of the clouds; there is no convincing evidence for filamentary structure. Miesch et al. (1998) conclude that turbulence in molecular clouds involves physical processes that are not adequately described by incompressible turbulence or mildly supersonic decay simulations (see also Mac Low & Ossenkopf 2000).
It is the principal goal of this paper to extend previous determinations of pdf’s from numerical models into a regime more applicable for interstellar turbulence by (1) by calculating fully supersonic flows, (2) by including self-gravity, and (3) by incorporating a (simple analytic) description of turbulent energy input. For comparison with molecular cloud observations, I discuss the dynamical properties of decaying and stationary (i.e. driven), supersonic, isotropic turbulence in self-gravitating isothermal gaseous media. The pdf’s for the density, for the line centroid velocity and for their increments are derived as function of time and evolutionary state of the turbulent model.
The structure of this paper is as follows: Section 2 introduces and defines the statistical tools applied in the current study. It is followed in Sec. 3 by a description of the numerical scheme used to compute the time evolution of the turbulent flows. Sec. 4 shows that already simple variance effects in random Gaussian fields are able to introduce strong non-Gaussian distortions to the pdf’s which makes a clear-cut interpretation difficult. Section 5 contains the analysis of decaying, initially highly supersonic turbulence without self-gravity. This effect is then added to the simulations presented in Sec. 6. The model most relevant for molecular cloud dynamics is discussed in Sec. 7. It includes a simple driving term to replenish the turbulent cascade. Finally, in Sec. 8 all results are summarized.
## 2 PDF’s and Their Interpretation
### 2.1 Turbulence and PDF’s
The Kolmogorov (1941) approach to incompressible turbulence is a purely phenomenological one and assumes the existence of a stationary turbulent cascade. Energy is injected into the system at large scales and cascades down in a self-similar way. At the smallest scales it gets converted into heat by molecular viscosity. The flow at large scales is essentially inviscid, hence for small wave numbers the equation of motion is dominated by the advection term. If the stationary state of fully developed turbulence results from random external forcing then one naïvely expects the velocity distribution in the fluid to be Gaussian on time scales larger than the correlation time of the forcing, irrespectively of the statistics of the forcing term which follows from the central limit theorem. However, the situation is more complex (e.g. Frisch 1995, Lesieur 1997). One of the most striking (and least understood) features of turbulence is its intermittent spatial and temporal behavior. The structures that arise in a turbulent flow manifest themselves as high peaks at random places and at random times. This is reflected in the pdf’s of dynamical variables or passively advected scalars. They are sensitive measures of deviations from Gaussian statistics. Rare strong fluctuations are responsible for extended tails, whereas the much larger regions of low intensity contribute to the peak of the pdf near zero (for an analytical approach see e.g. Forster, Nelson & Stephens 1977, Falkovich & Lebedev 1997, Chertkov, Kolokolov & Vergassola 1997, Balkovsky et al. 1997, Balkovsky & Falkovich 1998). For incompressible turbulence the theory predicts velocity pdf’s which are mainly Gaussian with only minor enhancement at the far ends of the tails. The distribution of velocity differences (between locations in the system separated by a given shift vector $`\mathrm{\Delta }\stackrel{}{r}`$) is expected to deviate considerably from being normal and is likely to resemble an exponential. This finding is supported by a variety of experimental and numerical determinations (e.g. Kida & Murakami 1989, Vincent & Meneguzzi 1991, Jayesh & Warhaft 1991, She 1991, She, Jackson & Orszag 1991, Cao, Chen, & She 1996, Vainshtein 1997, Lamballais, Lesieur, & Métais 1997, Machiels & Deville 1998). Compressible turbulence has remained to be too complex for a satisfying mathematical analysis.
### 2.2 PDF’s of Observable Quantities
It is not clear how to relate the analytical work on incompressible turbulence to molecular clouds. In addition to the fact that interstellar turbulence is highly supersonic and self-gravitating, there are also observational limitations. Unlike the analytical approach or numerical simulations, molecular cloud observations allow access only to dimensionally reduced information. Velocity measurements are possible only along the line-of-sight, and the spatial structure of a cloud is only seen in projection onto the plane of the sky, i.e. as variations of the column density. Although some methods can yield information about the 3-dimensional spatial structure of the cloud (see Stutzki & Güsten 1990, Williams, De Geus, & Blitz 1994), the result is always model dependent and equivocal (see also Ballesteros-Paredes, Vázquez-Semadeni, & Scalo 1999).
A common way of obtaining knowledge about the velocity structure of molecular clouds is to study individual line profiles at a large number of various positions across the cloud. In the optical thin case line shapes are in fact histograms of the radial velocities of gas sampled along the telescope beam. Falgarone & Phillips (1990) and Falgarone et al. (1994) showed that line profiles constructed from high-sensitivity CO maps exhibit non-Gaussian wings and attributed this to turbulent intermittency (see also Falgarone et al. 1998 on results from the IRAM-key project). Dubinski, Narayan, & Phillips (1995) demonstrated that non-Gaussian line profiles can be produced from any Gaussian random velocity field if variance effects become important (which is always the case for very steep or truncated power spectra). They concluded that non-Gaussian line profiles do not provide clear evidence for intermittency.
Another method of inferring properties of the velocity distribution in molecular clouds is to analyze the pdf of line centroid velocities obtained from a large number of individual measurements scanning the entire projected surface area of a cloud (Miesch & Scalo 1995, Lis et al. 1998, Miesch et al. 1998). Each line profile (i.e. the pdf along the line-of-sight) is collapsed into one single number, the centroid velocity, and then sampled perpendicular to the line-of-sight. Hence, the two functions differ in the direction of the sampling and in the quantity that is considered. A related statistical measure is the pdf of centroid velocity increments, it samples the velocity differences between the centroids for line measurements which are offset by a given separation. The observational advantage of using centroid and increment pdf’s is, that the line measurements can typically be taken with lower sensitivity as only the centroid has to be determined instead of the detailed line shape. These measures are also less dependent on large-scale systematic motions of the cloud and they are less effected by line broadening due to the possible presence of warm dilute gas. However, to allow for a meaningful analysis of the pdf’s especially in the tails, the number of measurements needs to be very large and should not be less than about 1000. In order to sample the entire volume of interstellar clouds, the molecular lines used to obtain the pdf’s are optically thin. I follow this approach in the present investigation and use a mass-weighted velocity sampling along the line-of-sight to determine the line centroid. This zero-opacity approximation does not require any explicit treatment of the radiation transfer process.
The observed pdf’s are obtained from averaged quantities (from column densities or line centroids). To relate these observational measures to quantities relevant for turbulence theory, i.e. to the full 3-dimensional pdf, numerical simulations are necessary as only they allow unlimited access to all variables in phase space. A first attempt to do this was presented by Lis et al. (1996, 1998) who analyzed a simulation of mildly supersonic decaying hydrodynamic turbulence by Porter et al. (1994). Since their model did neither include self-gravity nor consider flows at high Mach number or mechanisms to replenish turbulence, the applicability to the interstellar medium remained limited. This fact prompts the current investigation which extends the previous ones by calculating highly supersonic flows, and by including self-gravity and a turbulent driving scheme. The current study does not consider magnetic fields. Their influence on the pdf’s needs to be addressed separately. However, the overall importance of magnetic fields and MHD waves on the dynamical structure of molecular clouds may not be large. The energy associated with the observed fields is of the order of the (turbulent) kinetic energy content of molecular clouds (Crutcher 1999). Magnetic fields cannot prevent the decay of turbulence (e.g. Mac Low et al. 1998) which implies the presence of external driving mechanisms. These energy sources replenish the turbulent cascade and may excite MHD waves explaining the inferred equipartition between turbulent and magnetic energies.
### 2.3 Statistical Definitions
The one-point probability distribution function $`f(x)`$ of a variable $`x`$ is defined such that $`f(x)dx`$ measures the probability for the variable to be found in the interval $`[x,x+dx]`$. The density pdf ($`\rho `$-pdf) discussed in this paper is obtained from the local density associated with each SPH particle. It is basically the normalized histogram summed over all particles in the simulation, i.e. a mass-weighted sampling procedure is applied. The line-of-sight velocity centroid pdf ($`v`$-pdf) is more complicated to compute. The face of the simulated cube is divided into $`64^2`$ equal-sized cells. For each cell, the line profile is computed by sampling the normal (line-of-sight) velocity component of all gas particles that are projected into that cell. The line centroid is determined as the abscissa value of the peak of the distribution. This procedure corresponds to the formation of optically thin lines in molecular clouds, where all molecules within a certain column through the clouds contribute equally to the shape and intensity of the line. To reduce the sampling uncertainties, this procedure is repeated with the location of the cells shifted by half a cell size in each direction. Altogether about 20 000 lines contribute to the pdf. This is procedure is repeated for line-of-sights along all three system axes to identify projection effects. The line centroid increment pdf ($`\mathrm{\Delta }v`$-pdf) is obtained in a similar fashion. However, the sampled quantity is now the velocity difference between line centroids obtained at two distinct locations separated across the face of the cloud by a fixed shift vector $`\mathrm{\Delta }\stackrel{}{r}`$. The $`\mathrm{\Delta }v`$-pdf for a spatial lag $`\mathrm{\Delta }r`$ is obtained as azimuthal average, i.e. as superposition of all individual pdf’s with shift vectors of length $`\mathrm{\Delta }r`$.
Also statistical moments of the distribution can be used to quantify the spread and shape of pdf’s. For the current analysis I use the first four moments. Mean value $`\mu `$ and standard deviation $`\sigma `$ (the 1. and 2. moments) quantify the location and the width of the pdf and are given in units of the measured quantity. The third and fourth moments, skewness $`\theta `$ and kurtosis $`\kappa `$, are dimensionless quantities characterizing the shape of the distribution. The skewness $`\theta `$ describes the degree of asymmetry of a distribution around its mean. The kurtosis $`\kappa `$ measures the relative peakedness or flatness of the distribution. I use a definition where $`\kappa =3`$ corresponds to a normal distribution. Smaller values indicate existence of a flat peak compared to a Gaussian, larger values point towards a stronger peak or equivalently towards the existence of prominent tails in the distribution. A pure exponential results in $`\kappa =6`$. Gaussian random fields are statistically fully determined by their mean value and the 2-point correlation function, i.e. by their first two moments, $`\mu `$ and $`\sigma `$. All higher moments can be derived from those. The 2-point correlation function is equivalent to the power spectrum in Fourier space (e.g. Bronstein & Semendjajew 1979).
Besides using moments there are other possibilities of characterizing a distribution. Van den Marel & Franx (1993) and Dubinski et al. (1995) applied Gauss-Hermite expansion series to quantify non-normal contributions in line profiles. A more general approach has been suggested by Vio et al. (1994), who discuss alternatives to the histogram representation of pdf’s. However, as astrophysical data sets typically are histograms of various types and as histograms are the most commonly used method to describe pdf’s, this approach is also adopted here.
## 3 The Numerical Model
### 3.1 SPH in Combination with GRAPE
SPH (smoothed particle hydrodynamics) is a particle-based scheme to solve the equations of hydrodynamics. The fluid is represented by an ensemble of particles, each carrying mass, momentum, and hydrodynamic properties. The time evolution of the fluid is represented by the time evolution of the particles, governed by the equations of motion which are supplemented by a prescription to modify the hydrodynamic properties. At any location these properties are obtained by averaging over an appropriate set of neighboring particles. Excellent overviews over the method provide the reviews by Benz (1990) and Monaghan (1992). For the current study I use SPH because it is intrinsically Lagrangian and because it is able to resolve very high density contrasts. Another reason for choosing SPH is the possibility to use it in combination with the special-purpose hardware device GRAPE (Sugimoto et al. 1990, Ebisuzaki et al. 1993; and also Umemura et al. 1993, Steinmetz 1996). This allows calculations at supercomputer level on a normal workstation.
The code is based on a version originally developed by Benz (1990), and is used with a standard description of a von Neumann-type artificial viscosity (Monaghan & Gingold 1983) with the parameters $`\alpha _v=1`$ and $`\beta _v=2`$ for the linear and quadratic terms. The system is subject to periodic boundary conditions (Klessen 1997) and is integrated in time using a second-order Runge-Kutta-Fehlberg scheme, allowing individual time steps for each particle. Furthermore, the smoothing volume over which hydrodynamic quantities are averaged in the code is freely adjustable in space and time such that the number of neighbors for each particle remains approximately fifty. When including self-gravity, regions with masses exceeding the Jeans limit become unstable and collapse. Once a highly-condensed core has formed in the center of a collapsing gas clump, that core is substituted by a ‘sink’ particle (Bate, Bonnell, & Price 1995) which inherits the combined masses, linear and ‘spin’ angular momenta of the particles it replaces. It also has the ability to accrete further SPH particles from its infalling gaseous envelope.
For simulations of turbulent flows one also has to take into account that an explicit viscosity term is introduced in the SPH method. This fact demands attention when studying dissipative processes, especially in the subsonic regime. The current study focuses on the properties of highly supersonic turbulent flows. In this regime, direct comparison between SPH and grid-based methods has proven the close correspondence of both methods (Mac Low et al. 1998, Klessen, Heitsch & Mac Low 2000). If one bears the above caveats in mind, the SPH method calculates the time evolution of gaseous systems very reliably and accurately, and offers large spatial and dynamical flexibility.
### 3.2 Models
The numerical models discussed here describe isothermal gas. The hydrodynamic equations are extended to include self-gravity (in Sec.’s 6 and 7) and to incorporate a random turbulent driving mechanism (in Sec. 7). All physical constants are set to unity. The same applies to mass and length scales, i.e. the total mass is $`M=1`$ and the simulated volume is the cube $`[1,+1]^3`$. The mean density is thus $`\rho =1/8`$. The initial configuration of all dynamical systems discussed in this paper is a homogeneous gas distribution with a Gaussian velocity field. Without turbulence, the time evolution depends on one parameter, the ratio between internal and gravitational energy, $`\alpha ϵ_{\mathrm{int}}/|ϵ_{\mathrm{pot}}|`$. This quantity can be interpreted as dimensionless temperature and determines the number of thermal Jeans masses contained in the system. Molecular clouds are characterized by line widths which largely exceed the thermal broadening. The evolution away from the homogeneous initial state is thus strongly influenced by the adopted initial velocity distribution and depends on whether turbulence is decaying or driven. Large turbulent kinetic energy can considerable slow down or even prevent the collapse of thermally Jeans unstable gas. The situation is very complex and depends on the shape and strength of the turbulent velocity spectrum (Klessen et al. 2000; see also see Bonazzola et al. 1992 and Vázquez-Semadeni & Gazol 1995 for an analytical approach).
To generate and maintain turbulent flows Gaussian velocity fields are introduced. The spatial variations of each component of the velocity vector $`\stackrel{}{v}`$ are described as superpositions of plane waves with wave numbers $`\stackrel{}{k}=(k_x,k_y,k_z)`$, where the phase of each wave is random and sampled from a uniform distribution in the interval $`[0,2\pi [`$. Also the amplitude is random, but selected from a Gaussian distribution centered on zero and with a width determined by the power spectrum $`P(k)=A_kk^\alpha `$. Gaussian fields are isotropic and only depend on the absolute value of the wave vector $`k=|\stackrel{}{k}|`$. Only waves in the range $`1kk_{\mathrm{max}}`$ are considered. For large cut-off wave numbers $`k_{\mathrm{max}}`$ the Gaussian statistics is very well sampled. If only very few modes are used to generate the field, variance effects become strong and individual realizations of the field can deviate significant from the ensemble average (see Sec. 4). The field is then transformed back into real space and the resulting velocities are assigned to individual SPH particles using the ‘cloud-in-cell’ scheme (Hockney & Eastwood 1988). For the initial field, all velocities are multiplied by the appropriate factor to reach the desired rms Mach number of the flow. In case of driven turbulence, this velocity field is also used to ‘kick’ the SPH particles at every time step such that a constant level of kinetic energy is maintained (see Mac Low 1999)
## 4 PDF’s from Gaussian Velocity Fluctuations
Variance effects in poorly sampled Gaussian velocity fields can lead to considerable non-normal contributions to the $`v`$\- and $`\mathrm{\Delta }v`$-pdf’s. If a random process is the result of sequence of independent events (or variables), then in the limit of large numbers, its distribution function will be a Gaussian around some mean value. However, only the properties of a large ensemble of Gaussian fields are determined in a statistical sense. Individual realizations may exhibit considerable deviations from the mean. The effect is strongest when only few (spatial) modes contribute to the field or, almost equivalently, when the power spectrum falls off very steeply. In this case, most kinetic energy is in large-scale motions.
This is visualized in Fig. 1, it shows $`v`$-pdf’s for homogeneous gas (sampled by $`64^3`$ SPH particles placed on a regular grid) with Gaussian velocity fields with power spectra $`P(k)=\mathrm{const}.`$ which are truncated at different wave numbers $`k_{\mathrm{max}}`$ ranging from (a) $`k_{\mathrm{max}}=2`$ to (d) $`k_{\mathrm{max}}=32`$. Each realization is scaled such that the rms velocity dispersion is $`\sigma _v=0.5`$. The figure displays the pdf’s for the $`x`$-, $`y`$-, and $`z`$-component of the velocity. The pdf’s of the strongly truncated spectrum (Fig. 1a) do not at all resemble normal distributions. The Gaussian statistics of the field is very badly sampled with only very few modes. Note that the pdf’s of the same field may vary considerably for different velocity components, i.e. for different projections. With the inclusion of larger number of Fourier modes this situation improves, and in Fig. 1d the pdf’s of all projections sample the expected Gaussian distribution very well.
A similar conclusion can be derived for $`\mathrm{\Delta }v`$-pdf. This measure is even more sensitive to deviations from Gaussian statistics. Figure 2 plots the $`\mathrm{\Delta }v`$-pdf’s for the same sequence of velocity fields. For brevity, only the line-of-sight component parallel to the $`x`$-axis is considered. Furthermore, from the sequence of possible $`\mathrm{\Delta }v`$-pdf’s (defined by the spatial lag $`\mathrm{\Delta }r`$) only three are shown, at small ($`\mathrm{\Delta }r=1/32`$, upper curve), medium ($`\mathrm{\Delta }r=10/32`$, middle curve), and large spatial lags ($`\mathrm{\Delta }r=30/32`$, upper curve). Sampling the Gaussian field with only two modes (Fig. 2a) is again insufficient to yield increment pdf’s of normal shape. The velocity field is very smooth, and the line centroid velocity difference between neighboring cells is very small. Hence, for $`\mathrm{\Delta }r=1/32`$ the pdf is dominated by a distinct central peak at $`\mathrm{\Delta }v=0`$. The tails of the distribution are quite irregularly shaped. The situation becomes ‘better’ when sampling increasing distances, as regions of the fluid separated by larger $`\mathrm{\Delta }r`$ are less strongly correlated in velocity. For $`\mathrm{\Delta }r=10/32`$ and $`\mathrm{\Delta }r=30/32`$ the pdf’s follow the Gaussian distribution more closely although irregularities in the shapes are still present. In Fig.’s 2b and c the $`\mathrm{\Delta }v`$-pdf’s for medium to large lags are very well fit by Gaussians. Deviations occur only at small $`\mathrm{\Delta }r`$, the pdf’s are exponential (and the distribution for $`k_{\mathrm{max}}=4`$ is still a bit cuspy). Finally, Fig. 2d shows the three $`\mathrm{\Delta }v`$-pdf’s for the case where all available spatial modes contribute to the velocity field ($`1k32`$). The pdf’s follow a Gaussian for all spatial lags.
This behavior is also seen in the variation of the moments of the distribution as function of the spatial lag $`\mathrm{\Delta }r`$. Applied to the above sequence of Gaussian velocity fields, Fig. 3 displays the dispersion $`\sigma `$ and the kurtosis $`\kappa `$ of the distribution. The corresponding models are indicated at the right hand side of each plot. The width of the distribution, as indicated by the dispersion $`\sigma `$ (Fig. 3a), typically grows with increasing $`\mathrm{\Delta }r`$, reflecting the relative peakedness of the distribution at small lags. For example, the distribution (a) yields a slope of 0.3 in the range $`0.6\mathrm{log}_{10}\mathrm{\Delta }r0.4`$, and (b) leads to a value of 0.2 in relatively large interval $`1.5\mathrm{log}_{10}\mathrm{\Delta }r0.5`$. The effect disappears for the better sampled fields. Typical values for that slope in observed molecular clouds are $`0.3`$ to $`0.5`$ (Miesch et al. 1998).<sup>1</sup><sup>1</sup>1Note, that Miesch et al. (1998) are plotting the function $`\sigma ^2`$ versus the spatial lag $`\mathrm{\Delta }r`$. For a comparison with the present study, their numbers have to be divided by a factor of two. Furthermore, they use a relatively narrow range of $`\mathrm{\Delta }r`$-values to compute the slope of the function; larger intervals would on average tend to decrease these values (see their Fig. 14). In addition, Miesch et al. (1998) applied spatial filtering to remove large-scale velocity gradients in the clouds. These would lead to steeper slopes. The fact that in the present study the functions $`\sigma `$ and $`\kappa `$ level out for large spatial lags $`\mathrm{\Delta }r`$ is a consequence of the periodic boundary conditions which do not allow for large-scale gradients. A direct measure of the peakedness of the distribution is its fourth moment, the kurtosis $`\kappa `$ (Fig. 3b). At small lags $`\mathrm{\Delta }r`$, clearly the pdf’s of model (a) are more strongly peaked than exponential ($`\kappa =6`$). Comparing the entire sequence reveals again the tendency of the pdf’s to become Gaussian at decreasing $`\mathrm{\Delta }r`$ with increasing number of modes considered in the construction of the velocity field.
Taking all together, it is advisable to consider conclusions about interstellar turbulence derived from solely analyzing one-point probability distribution functions from molecular clouds with caution. Similar to what has been shown by Dubinski et al. (1995) for molecular line profiles, deviations from the regular Gaussian shape found in $`v`$\- and $`\mathrm{\Delta }v`$-pdf’s need not be the signpost of turbulent intermittency. Gaussian velocity fields which are dominated by only a small number of modes (either because the power spectrum falls off steeply towards larger wave numbers, or because small wave length distortions are cut away completely) will lead to very similar distortions. In addition, the properties of the pdf may vary considerably between different projections. The same velocity field may lead to smooth and Gaussian pdf’s for one velocity component, whereas another projection may result in strong non-Gaussian wings (see also Fig. 9).
## 5 Analysis of Decaying Supersonic Turbulence without Self-Gravity
In this section the pdf’s of freely decaying initially highly supersonic turbulence without self-gravity are discussed. They are calculated from an SPH simulation with $`\mathrm{350\hspace{0.17em}000}`$ particles (Mac Low et al. 1998, model G). Initially the system is homogeneous with a Gaussian velocity distribution with $`P(k)=\mathrm{const}.`$ in the interval $`1k8`$. The rms Mach number of the flow is $`M=5`$.
After the onset of the hydrodynamic evolution the flow quickly becomes fully turbulent resulting in rapid dissipation of kinetic energy. The energy decay is found to follow a power law $`t^\eta `$ with exponent $`\eta =1.1\pm 0.004`$. The overall evolution can be subdivided into several phases. The first phase is very short and is defined by the transition of the initially Gaussian velocity field into fully developed supersonic turbulence. It is determined by the formation of the first shocks which begin to interact with each other and build up a complex network of intersecting shock fronts. Energy gets transfered from large to small scales and the turbulent cascade builds up. The second phase is given by the subsequent self-similar evolution of the network of shocks. Even though individual features are transient, the overall properties of this network change only slowly. In this phase of highly supersonic turbulence the loss of kinetic energy is dominated by dissipation in shocked regions. In the transsonic regime, i.e. the transition from highly supersonic to fully subsonic flow, energy dissipation in vortices generated by shock interactions becomes more and more important. Only the strongest shocks remain in this phase. Surprisingly, the energy decay law does not change during this transition. It continues to follow a power law with exponent $`\eta 1`$. In the subsonic phase the flow closely resembles incompressible turbulence. Its properties are similar to those reported from numerous experiments and simulations (e.g. Porter et al. 1994, Lesieur 1997, Boratav et al. 1997). The simulation is stopped at $`t=20.0`$ when the flow has decayed to a rms Mach number of $`M=0.3`$. Since the energy loss rate follows a power law, the duration of each successive phase grows.
This sequence of evolutionary stages is seen in the pdf’s of the system. One noticeable effect is the decreasing width of the distribution functions as time progresses. As the kinetic energy decays the available range of velocities shrinks. This not only leads to ‘smaller’ $`v`$\- and $`\mathrm{\Delta }v`$-pdf’s, but also to a smaller $`\rho `$-pdf since compressible motions lose influence and the system becomes more homogeneous. This is indicated in Fig. 4, it displays (a) the $`rho`$-pdf and (b) $`v`$-pdf at the following stages of the dynamical evolution (from top to bottom): Shortly after the start, at $`t=0.2`$ when the first shocks occur, then at $`t=0.6`$ when the network of interacting shocks is established and supersonic turbulence is fully developed, during the transsonic transition at $`t=3.5`$, and finally at $`t=20.0`$ when the flow has progressed into the subsonic regime. The rms Mach numbers at these stages are $`M=5.0`$, $`M=2.5`$, $`M=1.0`$, and $`M=0.3`$, respectively. The density pdf always closely follows a log-normal distribution, i.e. it is Gaussian in the logarithm of the density. Also the distribution of line centroids at the four different evolutionary stages of the system is best described by a Gaussian with only minor deviations at the far ends of the velocity spectrum.
For the same points in time, Fig. 5 shows the $`\mathrm{\Delta }v`$-pdf’s for $`x`$-component of the velocity. The displayed spatial lags are selected in analogy to Fig. 2. Note the different velocity scaling in each plot reflecting the decay of turbulent energy as the system evolves in time. Throughout the entire sequence, spatial lags larger than about 10% of the system size always lead to $`\mathrm{\Delta }v`$-pdf’s very close to Gaussian shape (the middle and lower curves). Considerable deviations occur only at small spatial lags (the upper curves). For those, the increment pdf’s exhibit exponential wings during all stages of the evolution. When scaling the pdf’s to the same width, the distribution in the subsonic regime (d) appears to be more strongly peaked than during the supersonic or transsonic phase (a – c). There, the central parts of the pdf’s are still reasonably well described by the Gaussian obtained from the first two moments, whereas in (d) the peak is considerably narrower, or vice versa, the tails of the distribution are more pronounced.
These results can be compared with the findings by Lis et al. (1998). They report increment pdf’s for three snapshots of a high-resolution hydrodynamic simulation of decaying mildly super-sonic turbulence performed by Porter et al. (1994). They analyze the system at three different times corresponding to rms Mach numbers of $`M0.96`$, $`M0.88`$, and $`M0.52`$. Their first two data sets thus trace the transition from supersonic to subsonic flow and are comparable to phase (c) of the current model; their last data set corresponds to to phase (d). In the transsonic regime both studies agree: Lis et al. (1998) report enhanced tails in the increment pdf’s for the smallest spatial lags which they considered and near Gaussian distributions for larger lags (however, the largest separation they study is about 6% of the linear extent of the system). In the subsonic regime, Lis et al. (1998) find near Gaussian pdf’s for very small spatial lags ($`<1`$%), but extended wings in the pdf’s for lags of 3% and 6% of the system size. They associate this with the ‘disappearance’ of large-scale structure. Indeed, their Fig. 7 exhibits a high degree of fluctuations on small scales which they argue become averaged away when considering small spatial lags in the $`\mathrm{\Delta }v`$-pdf. Comparing the pdf with spatial lags of 3% (upper curves in Fig. 5, compared to the pdf’s labeled with $`\mathrm{\Delta }=15`$ in Lis et al. 1998) both studies come to the same result. At these scales the $`\mathrm{\Delta }v`$-pdf’s tend to exhibit more pronounced wings in the subsonic regime as in the supersonic regime. The SPH calculations reported here do not allow for a meaningful construction of $`\delta v`$-pdf’s for $`\mathrm{\Delta }r<3`$%. The Gaussian behavior of pdf’s for very small spatial lags reported by Lis et al. (1998) therefore cannot be examined. However, neither of the purely hydrodynamic simulations lead to pdf’s that are in good agreement with the observations. Observed pdf’s typically are much more centrally peaked at small spatial separation (see e.g. Fig. 4 in Lis et al 1998 and Miesch et al. 1998).
Figure 6 shows the spatial distribution of centroid velocity differences between cells separated by a vector lag of $`\mathrm{\Delta }\stackrel{}{r}=(1/32,1/32)`$ (i.e. between neighboring cells along the diagonal). Data are obtained at the same times as above. Each figure displays the array of the absolute values of the velocity increments $`\mathrm{\Delta }v_x`$ in linear scaling as indicated at the top. Note the decreasing velocity range reflecting the decay of turbulent energy. The distribution of $`\mathrm{\Delta }v_x`$ appears random, there is no clear indication for coherent structures. This is corresponds to most observations. Miesch et al. (1998) find for their sample of molecular clouds that high-amplitude velocity differences for very small spatial lags typically are well distributed resulting in a ‘spotty’ appearance. Note, however, that using azimuthal averaging Lis et al. 1998 report the finding of filamentary structures for the $`\rho `$-Ophiuchus cloud. Altogether, filamentary structure is difficult to define and a mathematical thorough analysis is seldomly performed (for an astrophysical approach see Adams & Wiseman 1994, for a discussion of the filamentary vortex structure in incompressible turbulence consult Frisch 1995 or Lesieur 1997). The visual inspection of maps is often misleading and influenced by the parameters used to display the image. Larger velocity bins for instance tend to produce a more ‘filamentary’ structure than very fine sampling of the velocity structure. Further uncertainty may be introduced by the fact that molecular clouds are only seen in one projection as the signatures of the dynamical state of the system can strongly depend on the viewing angle.
## 6 Analysis of Decaying Turbulence with Self-Gravity
In this section, I discuss the properties of decaying, initially supersonic turbulence in a self-gravitating medium. Figure 7 displays an SPH simulation with $`\mathrm{200\hspace{0.17em}000}`$ particles at six different times of its dynamical evolution. Since the model is subject to periodic boundary conditions, every figure has to be considered infinitely replicated in each direction. Analog to the previous model, the system is initially homogeneous and its velocity field is generated with $`P(k)=\mathrm{const}.`$ using modes with wave numbers $`1k8`$. From the choice $`\alpha =0.01`$ it follows that the system contains 120 thermal Jeans masses. The initial rms velocity dispersion is $`\sigma _v=0.5`$ and with the sound speed $`c_\mathrm{s}=0.082`$ the rms Mach number follows as $`M=6`$. These values imply that the initial turbulent velocity field contains sufficient energy to globally stabilize the system against gravitational collapse. Scaled to physical units using a density $`n(\mathrm{H}_2)=10^5`$cm<sup>-3</sup>, which is typical for massively star-forming regions (e.g. Williams et al. 2000), the system corresponds to a volume of $`[0.32\mathrm{pc}]^3`$ and contains a gas mass of 200 M. As the simulation starts, the system quickly becomes fully turbulent and loses kinetic energy. Like in the case without self-gravity a network of intersecting shocks develops leading to density fluctuations on all scales. If the mass of a fluctuation exceeds the (local) Jeans limit it begins to contract due to self-gravity. During the early evolution, there is enough kinetic energy to prevent this collapse process on all scales (Fig. 7b – $`t=0.5`$) and the properties of the system are similar to those of pure hydrodynamic turbulence. However, as time progresses and turbulent energy decays the effective Jeans mass decreases. Local collapse of shock generated density fluctuations sets in despite the fact that the system is still globally stabilized by turbulence (see also Klessen et al. 2000). The central high-density cores of collapsing clumps are indicated by black dots. The cores form mainly at the intersection of filaments, where the density is highest and local collapse is most likely to set in. When turbulence is decayed sufficiently also large-scale collapse becomes possible. Gas clumps follow the global flow pattern towards a common center of gravity where they may merge or sub-fragment. Gradually a cluster of dense cores is built up. In the isothermal model this process continues until all available gas is accreted onto the ‘protostellar’ cluster (for more details Klessen & Burkert 2000).
The pdf’s of (a) the density and of (b) the $`x`$-component of the line centroid velocities for the above six model snapshots are displayed in Fig. 8. The corresponding time is indicated by the letters at the right side of each panel. During the dynamical evolution of the system the density distribution develops a high density tail. This is the imprint of local collapse. The densities of compact cores are indicated by solid dots (at $`t=2.0`$ and $`t=2.5`$). Virtually all particles in the high density tails at earlier times (at $`t=1.0`$ and more so at $`t=1.5`$) are accreted onto these cores. The bulk of matter roughly follows a log-normal density distribution as indicated by the dotted parabola. The $`v`$-pdf’s are nearly Gaussian as long as the dynamical state of the system is dominated by turbulence. Also the width of the pdf remains roughly constant during this phase. This implies that the decay of turbulent kinetic energy is in balance with the gain of kinetic energy due to gravitational (‘quasi-static’) contraction on large scales. The time scale for this process is determined by the energy dissipation in shocks and turbulent eddies. However, once localized collapse is able to set in, accelerations on small scales increase dramatically and the evolution ‘speeds up’. For times $`t>2.0`$ the centroid pdf’s become wider and exhibit significant deviations from the original Gaussian shape. The properties of the pdf’s are similar to those observed in star-forming regions (Miesch & Scalo 1995, Lis et al. 1998, Miesch et al. 1998). This is expected since gravitational collapse is a necessary ingredient for forming stars.
Gravity creates non-isotropic density and velocity structure structures. When analyzing $`v`$\- and $`\mathrm{\Delta }v`$-pdf’s, their appearance and properties will strongly depend on the viewing angle. This is a serious point of caution when interpreting observational data, as molecular cloud are seen only in one projection. As illustration, Fig. 9 plots the centroid pdf at the time $`t=2.0`$ for the line-of-sight projection along all three axes of the system. Whereas the pdf’s for the $`x`$\- and the $`y`$-component of the velocity centroid are highly structured (upper and middle curve – the latter one is even double peaked), the distribution of the $`z`$-component (lowest curve) is smooth and much smaller in width, comparable to the ‘average’ pdf at earlier stages of the evolution. As the variations between different viewing angles or equivalently different velocity components can be very large, statements about the 3-dimensional velocity structure from only observing one projection can be misleading.
Gravity effects the $`\mathrm{\Delta }v`$-pdf. Figure 10 displays the increment pdf’s at small, intermediate and large spatial lags, analog to Fig.’s 2 and 5. Time ranges from (a) $`t=1.0`$ to (d) $`t=2.5`$ corresponding to Fig.’s 7c–f. The pdf’s for $`t=0.0`$ and $`t=0.5`$ are not shown since at these stages supersonic turbulence dominates the dynamic of the system and the pdf’s are comparable to the ones without gravity (Fig. 5). This still holds for $`t=1.0`$. The increment pdf’s for medium to large spatial lags appear Gaussian, however, the pdf for the smallest lag follows a perfect exponential all the way inwards to $`\mathrm{\Delta }v=0`$. Unlike in the case without gravity, the peak of the distribution is not ‘round’, i.e. is not Gaussian in the innermost parts (when scaled to the same width). It is a sign of self-gravitating systems that the increment pdf at smallest lags is very strongly peaked and remains exponential over the entire range of measured velocity increments. This behavior is also seen Fig.’s 10b–d. At these later stages of the evolution in addition non-Gaussian behavior is also found at medium lags. This results from the existence of large-scale filaments and streaming motions. The same behavior is found for the increment pdf’s from observed molecular clouds (for $`\rho `$-Ophiuchus see Lis et al 1998; for Orion, Mon R2, L1228, L1551, and HH83 see Miesch et al. 1998). In each case, the distribution for the smallest lag (one pixel size) is very strongly peaked at $`\mathrm{\Delta }v=0`$, in some cases even more than exponential. The deviations from the Gaussian shape remain for larger lags but are not so pronounced. The inclusion of self-gravity into models of interstellar turbulence leads to good agreement with the observed increment pdf’s. However, this result may not be unique as in molecular clouds additional processes are likely to be present that could also lead to strong deviations from Gaussianity.
The time evolution of the statistical moments of the $`\mathrm{\Delta }v`$-pdf’s for various spatial lags is presented in Fig. 11. It plots (a) the dispersion $`\sigma `$, and (b) the kurtosis $`\kappa `$. The letters on the right-hand side indicate the corresponding time in Fig. 7. At $`t=0.0`$ the width $`\sigma `$ of the pdf is approximately constant for all $`\mathrm{\Delta }r`$ and the kurtosis $`\kappa `$ is close to normal value of three. Both indicate that Gaussian statistics very well describes the initial velocity field. As turbulent energy decays, gravitational collapse sets in. Because of the gravitational acceleration, the amplitudes of centroid velocity differences between separate regions in the cloud grow larger, the width $`\sigma `$ of the $`\mathrm{\Delta }v`$-pdf’s increases. This becomes more important when sampling velocity differences on larger spatial scales, hence $`\sigma `$ also increases with $`\mathrm{\Delta }r`$. The slope is $`d\mathrm{log}_{10}\sigma /d\mathrm{log}_{10}\mathrm{\Delta }r\stackrel{<}{}\mathrm{\hspace{0.25em}0.2}`$. For $`\mathrm{log}_{10}\mathrm{\Delta }r>0.4`$ it levels out, which is a result of the adopted periodic boundary conditions. They do not allow for large-scale velocity gradients. The increasing ‘peakedness’ of $`\mathrm{\Delta }v`$-pdf is reflected in the large values of the kurtosis $`\kappa `$ at the later stages of the evolution. For small spatial lags the pdf’s are more centrally concentrated than exponential (i.e. $`\kappa >6`$), and even at large spatial separations they are still more strongly peaked than Gaussian ($`\kappa >3`$). The slope at $`t=2.5`$ is $`d\mathrm{log}_{10}\kappa /d\mathrm{log}_{10}\mathrm{\Delta }r0.4`$ which is indeed comparable to what is found in observed star-forming regions (Miesch et al. 1998).
For the above simulation of self-gravitating, decaying, supersonic turbulence, Fig.12 plots the 2-dimensional distribution of centroid increments for a vector lag $`\mathrm{\Delta }\stackrel{}{r}=(1/32,1/32)`$. The velocity profiles are sampled along the $`x`$-axis of the system. The magnitude of the velocity increment $`\mathrm{\Delta }v_x`$ is indicated at the top of each plot. The spatial distribution of velocity increments during the initial phases appears random. Later on, gravity gains influence over the flow and creates a network of intersecting filaments where gas streams onto and flows along towards local potential minima. At that stage, the velocity increments with the highest amplitudes tend to trace the large-scale filamentary structure. This is the sign of the anisotropic nature of gravitational collapse motions.
## 7 Analysis of Driven Turbulence with Self-Gravity
Figure 13 displays the gas distribution at different evolutionary stages of a simulation of driven, supersonic, self-gravitating turbulence. The number of SPH particles is $`\mathrm{205\hspace{0.17em}379}`$. Again, the system is initially homogeneous in space and has a random Gaussian velocity field with flat power spectrum in the wave number interval $`3k4`$. It contains 64 thermal Jeans masses and turbulence is continuously driven as described in Sec. 3.2. The initial evolution into equilibrium between the energy input by the driving force and the decay of turbulent kinetic energy is computed without self-gravity, then it is turned on. This phase is displayed in Fig. 13a. In this state the turbulent Jeans mass (on scales larger than the maximum driving wave length) exceeds the total mass in the system by a factor of two, the cloud is therefore stabilized by turbulence against gravitational collapse on global scales. However, local collapse (on scales at or below the driving scale) is still possible and does occur. As in the previous case without driving, the dynamical evolution of the system leads to the formation of a cluster of dense collapsed cores. This is shown in Fig.’s 13b–d, which display the system when 20%, 40%, and 60% of the gas mass has accumulated in dense collapsed cores (at time $`t=1.8`$, $`t=3.2`$, and $`t=4.8`$, respectively). However, in the presence of the driving source the time scales for accretion are longer and the cluster is less dense.
The pdf’s of (a) the density and (b) the $`x`$-component of the line centroid velocities corresponding to the above four snapshots are displayed in Fig. 14. As in the previous model, the bulk of gas particles that are not accreted onto cores build up an approximately log-normal $`\rho `$-pdf (indicated by the dotted lines). Also the $`v`$-pdf remains close to the Gaussian value. This is different from the case of purely decaying self-gravitating turbulence, where at some stage global collapse motions set in and lead to very wide and distorted centroid pdf’s. This is not possible in the simulation of driven turbulence, as it is stabilized on the largest scales by turbulence. Collapse occurs only locally which leaves the width of the pdf’s relatively unaffected and only mildly alters their shape.
Also the $`\mathrm{\Delta }v`$-pdf’s show no obvious sign of evolution. For the $`x`$-component of the velocity these functions are displayed in Fig. 15, again for three different spatial lags. The chosen times correspond (a) to the equilibrium state at $`t=0.0`$, and (b) to $`t=4.8`$ which is the final state of the simulation. The pdf’s only marginally grow in width. At every evolutionary stage, the pdf for the smallest spatial lag is exponential, whereas the pdf’s for medium and large shift vectors closely follow the Gaussian curve defined by the first two moments of the distribution (dotted lines). The functions are similar to the ones in the previous model before the large scale collapse motions set in (Fig. 10a, b). Only overall contraction will affect $`\mathrm{\Delta }v`$-pdf at medium to large lags. This behavior also follows from comparing the statistical moments. Figure 16 plots (a) the dispersion $`\sigma `$ and (b) the kurtosis $`\kappa `$ as function of the spatial lag $`\mathrm{\Delta }r`$. Figures 11a and 16a are very similar, as soon as turbulence is established the width $`\sigma `$ of the pdf increases with $`\mathrm{\Delta }r`$ with a slope of $`d\mathrm{log}_{10}\sigma /d\mathrm{log}_{10}\mathrm{\Delta }r\stackrel{<}{}\mathrm{\hspace{0.25em}0.2}`$ for small to medium lags and levels out for larger ones. However, when comparing the ‘peakedness’ of the pdf as indicated by $`\kappa `$ (Fig.’s 11b and 16b) the model of decaying self-gravitating turbulence yields much higher values since the pdf’s are more strongly peaked due to the presence of large-scale collapse motions.
Figure 17 finally shows the spatial distribution of the $`x`$-component of the line centroid increments for a vector lag $`\mathrm{\Delta }\stackrel{}{r}=(1/32,1/31)`$. Since the increment maps at different evolutionary times are statistically indistinguishable, only times (a) $`t=0.0`$ and (b) $`t=4.8`$ are displayed in the figure. As in the case of supersonic, purely hydrodynamic turbulence the spatial distribution of velocity increments appears random and uncorrelated.
The adopted driving mechanism prevents global collapse. The bulk properties of the system therefore resemble hydrodynamic, non-self-gravitating turbulence. However, local collapse motions do exist and are responsible for noticeable distortions away from the Gaussian statistics. As the non-local driving scheme adopted here introduces a bias towards Gaussian velocity fields, these distortions are not very large. There is a need to introduce other, more realistic driving agents into this analysis. These could lead to much stronger non-Gaussian signatures in the pdf’s.
## 8 Summary
SPH simulations of driven and decaying, supersonic, turbulent flows with and without self-gravity have been analyzed in this study . It extends previous investigations of mildly supersonic, decaying, non-self-gravitating turbulence (Lis et al. 1996, 1998) into a regime more relevant molecular clouds, by (a) considering highly supersonic flows and by including (b) self-gravity and (c) a driving source for turbulence.
The flow properties are characterized by using the probability distribution functions of the density, of the line-of-sight velocity centroids, and of their increments. Furthermore the dispersion and the kurtosis of the increment pdf’s are discussed, as well as the spatial distribution of the velocity increments for the smallest spatial lags.
(1) To asses the influence of variance effects, simple Gaussian velocity fluctuations are studied. The insufficient sampling of random Gaussian ensembles leads to distorted pdf’s similar to the observed ones. For line profiles this has been shown by Dubinski et al. (1995).
(2) Decaying, initially highly supersonic turbulence without self-gravity leads to pdf’s which also exhibit deviations from Gaussianity. For the trans- and subsonic regime this has been reported by Lis et al. (1996, 1998). However, neglecting gravity and thus not allowing for the occurance of collapse motions, these distortions are not very pronounced and cannot account well for the observational data (Lis et al. 1998, Miesch et al. 1998).
(3) When including gravity into the models of decaying initially supersonic turbulence, the pdf’s get into better agreement with the observations. During the early dynamical evolution of the system turbulence carries enough kinetic energy to prevent collapse on all scales. In this phase the properties of the system are similar to those of non-gravitating hydrodynamic supersonic turbulence. However, as turbulent energy decays gravitational collapse sets in. First localized and on small scales, but as the turbulent support continues to diminish collapse motions include increasingly larger spatial scales. The evolution leads to the formation of an embedded cluster of dense protostellar cores (see also Klessen & Burkert 2000). As the collapse scale grows, the $`\rho `$-, $`v`$-, and $`\mathrm{\Delta }v`$-pdf’s get increasingly distorted. In particular, the $`\mathrm{\Delta }v`$-pdf’s for small spatial lags are strongly peaked and exponential over the entire range of measured velocities. This is very similar to what is observed in molecular clouds (for $`\rho `$-Ophiuchus see Lis et al 1998; for Orion, Mon R2, L1228, L1551, and HH83 see Miesch et al. 1998).
(4) The most realistic model for interstellar turbulence considered here includes a simple (non-local) driving scheme. It is used to stabilize the system against collapse on large scales. Again non-Gaussian pdf’s are observed. Despite global stability, local collapse is possible and the system again evolves towards the formation of an embedded cluster of accreting protostellar cores. As the adopted driving scheme introduces a bias towards maintaining a Gaussian velocity distribution, the properties of the pdf’s fall in between the ones of pure hydrodynamic supersonic turbulence and the ones observed in systems where self-gravity dominates after sufficient turbulent decay. This situation may change when considering more realistic driving schemes.
(5) A point of caution: The use of $`v`$\- and $`\mathrm{\Delta }v`$-pdf’s to unambiguously characterize interstellar turbulence and to identify possible physical driving mechanisms may be limited. All models considered in the current analysis lead to non-Gaussian signatures in the pdf’s, differences are only gradual. In molecular clouds the number of physical processes that are expected to give rise to deviations from Gaussian statistics is large. Simple statistical sampling effects (Sec. 4) and turbulent intermittency caused by vortex motion (Lis et al. 1996, 1998), as well as the effect self-gravity (Sec. 6) and of shock interaction in highly supersonic flows (Mac Low & Ossenkopf 2000), all will lead to non-Gaussian signatures in the observed pdf’s. Also stellar feedback processes, galactic shear and the presence of magnetic fields will influence the interstellar medium and create distortions in the velocity field. This needs to be studied in further detail. In addition, the full 3-dimensional spatial and kinematical information is not accessible in molecular clouds, measured quantities are always projections along the line-of-sight. As the structure of molecular clouds is extremely complex, the properties of the pdf’s may vary considerably with the viewing angle. Attempts to disentangle the different physical processes influencing interstellar turbulence therefore should no rely on analyzing velocity pdf’s alone, they require additional statistical information.
I thank A. Burkert, F. Heitsch, and M.-M. Mac Low for many fruitful and stimulating discussions, and the editor S. Shore for his comments on the paper and his help with an extremely slow and non-responsive (anonymous) referee. |
warning/0001/hep-th0001131.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The conjectured equivalence between string theory on anti-de Sitter (AdS) spaces (times some compact manifold) and certain superconformal gauge theories living on the boundary of AdS has led to an increasing interest in black objects in asymptotically anti-de Sitter spaces. On one hand, this interest is based on the fact that the classical supergravity solution can furnish important information on the dual gauge theory in the large $`N`$ limit, $`N`$ being the rank of the gauge group. An example of this is the Hawking-Page phase transition from thermal AdS space to the Schwarzschild-AdS black hole, which was later reconsidered by Witten in the spirit of the AdS/CFT correspondence . There it was observed that it can be interpreted as a transition from a low-temperature confining to a high temperature deconfining phase in the dual field theory.
On the other hand, the proposed AdS/CFT equivalence opens the possibility to a microscopic understanding of the Bekenstein-Hawking entropy of asymptotically anti-de Sitter black holes. This route was pioneered by Strominger , who used the central charge of the AdS<sub>3</sub> asymptotic symmetry algebra to count the microstates giving rise to the BTZ black hole entropy.
Of particular interest in this context are black objects in AdS space which preserve some fraction of supersymmetry. On the CFT side, these supergravity vacua could correspond to an expansion around non-zero vacuum expectation values of certain operators. Supersymmetry of Reissner-Nordström-AdS black holes in four dimensions was first studied by Romans in the context of $`N=2`$ gauged supergravity . These considerations have been extended later in various directions . One common feature of all results is the appearance of naked singularities in the BPS limit<sup>1</sup><sup>1</sup>1It is worth pointing out, however, that by including rotation , or by allowing for different event horizon geometries , one can get genuine BPS black holes in anti-de Sitter space.. This means that the theory on the bulk side is ill-defined in the limit of small distances, and stringy corrections have to be taken into account.
Although many string- and brane solutions in ungauged supergravity theories are known, very little is known on the corresponding objects in the gauged case. To remedy this will be the main purpose of the present paper. In particular, we will derive supersymmetric black string solutions with various topologies in five dimensional $`N=2`$ supergravity theories coupled to vector multiplets . The theory of ungauged five-dimensional $`N=2`$ supergravity coupled to abelian vector supermultiplets can be obtained by compactifying eleven-dimensional supergravity, the low-energy theory of M-theory, on a Calabi-Yau three-fold . Gauged supergravity theories are obtained by gauging a subgroup of the R-symmetry group, the automorphism group of the supersymmetry algebra. The gauged $`D=5`$, $`N=2`$ supergravity theories are obtained by gauging the $`U(1)`$ subgroup of the $`SU(2)`$ automorphism group of the superalgebra . The Lagrangian of the theory is obtained by introducing a linear combination of the abelian vector fields already present in the ungauged theory, i. e. $`A_\mu =V_IA_\mu ^I`$, with a coupling constant $`g`$. The coupling of the Fermi-fields to the $`U(1)`$ vector field breaks supersymmetry, and therefore gauge-invariant $`g`$-dependent terms must be introduced in order to preserve $`N=2`$ supersymmetry. In a bosonic background, this amounts to the addition of a $`g^2`$-dependent scalar potential $`V`$ .
Our work in this paper is organized as follows. Section 2 contains a brief review of $`D=5`$, $`N=2`$ gauged supergravity. In 3, the supersymmetric string solutions are derived, and their near horizon limit is studied. In section 4, general supersymmetric product space compactifications of $`D=5`$, $`N=2`$ gauged supergravity are considered. Finally, our results are summarized and discussed in 5.
## 2 $`D=5`$, $`N=2`$ Gauged Supergravity
The bosonic part of the gauged supersymmetric $`N=2`$ Lagrangian which describes the coupling of vector multiplets to supergravity is given by
$`e^1`$ $`=`$ $`{\displaystyle \frac{1}{2}}R+g^2V{\displaystyle \frac{1}{4}}G_{IJ}F_{\mu \nu }{}_{}{}^{I}F_{}^{\mu \nu J}{\displaystyle \frac{1}{2}}𝒢_{ij}_\mu \varphi ^i^\mu \varphi ^j`$ (1)
$`+{\displaystyle \frac{e^1}{48}}ϵ^{\mu \nu \rho \sigma \lambda }C_{IJK}F_{\mu \nu }^IF_{\rho \sigma }^JA_\lambda ^K,`$
where $`\mu ,\nu `$ are spacetime indices, $`R`$ is the scalar curvature, $`F_{\mu \nu }^I`$ denote the abelian field-strength tensors, and $`e=\sqrt{g}`$ is the determinant of the Fünfbein $`e_\mu ^a`$. The scalar potential $`V`$ is given by
$$V(X)=V_IV_J\left(6X^IX^J\frac{9}{2}𝒢^{ij}_iX^I_jX^J\right),$$
where $`X^I`$ represent the real scalar fields satisfying the condition $`𝒱=\frac{1}{6}C_{IJK}X^IX^JX^K=1.`$ The physical quantities in (1) can all be expressed in terms of the homogeneous cubic polynomial $`𝒱`$ which defines a “very special geometry”. We also have the relations
$`G_{IJ}`$ $`=`$ $`{\displaystyle \frac{1}{2}}_I_J\mathrm{log}𝒱|_{𝒱=1},`$
$`𝒢_{ij}`$ $`=`$ $`_iX^I_jX^JG_{IJ}|_{𝒱=1},`$ (2)
where $`_i`$ and $`_I`$ refer, respectively, to a partial derivative with respect to the scalar field $`\varphi ^i`$ and $`X^I=X^I(\varphi ^i)`$.
Note that for Calabi-Yau compactification of M-theory, $`𝒱`$ is the intersection form, $`X^I`$ and $`X_I=\frac{1}{6}C_{IJK}X^JX^K`$ correspond to the size of the two- and four-cycles and $`C_{IJK}`$ are the intersection numbers of the Calabi-Yau threefold.
The supersymmetry transformations of the gravitino $`\psi _\mu `$ and the gauginos $`\lambda _i`$ in a bosonic background read
$`\delta \psi _\mu `$ $`=`$ $`(𝒟_\mu +{\displaystyle \frac{i}{8}}X_I(\mathrm{\Gamma }_\mu {}_{}{}^{\nu \rho }4\delta _\mu {}_{}{}^{\nu }\mathrm{\Gamma }_{}^{\rho })F_{\nu \rho }{}_{}{}^{I}+{\displaystyle \frac{1}{2}}g\mathrm{\Gamma }_\mu X^IV_I{\displaystyle \frac{3i}{2}}gV_IA_\mu ^I)ϵ,`$ (3)
$`\delta \lambda _i`$ $`=`$ $`\left({\displaystyle \frac{3}{8}}\mathrm{\Gamma }^{\mu \nu }F_{\mu \nu }^I_iX_I{\displaystyle \frac{i}{2}}𝒢_{ij}\mathrm{\Gamma }^\mu _\mu \varphi ^j+{\displaystyle \frac{3i}{2}}gV_I_iX^I\right)ϵ,`$ (4)
where $`ϵ`$ is the supersymmetry parameter and $`𝒟_\mu `$ is the covariant derivative<sup>2</sup><sup>2</sup>2We use the metric $`\eta ^{ab}=(,+,+,+,+)`$, $`\{\mathrm{\Gamma }^a,\mathrm{\Gamma }^b\}=2\eta ^{ab}`$, $`𝒟_\mu =_\mu +\frac{1}{4}\omega _{\mu ab}\mathrm{\Gamma }^{ab}`$, $`\omega _{\mu ab}`$ is the spin connection, and $`\mathrm{\Gamma }^{a_1a_2\mathrm{}a_n}=\frac{1}{n!}\mathrm{\Gamma }^{[a_1}\mathrm{\Gamma }^{a_2}\mathrm{}\mathrm{\Gamma }^{a_n]}`$..
## 3 Supersymmetric String Solutions
As a general ansatz for the supersymmetric solutions we consider metrics of the form
$$ds^2=e^{2V}dt^2+e^{2T}dz^2+e^{2U}dr^2+F(r)^2d\sigma ^2,$$
(5)
where $`V,T,U`$ are functions of $`r`$ only, and we consider either $`F(r)=r`$ or $`F(r)=R=`$ constant. $`d\sigma ^2`$ denotes the metric of a two-manifold $`𝒮`$ of constant Gaussian curvature $`k`$. Without loss of generality we restrict ourselves to the cases $`k=0,\pm 1`$. Clearly $`𝒮`$ is a quotient space of the universal coverings $`S^2`$ ($`k=1`$), $`H^2`$ ($`k=1`$) or $`E^2`$ ($`k=0`$). Explicitly, we choose
$$d\sigma ^2=d\theta ^2+f(\theta )^2d\varphi ^2,$$
(6)
where
$$f(\theta )=\{\begin{array}{cc}\mathrm{sin}\theta ,\hfill & k=1,\hfill \\ 1,\hfill & k=0,\hfill \\ \mathrm{sinh}\theta ,\hfill & k=1.\hfill \end{array}$$
(7)
The case $`k=1`$, $`F=r`$ has been considered in . There, a spherically symmetric magnetic string solution was obtained, which contains a naked singularity.
For the metric (5), the fünfbein and its inverse can be chosen as
$`e_t^0`$ $`=`$ $`e^V,e_z^1=e^T,e_r^2=e^U,e_\theta ^3=F,e_\varphi ^4=Ff,`$
$`e_0^t`$ $`=`$ $`e^V,e_1^z=e^T,e_2^r=e^U,e_3^\theta ={\displaystyle \frac{1}{F}},e_4^\varphi ={\displaystyle \frac{1}{Ff}}.`$ (8)
The nonvanishing components of the spin connection are given by
$`\omega _t^{02}`$ $`=`$ $`V^{}e^{VU},`$
$`\omega _z^{12}`$ $`=`$ $`T^{}e^{TU},`$
$`\omega _\theta ^{23}`$ $`=`$ $`F^{}e^U,`$
$`\omega _\varphi ^{24}`$ $`=`$ $`F^{}fe^U,`$
$`\omega _\varphi ^{34}`$ $`=`$ $`f^{}.`$ (9)
In five dimensions, strings can carry magnetic charges under the one-form potentials $`A^I`$, so we assume that the gauge fields have only a magnetic part, i. e.
$$F_{\theta \varphi }^I=kq^If(\theta ),A_\varphi ^I=kq^If(\theta )𝑑\theta .$$
(10)
Furthermore, we are looking for solutions with constant moduli $`X^I`$, which are chosen to minimize the magnetic central charge $`Z=q^IX_I`$, as in the case of the double extreme solutions in the ungauged theory . This means that one has
$$_i(Z)=_i(q^IX_I)=\frac{1}{3}C_{IJK}X^I_i(X^J)q^K=0.$$
(11)
Moreover, we make the choice
$$X^IV_I=1.$$
(12)
Using (11) and (12), the gaugino transformations (4) can be easily seen to vanish identically.
Plugging the spin connection (9) and the magnetic fields (10) into (3), we obtain for the supersymmetry variation of the gravitino
$`\delta \psi _t`$ $`=`$ $`\left(_t+{\displaystyle \frac{1}{2}}V^{}e^{VU}\mathrm{\Gamma }_{02}+{\displaystyle \frac{ik}{4}}Z{\displaystyle \frac{e^V}{F^2}}\mathrm{\Gamma }_{034}+{\displaystyle \frac{g}{2}}e^V\mathrm{\Gamma }_0\right)ϵ,`$
$`\delta \psi _z`$ $`=`$ $`\left(_z+{\displaystyle \frac{1}{2}}T^{}e^{TU}\mathrm{\Gamma }_{12}+{\displaystyle \frac{ik}{4}}Z{\displaystyle \frac{e^T}{F^2}}\mathrm{\Gamma }_{134}+{\displaystyle \frac{g}{2}}e^T\mathrm{\Gamma }_1\right)ϵ,`$
$`\delta \psi _r`$ $`=`$ $`\left(_r+{\displaystyle \frac{ike^U}{4}}{\displaystyle \frac{Z}{F^2}}\mathrm{\Gamma }_{234}+{\displaystyle \frac{g}{2}}e^U\mathrm{\Gamma }_2\right)ϵ,`$ (13)
$`\delta \psi _\theta `$ $`=`$ $`\left(_\theta {\displaystyle \frac{1}{2}}F^{}e^U\mathrm{\Gamma }_{23}{\displaystyle \frac{ikZ}{2F}}\mathrm{\Gamma }_4+{\displaystyle \frac{g}{2}}F\mathrm{\Gamma }_3\right)ϵ,`$
$`\delta \psi _\varphi `$ $`=`$ $`\left(_\varphi {\displaystyle \frac{1}{2}}F^{}fe^U\mathrm{\Gamma }_{24}{\displaystyle \frac{1}{2}}f^{}\mathrm{\Gamma }_{34}+{\displaystyle \frac{ikZf}{2F}}\mathrm{\Gamma }_3+{\displaystyle \frac{g}{2}}Ff\mathrm{\Gamma }_4{\displaystyle \frac{3i}{2}}gV_Iq^Ik{\displaystyle f𝑑\theta }\right)ϵ.`$
In what follows, we will consider three different cases.
### 3.1 Solutions With Flat Transverse Space
Let us first consider the case $`k=0`$, $`F=r`$, i. e. flat transverse space. Our choice (10) implies that we have vanishing gauge fields for $`k=0`$. The vanishing of the gravitino supersymmetry transformations (13) then yields the Killing spinor equations
$`\left(_t+{\displaystyle \frac{1}{2}}e^V\mathrm{\Gamma }_0(V^{}e^U\mathrm{\Gamma }_2+g)\right)ϵ`$ $`=`$ $`0,`$
$`\left(_z+{\displaystyle \frac{1}{2}}e^T\mathrm{\Gamma }_1(T^{}e^U\mathrm{\Gamma }_2+g)\right)ϵ`$ $`=`$ $`0,`$
$`\left(_r+{\displaystyle \frac{g}{2}}e^U\mathrm{\Gamma }_2\right)ϵ`$ $`=`$ $`0,`$
$`\left(_\theta +{\displaystyle \frac{1}{2}}\mathrm{\Gamma }_3(e^U\mathrm{\Gamma }_2+gr)\right)ϵ`$ $`=`$ $`0,`$
$`\left(_\varphi +{\displaystyle \frac{1}{2}}\mathrm{\Gamma }_4(e^U\mathrm{\Gamma }_2+gr)\right)ϵ`$ $`=`$ $`0.`$ (14)
From the integrability conditions of these equations one gets
$$e^V=e^T=e^U=gr.$$
(15)
Plugging these results into (14) and introducing the projectors $`P_\pm =\frac{1}{2}(1\pm \mathrm{\Gamma }_2)`$, we obtain for the Killing spinors
$$ϵ=r^{\frac{1}{2}}ϵ_+^0gr^{\frac{1}{2}}(g\mathrm{\Gamma }_0t+g\mathrm{\Gamma }_1z+\mathrm{\Gamma }_3\theta +\mathrm{\Gamma }_4\varphi )ϵ_+^0+r^{\frac{1}{2}}ϵ_{}^0,$$
(16)
where $`ϵ_\pm ^0`$ are constant spinors satisfying $`P_{}ϵ_\pm =0`$. From (15) it is clear that the solution we found is locally $`AdS_5`$ (written in horospherical coordinates). (16) tells us that this spacetime, as it should be, is fully supersymmetric. However, one may wish to compactify the $`(\theta ,\varphi )`$ sector to a cylinder or a torus, considering thus a quotient space of $`AdS_5`$. In this case, the surviving Killing spinors are those which respect the identifications performed in the $`(\theta ,\varphi )`$ sector. These are
$$ϵ=r^{\frac{1}{2}}ϵ_{}^0,$$
(17)
so that the considered $`AdS_5`$ quotient space preserves half of the supersymmetries.
Note that the above supersymmetric string solution is a limiting case of a family of nonextremal black strings, whose metric is given by
$$ds^2=e^{2V}dt^2+r^2dz^2+e^{2V}dr^2+r^2(d\theta ^2+d\varphi ^2),$$
(18)
where
$$e^{2V}=\frac{m}{r^2}+g^2r^2,$$
(19)
$`m`$ denoting an integration constant related to the mass of the black string. Considering $`z`$ as a coordinate of transverse space, (18) clearly can also be interpreted as a black hole, it is the solution found in .
### 3.2 Hyperbolic Transverse Space
We turn now to the more interesting case of hyperbolic transverse space, i. e. $`k=1`$, $`F=r`$. As supersymmetry breaking conditions we take
$`\mathrm{\Gamma }_3\mathrm{\Gamma }_4ϵ`$ $`=`$ $`iϵ,`$
$`\mathrm{\Gamma }_2ϵ`$ $`=`$ $`ϵ.`$ (20)
Then, the transformations (3) reduce to
$`\delta \psi _t`$ $`=`$ $`\left(_t{\displaystyle \frac{1}{2}}(V^{}e^{VU}Z{\displaystyle \frac{e^V}{2r^2}}ge^V)\mathrm{\Gamma }_0\right)ϵ,`$
$`\delta \psi _z`$ $`=`$ $`\left(_z{\displaystyle \frac{1}{2}}(T^{}e^{TU}Z{\displaystyle \frac{e^T}{2r^2}}ge^T)\mathrm{\Gamma }_1\right)ϵ,`$
$`\delta \psi _r`$ $`=`$ $`\left(_r{\displaystyle \frac{e^U}{2}}({\displaystyle \frac{Z}{2r^2}}+g)\right)ϵ,`$ (21)
$`\delta \psi _\theta `$ $`=`$ $`\left(_\theta {\displaystyle \frac{1}{2}}(e^U+{\displaystyle \frac{Z}{r}}gr)\mathrm{\Gamma }_3\right)ϵ,`$
$`\delta \psi _\varphi `$ $`=`$ $`\left(_\varphi {\displaystyle \frac{i}{2}}(13gV_Iq^I)\mathrm{cosh}\theta {\displaystyle \frac{1}{2}}(e^U+{\displaystyle \frac{Z}{r}}gr)\mathrm{sinh}\theta \mathrm{\Gamma }_4\right)ϵ.`$
The vanishing of the above equations implies the following conditions on the supersymmetry spinor $`ϵ`$,
$`_tϵ`$ $`=`$ $`0,`$
$`_zϵ`$ $`=`$ $`0,`$
$`_\theta ϵ`$ $`=`$ $`0,`$
$`_\varphi ϵ`$ $`=`$ $`0,`$
$`3gq^IV_I`$ $`=`$ $`1,`$
$`e^U{\displaystyle \frac{Z}{r}}+gr`$ $`=`$ $`0,`$
$`e^UT^{}+{\displaystyle \frac{Z}{2r^2}}+g`$ $`=`$ $`0,`$
$`e^UV^{}+{\displaystyle \frac{Z}{2r^2}}+g`$ $`=`$ $`0.`$ (22)
The last two equations in (22) imply that one should set $`T=V`$. From the sixth equation of (22), one immediately obtains
$$e^U=\frac{Z}{r}+gr.$$
(23)
Using the last equation of (22), we obtain a differential equation for $`V`$,
$$V^{}e^U=g+\frac{Z}{2r^2}.$$
(24)
The above differential equation can be easily solved by noticing that it can be rewritten in the form
$$\frac{dV}{dr}=\frac{d}{dr}\mathrm{log}(e^U)\frac{1}{4}\frac{d}{dr}\mathrm{log}(\frac{e^U}{gr}).$$
(25)
(Recall that the central charge $`Z`$ takes a constant value to be determined). (25) yields the following solution for $`V`$,
$$e^V=e^{3\frac{U}{4}}(gr)^{\frac{1}{4}}.$$
(26)
Let us now return to the minimization condition (11) of the magnetic central charge. It implies that the critical values of $`X^I`$ and its dual are given by
$$X^I=\frac{q^I}{Z},X_I=\frac{1}{6}\frac{C_{IJK}q^Jq^K}{Z^2},$$
(27)
and thus the critical value of the magnetic central charge is
$$Z=\left(\frac{1}{6}C_{IJK}q^Iq^Jq^K\right)^{\frac{1}{3}}.$$
(28)
Using the conditions $`X^IV_I=1`$ and the fifth relation of (22), one obtains a generalized Dirac quantization condition
$$\left(\frac{1}{6}C_{IJK}q^Iq^Jq^K\right)^{\frac{1}{3}}=\frac{1}{3g}.$$
(29)
For the case of pure supergravity where only the graviphoton charge $`q^0`$ is present, one obtains $`q^0=1/(3g)`$. A similar condition was obtained in .
To summarize, the BPS magnetic black string solution to $`D=5`$, $`N=2`$ gauged supergravity coupled to vector multiplets is given by
$`ds^2`$ $`=`$ $`(gr)^{\frac{1}{2}}e^{\frac{3}{2}U}(dt^2+dz^2)+e^{2U}dr^2+r^2\left(d\theta ^2+\mathrm{sinh}^2\theta d\varphi ^2\right),`$
$`e^U`$ $`=`$ $`{\displaystyle \frac{1}{3gr}}+gr,`$ (30)
while the gauge fields and the scalars are
$`A_\varphi ^I`$ $`=`$ $`q^I\mathrm{cosh}\theta ,`$
$`X^I`$ $`=`$ $`3gq^I.`$ (31)
The Killing spinor is independent of the angular variables, and its radial dependence is obtained by solving for its radial differential equation, which reads
$$\left(_r\frac{e^U}{2}(g+\frac{1}{6gr^2})\right)ϵ=0.$$
(32)
Using the relation (24), the above differential equation can be written in the simple form
$$\left(_r\frac{1}{2}V^{}\right)ϵ=0,$$
(33)
and thus we get
$$ϵ(r)=e^{\frac{1}{2}V}ϵ_0,$$
(34)
where $`ϵ_0`$ is a constant spinor satisfying the constraints
$$\mathrm{\Gamma }_3\mathrm{\Gamma }_4ϵ_0=iϵ_0,\mathrm{\Gamma }_2ϵ_0=ϵ_0.$$
(35)
As the Killing spinors do not depend on the coordinates $`\theta ,\varphi `$ of the transverse hyperbolic space, one could also compactify the $`H^2`$ to a Riemann surface $`𝒮_n`$ of genus $`n`$, and the resulting solution would still preserve one quarter of supersymmetry.
Whereas the spherical BPS magnetic string found in contains a naked singularity, the hyperbolic black string (30) has an event horizon at $`r=r_+=1/(g\sqrt{3})`$. This is analogous to the $`AdS_4`$ case, where BPS magnetic black holes with hyperbolic event horizons have been found , whereas for spherical topology one gets supersymmetric naked singularities .
Note that the black strings (30) are solitonic objects in the sense that the limit $`g0`$ (we recall that $`g`$ is the coupling constant of the gauged theory, coupling the $`U(1)`$ vector fields to the fermions) does not exist.
In the near horizon region, (30) reduces to the product manifold $`AdS_3\times H^2`$. This is easily seen by introducing the new radial coordinate $`\rho =(rr_+)^{1/4}`$. In the next subsection, we will see that in the near horizon limit, supersymmetry is enhanced.
### 3.3 Hyperbolic Transverse Space and Constant Warping Function
Let us now consider the case $`k=1`$ and $`F=R`$, where $`R`$ is a constant. We also choose $`T=V`$. As supersymmetry breaking condition we take
$$\mathrm{\Gamma }_3\mathrm{\Gamma }_4ϵ=iϵ.$$
(36)
The Killing spinor equations following from (3) are then
$`\delta \psi _t`$ $`=`$ $`\left(_t+{\displaystyle \frac{1}{2}}V^{}e^{VU}\mathrm{\Gamma }_{02}+{\displaystyle \frac{1}{4}}Z{\displaystyle \frac{e^V}{R^2}}\mathrm{\Gamma }_0+{\displaystyle \frac{g}{2}}e^V\mathrm{\Gamma }_0\right)ϵ=0,`$
$`\delta \psi _z`$ $`=`$ $`\left(_z+{\displaystyle \frac{1}{2}}V^{}e^{VU}\mathrm{\Gamma }_{12}+{\displaystyle \frac{1}{4}}Z{\displaystyle \frac{e^V}{R^2}}\mathrm{\Gamma }_1+{\displaystyle \frac{g}{2}}e^V\mathrm{\Gamma }_1\right)ϵ=0,`$
$`\delta \psi _r`$ $`=`$ $`\left(_r+{\displaystyle \frac{1}{4}}Z{\displaystyle \frac{e^U}{R^2}}\mathrm{\Gamma }_2+{\displaystyle \frac{g}{2}}e^U\mathrm{\Gamma }_2\right)ϵ=0,`$ (37)
$`\delta \psi _\theta `$ $`=`$ $`\left(_\theta {\displaystyle \frac{Z}{2R}}\mathrm{\Gamma }_3+{\displaystyle \frac{g}{2}}R\mathrm{\Gamma }_3\right)ϵ=0,`$
$`\delta \psi _\varphi `$ $`=`$ $`\left(_\varphi {\displaystyle \frac{i}{2}}\mathrm{cosh}\theta {\displaystyle \frac{Z}{2R}}\mathrm{sinh}\theta \mathrm{\Gamma }_4+{\displaystyle \frac{gR}{2}}\mathrm{sinh}\theta \mathrm{\Gamma }_4+{\displaystyle \frac{3i}{2}}gV_Iq^I\mathrm{cosh}\theta \right)ϵ=0.`$
The integrability conditions for these equations imply that
$$3gV_Iq^I=1,$$
(38)
and that the central charge is related to the compactification radius by $`Z=gR^2`$. Furthermore, one obtains
$$_r(V^{}e^{VU})=\frac{9}{4}g^2e^{V+U},V^{}e^U=\frac{3}{2}g.$$
(39)
Defining a new radial coordinate $`\rho `$ by $`g\rho =e^V`$, one immediately sees that the three-dimensional part of spacetime is $`AdS_3`$ in horospherical coordinates. We have thus obtained a supersymmetric product space $`AdS_3\times H^2`$, i. e. the near-horizon geometry of (30). Plugging the relations (38) and $`Z=gR^2`$ into (37), one obtains that the Killing spinors are independent of $`\theta ,\varphi `$. The remaining system is solved by
$$ϵ=e^{\frac{1}{2}V}ϵ_+^0\frac{3}{2}ge^{\frac{1}{2}V}(\mathrm{\Gamma }_0t+\mathrm{\Gamma }_1z)ϵ_+^0+e^{\frac{1}{2}V}ϵ_{}^0,$$
(40)
where $`ϵ_\pm ^0`$ are constant spinors satisfying
$$(1\mathrm{\Gamma }_2)ϵ_\pm ^0=0,\mathrm{\Gamma }_3\mathrm{\Gamma }_4ϵ_\pm ^0=iϵ_\pm ^0.$$
(41)
The product space $`AdS_3\times H^2`$ is thus one half supersymmetric. This supersymmetry enhancement near the horizon of the BPS black string is analogous to the case of ungauged supergravity theories, where usually in the near-horizon limit supersymmetry is fully restored.
## 4 General Product Space Compactifications
In this section we consider general product space compactifications of gauged $`D=5`$, $`N=2`$ supergravity coupled to vector multiplets. Spacetime is assumed to be a product $`M_3\times M_2`$, where, as above, $`M_2`$ denotes a two-manifold of constant curvature. We are interested in the conditions imposed by supersymmetry on $`M_{2,3}`$. To this end, we perform a $`3+2`$ decomposition of the gamma matrices<sup>3</sup><sup>3</sup>3See the appendix of for a nice summary of gamma matrix decomposition in Kaluza-Klein compactifications. in the following way
$$\mathrm{\Gamma }^a=(\mathrm{\Gamma }^{\widehat{\alpha }},\mathrm{\Gamma }^{\widehat{ı}})=(\gamma ^{\widehat{\alpha }}\sigma ^3,1\mathrm{\Sigma }^{\widehat{ı}}),$$
(42)
where early Greek letters $`\alpha ,\beta ,\mathrm{}`$ are $`M_3`$ spacetime indices, and $`i,j,\mathrm{}`$ are $`M_2`$ spacetime indices. The hatted indices refer to the corresponding tangent spaces. The $`\gamma ^{\widehat{\alpha }}`$ and $`\mathrm{\Sigma }^{\widehat{ı}}`$ denote Dirac matrices in three and two dimensions respectively. To be concrete, we make the choice $`\mathrm{\Sigma }^1=\sigma ^2`$, $`\mathrm{\Sigma }^2=\sigma ^1`$, where the Pauli matrices are chosen to be
$$\sigma ^1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\sigma ^2=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),\sigma ^3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
(43)
The supersymmetry parameter $`ϵ`$ in five dimensions is decomposed as $`ϵ=\eta \chi `$. Note that $`\sigma ^3`$ plays the role of a chirality operator for the spinors $`\chi `$ in two dimensions. Some useful relations needed below are
$$\mathrm{\Gamma }^{\widehat{\alpha }\widehat{\beta }}=\gamma ^{\widehat{\alpha }\widehat{\beta }}1,\mathrm{\Gamma }^{\widehat{ı}\widehat{ȷ}}=1\mathrm{\Sigma }^{\widehat{ı}\widehat{ȷ}},\mathrm{\Gamma }^{\widehat{\alpha }}\mathrm{\Gamma }^{\widehat{ı}}\mathrm{\Gamma }^{\widehat{ȷ}}=\gamma ^{\widehat{\alpha }}\sigma ^3\mathrm{\Sigma }^{\widehat{ı}}\mathrm{\Sigma }^{\widehat{ȷ}}.$$
(44)
For the field strength of the abelian vectors we make the ansatz
$$F^I=q^Iϵ,$$
(45)
where $`ϵ`$ denotes the volume form on $`M_2`$<sup>4</sup><sup>4</sup>4We apologize for using the same symbol for the volume form and the supersymmetry parameter, but the meaning should be clear from the context.. The gaugino variation (4) vanishes as before, provided that (11) and $`V_IX^I=1`$ are satisfied. (Note that we still assume the moduli $`X^I`$ to be constant). Inserting the decomposition of the Dirac matrices and the supersymmetry parameter, as well as the ansatz for the field strength into the gravitino variation (3) yields the relations
$`\left(_\alpha +{\displaystyle \frac{1}{4}}\omega _\alpha ^{\widehat{\alpha }\widehat{\beta }}\gamma _{\widehat{\alpha }\widehat{\beta }}+{\displaystyle \frac{1}{4}}X_Iq^I\gamma _\alpha +{\displaystyle \frac{g}{2}}\gamma _\alpha \right)\eta \chi _+`$ (46)
$`+\left(_\alpha +{\displaystyle \frac{1}{4}}\omega _\alpha ^{\widehat{\alpha }\widehat{\beta }}\gamma _{\widehat{\alpha }\widehat{\beta }}+{\displaystyle \frac{1}{4}}X_Iq^I\gamma _\alpha {\displaystyle \frac{g}{2}}\gamma _\alpha \right)\eta \chi _{}=0`$
and
$$\left(_i+\frac{1}{4}\omega _i^{\widehat{ı}\widehat{ȷ}}\mathrm{\Sigma }_{\widehat{ı}\widehat{ȷ}}\frac{i}{2}X_Iq^Iϵ_{ij}\mathrm{\Sigma }^j+\frac{g}{2}\mathrm{\Sigma }_i\frac{3i}{2}gV_IA_i^I\right)\chi =0,$$
(47)
where $`\chi _\pm =\frac{1}{2}(1\pm \sigma ^3)\chi `$ denote the chirality projections of the two-dimensional spinor $`\chi `$. From (46) we get that either $`\chi _{}`$ and the bracketed expression in the first line have to vanish, or $`\chi _+`$ and the term in brackets in the second line are zero. Without loss of generality, we assume the first possibility. Plugging $`\chi _{}=0`$ into (47), one obtains
$$(g\mathrm{\Sigma }_iiX_Iq^Iϵ_{ij}\mathrm{\Sigma }^j)\chi _+=0.$$
(48)
In order to have nontrivial solutions to this equation, the determinant of $`g\mathrm{\Sigma }_iiZϵ_{ij}\mathrm{\Sigma }^j`$ has to vanish. This implies
$$Z=\pm g$$
(49)
for the magnetic central charge $`Z`$. One can easily show that the lower sign is incompatible with the condition $`\sigma ^3\chi _+=\chi _+`$, so the upper positive sign is chosen and (48) is then automatically satisfied. The remaining Killing spinor equation for $`\chi _+`$ reads
$$𝒟_i\chi _+=0,$$
(50)
with the gauge- and Lorentz-covariant derivative $`𝒟_i`$ given by
$$𝒟_i=_i+\frac{1}{4}\omega _i^{\widehat{ı}\widehat{ȷ}}\mathrm{\Sigma }_{\widehat{ı}\widehat{ȷ}}\frac{3i}{2}gV_IA_i^I.$$
(51)
The integrability condition for (50) is
$$[𝒟_i,𝒟_j]\chi _+=0,$$
(52)
or, equivalently,
$$(\frac{1}{4}R_{ij\widehat{k}\widehat{l}}\mathrm{\Sigma }^{\widehat{k}\widehat{l}}\frac{3i}{2}gV_Iq^Iϵ_{ij})\chi _+=0.$$
(53)
Taking into account that $`M_2`$ is of constant curvature, we have
$$R_{ijkl}=\frac{R}{2}(g_{ik}g_{jl}g_{il}g_{jk})$$
(54)
for the Riemann tensor of $`M_2`$. Using this in (53), one immediately obtains
$$R+6gV_Iq^I=0$$
(55)
for the scalar curvature $`R`$ of $`M_2`$. From $`q^IX_I=g`$ we have
$$X^I=\frac{q^I}{g}$$
(56)
for the moduli, and thus
$$R=6gV_Iq^I=6g^2V_IX^I=6g^2<0$$
(57)
for the scalar curvature $`R`$. This means that, to preserve some supersymmetry, $`M_2`$ must be diffeomorphic to hyperbolic space $`H^2`$ or to a quotient thereof.
Using the second Cartan equation, one obtains that the spin connection $`\omega ^{12}`$ on $`M_2`$ is related to the vector potential $`A^I`$ by
$$q^I\omega ^{12}=A^I\frac{R}{2}.$$
(58)
Using this, (50) reduces to $`_i\chi _+=0`$, so that $`\chi _+`$ is independent of the coordinates on $`M_2`$. The remaining equation to solve is the Killing spinor equation on $`M_3`$ for $`\eta `$, i. e.
$$\left(_\alpha +\frac{1}{4}\omega _\alpha ^{\widehat{\alpha }\widehat{\beta }}\gamma _{\widehat{\alpha }\widehat{\beta }}+\frac{3}{4}g\gamma _\alpha \right)\eta =0.$$
(59)
The integrability conditions for (59) yield that $`M_3`$ must be an Einstein space with cosmological constant $`\mathrm{\Lambda }=(3g/2)^2`$. As we are in three dimensions, this means that $`M_3`$ is also of constant curvature, i. e. a quotient space of $`AdS_3`$. Note that the chirality condition $`\chi =\chi _+`$ breaks half of supersymmetry. The amount of supersymmetry preserved by $`M_3\times M_2`$ is then determined by the solutions of (59). If $`M_3=AdS_3`$, then the whole solution $`AdS_3\times H^2`$ is half supersymmetric, in agreement with what we found above. However, we can also choose $`M_3`$ to be e. g. the BTZ black hole. If we take the extremal rotating BTZ black hole, which preserves one half of the $`AdS_3`$ supersymmetries , then the solution BTZ$`{}_{extr}{}^{}\times H^2`$ preserves one quarter of the supersymmetries. We would like to point out that BTZ$`\times H^2`$ compactifications of $`D=5`$ anti-de Sitter gravity without gauge fields were previously considered in . However, as we showed above, due to the relation $`Z=g`$ these configurations cannot be supersymmetric unless some gauge fields are turned on.
## 5 Summary and Discussion
To sum up, we presented supersymmetric string solutions of gauged $`D=5`$, $`N=2`$ supergravity coupled to abelian vector multiplets. The main result is the construction of a BPS black string with hyperbolic transverse space, preserving one quarter of supersymmetry. The curvature of the $`H^2`$ is supported by a nonvanishing field strength of the vector fields. The black strings are thus magnetically charged. In the near-horizon limit, their geometry approaches the half-supersymmetric product space $`AdS_3\times H^2`$, so we encounter supersymmetry enhancement near the horizon. This behaviour is similar to the case of ungauged supergravity theories. Note however, that in the ungauged case, usually supersymmetry is fully restored near the event horizon.
As the near horizon geometry contains an $`AdS_3`$ factor, it should be possible to use the $`AdS_3`$ asymptotic symmetry algebra in order to count the microstates yielding the Bekenstein-Hawking entropy of the extremal black string. This was done in for the BTZ black hole, and subsequently generalized in to higher-dimensional black holes containing a BTZ factor near the horizon. In our case, a similar procedure is hindered by the fact that we get the $`M=J=0`$ BTZ black hole in the near-horizon limit, so using Cardy’s formula we would obtain an incorrect result for the entropy. Similar difficulties have been encountered in , where a state counting for extremal black strings in three dimensions was performed. This suggests that a similar approach to that in must be used in our case, in order to overcome the above mentioned difficulties.
It is also of interest to investigate the role of the BPS magnetic black strings in the AdS/CFT correspondence. Note that the $`U(1)^3`$ truncation of gauged $`D=5`$, $`N=2`$ supergravity can be embedded into type IIB supergravity . This means that our solutions can be lifted to ten dimensions, with an internal five-sphere that is distorted by the one-form gauge fields $`A_\mu ^I`$. That breaks the isometry group $`SO(6)`$ of the $`S^5`$ down to a smaller subgroup. In the dual CFT, which is $`𝒩=4`$ SYM on $`^2`$ $`\times `$ $`H^2`$(or $`R\times `$ $`S^1\times H^2`$, if the coordinate $`z`$ parametrizes a compact space), the $`S^5`$ isometry group becomes the R-symmetry. On the CFT side, we are now dealing with the presence of nonvanishing background $`U(1)^3`$ currents, which break the global $`SO(6)`$ R-symmetry. In principle, it should also be possible to count the microstates giving rise to the black string entropy using the dual SYM theory on $`^2\times H^2`$ in the presence of these global background $`U(1)^3`$ currents. As the near-horizon geometry is $`AdS_3\times H^2`$, the presented magnetic black string solutions may also have a holographic interpretation in the sense that the four-dimensional field theory discussed above flows to a two-dimensional CFT in the IR. These issues are currently under investigation .
Finally, in the present paper, we also investigated $`3+2`$ product compactifications, and showed that only if the internal two-manifold is diffeomorphic to $`H^2`$, some amount of supersymmetry can be preserved. As an example we found the one quarter supersymmetric BTZ$`{}_{extr}{}^{}\times H^2`$ configuration, where BTZ<sub>extr</sub> denotes the extremal rotating BTZ black hole. Considering the equations of motion following from (1), one easily sees that one also can have the nonextremal BTZ black hole tensored by $`H^2`$ as a solution in presence of magnetic gauge fields. Presumably, these configurations arise as the near-horizon limit of the nonextremal generalization of (30).
## Acknowledgements
The authors would like to thank K. Khuri-Makdisi and A. Zaffaroni for useful discussions. |
warning/0001/quant-ph0001067.html | ar5iv | text | # Approach to the semiconductor cavity QED in high-Q regimes with q-deformed boson
## I Introduction
It’s well known that the exciton system is a quasi-particle system. At the low density, excitons are approximately treated as bosons which obey Bose statistics . But when the density of the excitons become higher, the excitons, which somewhat deviate ideal bosons, are no longer ideal bosons. There are two ways of dealing with this problem : one way is to put these deviations into the effective interaction between the hypothetical ideal bosons and the exciton operators are still presented by the bosonic operators . Another way is the implement of the atomic operators model . The question naturelly arises whether the exciton system is equivalent to the atomic system.
In this paper we show that the exciton system could be described by $`q`$ deformed boson system which interpolates between Bose system and Fermi system ; and the deformation parameter $`q`$ is well defined by the total atomic particle number $`N,`$rather than it is phenomenological in the previous discussions .The concept of the q-deformed boson was even extensively applied due to the q-deformed boson realization of quantum group theory by different authors ten years ago . Since then, many physicists make efforts to find its real physical realizations. For example, they give some phenomenological investigations to fit the deformed spectra of rotation and oscillation for molecules and nuclei . In our opinion, those investigations can be regarded as merely phenomenological because a $`q`$-deformed structure is postulated in advance without giving it a microscopic mechanism. In this paper it will be shown that a physical and natural realization of the $`q`$-deformed boson is provided by the exciton operators, which was presented recently by Gardiner for the description of Bose-Einstein condensation (BEC). In fact, the similar quasiparticles scheme for particle-number conservation has already been introduced by Girardeau and Arnowitt almost 40 years ago . The relationship between Gardiner’s phonons and these quasiparticles has been discussed in a recent comment .
Following these ways, we find that the exciton system also could be described by the q-deformed boson. When the density of the exciton (the particles excited in upper state) is low enough, it return to the ideal boson case. Using this theory, we could give a good explanation on the semiconductor cavity QED in high-Q regimes.What will be investigated here is the case that the total atomic particle number $`N`$ is very large but not infinite. That is, we shall consider the effects of order $`o(1/N)`$. And we shall focus on an algebraic method, a q-deformed boson algebra, of treating the effects of finite particle number. As it turns out, the commutation relations for the exciton operators will no longer obey the commutation relation of the Heisenberg-Weyl algebra but the $`q`$-deformed bosonic commutation relation
$$[b_q,b_q^{}]_qb_qb_q^{}qb_q^{}b_q=1,$$
(1)
where the deformation constant $`q`$ depends on the total atomic particle number.
This paper is organized as follows. In section 2 we firstly deduce the q-deformed commutation relation for the exciton in the high-Q cavity in case of the large but finite lattice molecule number $`N`$ . In section 3, only keeping the first order term of $`\frac{1}{N},`$we model the Frenkel excitons in a micro-cavity as the dressed q-deformed boson system. In section 4, the quantum approach for angular momentum is used to obtain the eigen- values and eigen-function of the system under the one order approximation. The stationary physical spectrum of the system is calculated in the section 5. Finally we summarize our results with some comments.
## II $`q`$-deformed bosonic algebra for exciton
Gardiner’s starting point to introduce the exciton operators is to consider a system of the weakly interacting Bose gas. Without losing generality, we consider a thin molecular crystal film containing $`N`$ identical two-level molecules interacting resonantly with a single mode quantum field. The intermolecular interaction is neglected. We assume that all molecules have equivalent mode positions, so they have the same coupling constant $`\kappa `$. By using Dick model , we could write the Hamiltonian under the rotating wave approximation as following:
$$H=\mathrm{}\mathrm{\Omega }(S_z+a^+a)+\mathrm{}\kappa (aS_++a^+S_{}),$$
(2)
where, $`a`$ is annihilation operator of the quantum cavity field and
$$\begin{array}{ccc}S_Z\hfill & =\hfill & _{n=1}^Ns_z(n),\hfill \\ S_+\hfill & =\hfill & _{n=1}^Ns_+(n),\hfill \\ S_{}\hfill & =\hfill & _{n=1}^Ns_{}(n),\hfill \end{array}$$
(3)
where, $`s_z(n)=\frac{1}{2}(|e_n><e_n||g_n><g_n|)`$, $`s_+(n)=|e_n><g_n|`$ and $`s_{}(n)=|g_n><e_n|`$ are quasi spin operators of the nth molecule. $`|e_n>`$ and $`|g_n>`$ are the excited state and the ground state of n’th molecule.
Consider the second quantization of the above model. Let $`b_e^{}`$ and $`b_e`$ denote the creation and annihilation operators for the atoms in the excited state and $`b_g^{}`$ and $`b_g`$ for the creation and annihilation operators of the atoms in the ground state. The simplified Hamiltonian in second quantization reads
$$H=\mathrm{}\mathrm{\Omega }(b_e^{}b_eb_g^{}b_g+a^+a)+\mathrm{}\kappa [ab_e^{}b_g+H.\mathrm{c}.].$$
(4)
Note that the total atomic particle number $`𝐍=b_e^{}b_e+b_g^{}b_g`$ is conserved. For convenience we define $`\eta =1/N`$ for large particle number.
In the thermodynamical limit $`N\mathrm{}`$, the Bogoliubov approximation is usually applied, in which the ladder operators $`b_g^{},b_g`$ of the ground state are replaced by a $`c`$-number $`\sqrt{N_c}`$, where $`N_c`$ is the average number of the ininital condensated atoms. As a result Hamiltonian Eq.(4) becomes a two-coupling harmonic oscillator system
$$H_b=\mathrm{}\mathrm{\Omega }(b_e^{}b_e+a^+a)+\mathrm{}\kappa \sqrt{N_c}[ab_e^{}+H.\mathrm{c}.].$$
(5)
However, this apporoximation destroyes a symmetry of the Hamiltonian Eq.(4), i.e., the conservation of the total particle number is violated because of $`[N,H_b]0`$.
To avoide this problem, the exciton operators are defined as:
$$b_q=\frac{1}{\sqrt{N}}b_g^{}b_e,b_q^{}=\frac{1}{\sqrt{N}}b_gb_e^{}.$$
(6)
according to Gardiner.These operators act invariantly on the subspace $`V^N`$ spanned by bases $`|N;n|Nn,n`$ $`(n=0,1,\mathrm{},N)`$, where Fock sates $`(m,n=0,1,2,\mathrm{})`$
$`|m,n={\displaystyle \frac{1}{\sqrt{m!n!}}}b_e^mb_g^n|0`$
spann the Fock space $`H_{2b}`$ of a two mode boson.
A straightforward calculation leads to the following commutation relation between the exciton operaotr and its Hermitian conjugate:
$$[b_q,b_q^{}]=1\frac{2}{N}b_e^{}b_e=f(b_q^{}b_q;\eta ),$$
(7)
with $`f(x;\eta )=\sqrt{1+2(12x)\eta +\eta ^2}\eta `$. Keeping only the lowest order of $`\eta `$ for a very large total particle number, the commutator above becomes
$$[b_q,b_q^{}]=12\eta b_q^{}b_q$$
(9)
or
$$[b_q,b_q^{}]_q=1,$$
(10)
with $`q=12\eta `$. This is exactly a typical $`q`$-deformed commutation relation. As $`N\mathrm{}`$ or $`q1`$, the usual commutation relation of Heisenberg-Weyl algebra is regained.
In the above discussion about the phonon excitation, we have linearized commutator $`hf(b^{}b;\eta )`$ so that a $`q`$-deformed commutation rule was obtained. Essentially this linearization establishes a physical realization of the $`q`$-deformed algebra. However, if the total particle number $`N`$ is not large enough, then $`h`$ can not be approximated by a linear function. From the commutation relations between $`h`$ and $`b_q,b_q^{}`$
$$[h,b_q^{}]=\frac{2}{N}b_q^{},[h,b_q]=\frac{2}{N}b_q,$$
(11)
we see that the algebra of exciton operators is a rescaling of algebra $`su(2)`$ with factor $`N`$.
## III Theoretical Model
Based on the above analysis about the algebraic structure of exciton operator, we consider the case of the low density of atoms in excited sate for the Hamiltonian (2).
Since the second quantization forms of $`S_+`$ and $`S_{}`$ are $`S_+=b_e^{}b_g,S_{}=b_g^{}b_e`$, it is straighforward to prove that the collective operator $`\frac{S_+}{\sqrt{N}}`$ and $`\frac{S_{}}{\sqrt{N}}`$ are approximately considered as the simple bosonic operators as $`N\mathrm{}.`$ These collective operators are called exciton operators. But in case of the high density of molecules in excited state with finite $`N`$, many molecules are in the excited states, the bosonic approximation could no longer work well. The Hamiltonian (2) is rewritten as the effective Hamiltonian in the form of q-deformed boson:
$$H=\mathrm{}\mathrm{\Omega }(a^+a+b_q^+b_q)+\mathrm{}g(a^+b_q+b_q^+a)$$
(12)
with $`g=\sqrt{N}k`$, $`b_q`$ and $`b_q^+`$ satisfy q-deformed relation:
$$[b_q,b_q^+]=b_qb_q^+qb_q^+b_q=1,$$
(13)
where,
$$q=1\frac{2}{N}$$
(14)
Here, the deformation parameter $`q`$ is determined by the lattice molecule number. So $`q`$ is no longer phenomenological.
Up to the first order approximation, $`b_q^+(b_q)`$ could be expressed as following
$`b_q^+`$ $`=`$ $`b^++{\displaystyle \frac{b^+b^+b}{2N}},`$ (15)
$`b_q`$ $`=`$ $`b+{\displaystyle \frac{b^+bb}{2N}}.`$ (16)
in terms of the general bosonic operators $`b^+(b)`$ .So the Hamitonian $`H`$ is rewritten in form of perturbation
$$H=H_0+H^{}$$
(17)
where
$`H_0`$ $`=`$ $`\mathrm{}\mathrm{\Omega }(a^+a+b^+b)+\mathrm{}g(a^+b+b^+a),`$ (18)
$`H^{}`$ $`=`$ $`{\displaystyle \frac{\mathrm{}}{2N}}(2\mathrm{\Omega }b^+b^+bb+gb^+b^+ab+a^+b^+bb).`$ (19)
It is clearly that $`H^{}`$ is equivalent to the attractive exciton-exciton collisions due to the bi-exciton effect and decreased exciton-photon coupling constants due to the phase phase filling effect .
## IV Approximate analytical solutions
To solve the Schroedinger equation governed by eq.(15),we implement the quantum angular momentum theory
If we define the angular momentum operators
$`J_z`$ $`=`$ $`{\displaystyle \frac{1}{2}}(a^+ab^+b),`$ (20)
$`J_+`$ $`=`$ $`a^+b,J_{}=ab^+.`$ (21)
then
$`J_x`$ $`=`$ $`{\displaystyle \frac{1}{2}}(a^+b+ab^+),`$ (22)
$`J_y`$ $`=`$ $`{\displaystyle \frac{1}{2i}}(a^+bab^+).`$ (23)
We rewrite the Hamitonian (5)
$$H_0=\mathrm{}\mathrm{\Omega }\widehat{N}+2\mathrm{}gJ_x=\mathrm{}\mathrm{\Omega }\widehat{N}+2\mathrm{}ge^{i\frac{\pi }{2}J_y}J_ze^{i\frac{\pi }{2}J_y}.$$
(24)
In terms of a $`SO(3)`$ rotation $`\mathrm{}(\mathrm{\Omega }\widehat{N}+2gJ_x)`$ by $`e^{i\frac{\pi }{2}J_y}`$. Note that the excitation number operator $`\widehat{N}=a^+a+b^+b`$ is a constant under the a $`SO(3)`$ rotation and
$$J^2=J_x^2+J_y^2+J_z^2=\frac{\widehat{N}}{2}(\frac{\widehat{N}}{2}1)$$
(25)
is the total angular momentum operator. In terms of the eigen states of $`J^2`$ and $`J_z`$
$$|jm>=\frac{(a^+)^{j+m}(b^+)^{jm}}{\sqrt{(j+m)!(jm)!}}|0>,$$
(26)
where the eigen values of the $`J^2`$ and $`J_z`$ are
$$j=\frac{𝒩}{2},m=\frac{𝒩}{2},\mathrm{},\frac{𝒩}{2}$$
(27)
the eigen functions and the eigen values of $`H_0`$ are constructed as
$$\psi _{jm}^0=e^{i\frac{\pi }{2}J_y}|jm>,E_{jm}^{(0)}=\mathrm{}\mathrm{\Omega }𝒩+2\mathrm{}gm.$$
(28)
Up to the first order approximation, the eigen values of $`H`$ are obtained as
$$E_{jm}=E_{jm}^{(0)}+<mj||e^{i\frac{\pi }{2}J_y}H^{}e^{i\frac{\pi }{2}J_y}|jm>,$$
(29)
whith their corresponding eigen functions
$$\psi _{jk}=\psi _{jk}^{(0)}+\underset{nk}{}\frac{H_{nk}^{}}{E_{jk}^{(0)}E_{jn}^{(0)}}\psi _{jn}^{(0)}.$$
(30)
where,$`n`$ and $`k`$ present the subscript $`(jm^{})`$ and $`(jm)`$. We calculate the matrix elements of the pretubation Hamiltonian $`H^{}`$:
$`<`$ $`m^{}j||e^{i\frac{\pi }{2}J_y}H^{}e^{i\frac{\pi }{2}J_y}|jm>`$ (31)
$`=`$ $`{\displaystyle \frac{\mathrm{}}{4N}}\mathrm{\Omega }\sqrt{(j+m)(j+m1)}`$ (32)
$`\times `$ $`\sqrt{(jm+1)(jm+2)}\delta _{m2,m^{}}`$ (33)
$`+`$ $`{\displaystyle \frac{\mathrm{}}{4N}}\mathrm{\Omega }\sqrt{(j+m+1)(j+m+2)}`$ (34)
$`\times `$ $`\sqrt{(jm)(jm1)}\delta _{m+2,m^{}}`$ (35)
$`+`$ $`{\displaystyle \frac{\mathrm{}}{4N}}(2\mathrm{\Omega }g)\sqrt{(jm)(j+m+1)}`$ (36)
$`\times `$ $`(jm1)\delta _{m+1,m^{}}`$ (37)
$`+`$ $`{\displaystyle \frac{\mathrm{}}{4N}}(2\mathrm{\Omega }+g)\sqrt{(jm)(j+m+1)}`$ (38)
$`\times `$ $`(j+m)\delta _{m+1,m^{}}`$ (39)
$`+`$ $`{\displaystyle \frac{\mathrm{}}{4N}}(2\mathrm{\Omega }+g)\sqrt{(j+m)(jm+1)}`$ (40)
$`\times `$ $`(j+m1)\delta _{m1,m^{}}`$ (41)
$`+`$ $`{\displaystyle \frac{\mathrm{}}{N}}(2\mathrm{\Omega }g)\sqrt{(j+m)(jm+1)}`$ (42)
$`\times `$ $`(jm)\delta _{m1,m^{}}`$ (43)
$`+`$ $`{\displaystyle \frac{\mathrm{}}{4N}}(\mathrm{\Omega }+g)(j+m)(j+m1)\delta _{m,m^{}}`$ (44)
$`+`$ $`{\displaystyle \frac{\mathrm{}}{4N}}(\mathrm{\Omega }g)(jm)(jm1)\delta _{m,m^{}}`$ (45)
$`+`$ $`{\displaystyle \frac{\mathrm{}}{N}}\mathrm{\Omega }(j^2m^2)\delta _{m,m^{}}`$ (46)
so the eigen values of $`H`$ are
$`E_{jm}`$ $`=`$ $`\mathrm{}\mathrm{\Omega }𝒩+2m\mathrm{}g+{\displaystyle \frac{\mathrm{}}{N}}\mathrm{\Omega }(j^2m^2)`$ (47)
$`+`$ $`{\displaystyle \frac{\mathrm{}}{4N}}(\mathrm{\Omega }+g)(j+m)(j+m1)`$ (48)
$`+`$ $`{\displaystyle \frac{\mathrm{}}{4N}}(\mathrm{\Omega }g)(jm)(jm1)`$ (49)
In general, we could obtain all eigen functions of $`H`$ under one order approximation by using eq.(24), eq.(26), and eq.(27). So the time evolution operator of the system is written as:
$$U(t)=exp(i\frac{H}{\mathrm{}}t)=\underset{j=0}{}\underset{m=j}{\overset{j}{}}exp(i\frac{E_{jm}}{\mathrm{}}t)|\psi _{jm}><\psi _{jm}|$$
(50)
## V Fluorescence spectrum of the exciton
We firstly give an analytic expression for the physical spectrum of the q-deformed exciton in terms of the Fock state of the quantum field and the exciton. The standard definition of the physical spectrum is
$$S(\omega )=2\gamma _0^tdt_1_0^tdt_2e^{(\gamma i\omega )(tt_2)}e^{(\gamma +i\omega )(tt_1)}G(t_1,t_2)$$
(51)
where $`\gamma `$ is the half-bandwith of spectrometer which is being used to measure the spectrum, and $`t`$ is time length of the excitation in the cavity. $`G(t_1,t_2)`$ is dipole correlation function, and
$$G(t_1,t_2)=<i|U^+(t_2)b_q^+U(t_2)U^+(t_1)b_qU(t_1)|i>,$$
(52)
where, $`|i>`$ is any initial state of the system. By using the bosonic approximation and substituting eq.(29) into eq.(31), the dipole correlation function is rewritten as following:
$`G(t_1,t_2)=<i|U^+(t_2)b_q^+U(t_2)U^+(t_1)b_qU(t_1)|i>`$ (53)
$`=`$ $`<i|U^+(t_2)(b^++{\displaystyle \frac{b^+b^+b}{2N}})U(t_2)U^+(t_1)(b+{\displaystyle \frac{b^+bb}{2N}}))U(t_1)|i>`$ (54)
$`=`$ $`{\displaystyle \underset{j,k,l,m,n,}{}}<i|\psi _{jl}><\psi _{jl}|(b^++{\displaystyle \frac{b^+b^+b}{2N}})|\psi _{km}>`$ (56)
$`\times <\psi _{km}|(b+{\displaystyle \frac{b^+bb}{2N}}))|\psi _{jn}><\psi _{jn}|i>e^{i\omega _{lm}t_2}e^{i\omega _{nm}t_1}`$
where $`\omega _{lm}=\frac{E_{jl}E_{km}}{\mathrm{}}`$ and $`\omega _{nm}=\frac{E_{jn}E_{km}}{\mathrm{}}`$. It’s evident that $`j`$ is determined only by the initial state $`|i>`$. So we have
$`S(\omega )`$ $`=`$ $`{\displaystyle \underset{l,k,m}{}}{\displaystyle \frac{2\gamma }{\gamma ^2+(\omega \omega _{lm})^2}}|<i|\psi _{jl}>|^2`$ (57)
$`|`$ $`<`$ $`\psi _{jl}|(b^++{\displaystyle \frac{b^+b^+b}{2N}})|\psi _{km}>|^2`$ (58)
Noting that we have passed the transient terms and slowly variation terms. This equation gives the stationary physical spectrum in terms of the system eigenvalues and eigenstaes. If $`<m^{}j^{}|b^+|jm>0`$, then we have $`j^{}=j+\frac{1}{2}`$ and $`m^{}=m\frac{1}{2}`$. So the equation (25) is rewritten as:
$`S(\omega )`$ $`=`$ $`{\displaystyle \underset{l,m}{}}{\displaystyle \frac{2\gamma }{\gamma ^2+(\omega \omega _{lm})^2}}|<i|\psi _{jl}>|^2`$ (59)
$`|`$ $`<`$ $`\psi _{jl}|(b^++{\displaystyle \frac{b^+b^+b}{2N}})|\psi _{(j\frac{1}{2})m}>|^2`$ (60)
The eigenvalues determine the position of the spectral component and $`|<i|\psi _{jl}>|^2|<\psi _{jl}|(b^++\frac{b^+b^+b}{2N})|\psi _{(j\frac{1}{2})m}>|^2`$ determine the intensity of the spectral lines.
In terms of the experimental condition of the reference , the bare excitons could be prepared by resonant femtosecond pulse pumping. If we prepare the initial state being in $`𝒩=1`$, then the eq.(34) shows that the emission spectrum of the $`𝒩=1`$ to the $`𝒩=0`$ transition has double peaks structure which is exactly equal to that of the two-level atomic system. When the pumping power is increased , the emission spectrum is quiet different from the case of the two-level atomic system. For example, if the system initially is in Fock state $`𝒩=2`$ , then we have:
$`S(\omega )`$ $`=`$ $`{\displaystyle \underset{l,m}{}}{\displaystyle \frac{2\gamma }{\gamma ^2+(\omega \omega _{lm})^2}}|<i|\psi _{1l}>|^2`$ (61)
$`|`$ $`<`$ $`\psi _{1l}|(b^++{\displaystyle \frac{b^+b^+b}{2N}})|\psi _{\frac{1}{2}m}>|^2`$ (62)
From this equation we know when $`𝒩=2`$ there are six peaks in the emission spectrum which are expected. Although there are three different initial state, they have similar spectrum shape. As is shown in Fig.1, this sextet structure is contrast to the triplet structure in the emission spectrum from strong pumped two-level system .
When the molecular number of the system is increased, Such as there are 10000 molecular in the system, the other conditions are the same as that of the Fig.1, then the coupling between the molecular and the cavity field is weak. There are two peaks in the emission spectrum (Fig. 2). In this case Bose approximation is good.
## VI Conclusion
It has been shown that the higher density Frenkel exciton coupling to a single mode high-Q microcavity field can be described by the quantum dynamics for the dressed q-deformed boson. Keeping the first order term of $`\frac{1}{N},`$the high density Frenkel exciton naturally obyes the q-deformed commutation relation. Based on this observation the quantum theory of angular momentum is used to obtain the eigen- values and eigen-function of the exciton system under the one order approximation. Comparing with the usual approach for Frenkel exciton dynamics our Hamiltonian is Hermitian and closed in form. An analytical expression for the stationaryphysical spectrum for the exciton is obtained by using of the dressed states of the cavity field and the exciton.
Acknowledgment
This work is supported by the NSF of China and ”9-5 Pandeng ” project. |
warning/0001/hep-ex0001045.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Intensity interferometry was applied in 1953 by Hanbury-Brown and Twiss in radio astronomy in order to estimate the spatial extension of stars (HBT effect). In particle reactions which lead to multi-hadronic final states the HBT effect manifests itself as a constructive interference between two identical bosons — the so-called Bose-Einstein (BE) correlation, which is now well known and was first observed by Goldhaber et al. in $`\overline{\mathrm{p}}`$p annihilations. There is interest in the quantum mechanical aspects of the BE correlations, but they are also used to estimate the dimensions of the source of the identical bosons. BE correlation studies of pion pairs have been carried out in a large variety of particle interactions and over a wide range of energies . Recently, pion BE correlations have been investigated in connection with the $`\mathrm{W}`$ mass measurement in $`\mathrm{e}^+\mathrm{e}^{}\mathrm{W}^+\mathrm{W}^{}`$ reactions at LEP at centre-of-mass energies above 161 GeV .
Compared to the abundant information now available on BE correlations in pion pairs, knowledge of the correlations in identical strange boson pairs is scarce and, until recently, was mainly limited to the $`\mathrm{K}_s^0\mathrm{K}_s^0`$ system. $`\mathrm{K}_s^0\mathrm{K}_s^0`$ pairs may exhibit a BE correlation enhancement near threshold even if the origin is not a $`\mathrm{K}^0\mathrm{K}^0`$ or $`\overline{\mathrm{K}^0}\overline{\mathrm{K}^0}`$ system but a $`\mathrm{K}^0\overline{\mathrm{K}^0}`$ boson-antiboson system. A $`\mathrm{K}_s^0\mathrm{K}_s^0`$ low-mass enhancement has recently been observed in hadronic $`\mathrm{Z}^0`$ decays . However, the interpretation of this enhancement as a pure BE correlation effect is complicated by the possible presence of the f$`{}_{0}{}^{}(980)`$ decay into $`\mathrm{K}^0\overline{\mathrm{K}^0}`$. Recently it has been pointed out that the information from BE correlation studies of the $`\mathrm{K}^\pm \mathrm{K}^\pm `$ system, which cannot result from the f$`{}_{0}{}^{}(980)`$ decay, can serve as an effective tool in setting a limit on the resonant fraction of the $`\mathrm{K}_s^0\mathrm{K}_s^0`$ BE enhancement .
The study of $`\mathrm{K}^\pm \mathrm{K}^\pm `$ BE correlations also has a bearing on the relation between the dimension of the emitting source and the mass of the emitted particles; from recent measurements it has been pointed out that the source dimension seems to decrease as the mass increases . Several models have been proposed to account for this behaviour .
This paper presents a study of $`\mathrm{K}^\pm \mathrm{K}^\pm `$ BE correlations using the high statistics sample of $`\mathrm{Z}^0`$ hadronic events recorded by the OPAL detector at LEP. The paper is organised as follows. In Section 2 the methodology used for measuring the BE correlations is presented. The event and track selections are described in Section 3. In Section 4 the analysis of the data is presented, and in Section 5 the systematic effects are studied. Finally, the conclusions are drawn in Section 6.
## 2 Analysis method
The BE correlation function for two identical bosons is defined as:
$`C(p_1,p_2)={\displaystyle \frac{\rho (p_1,p_2)}{\rho (p_1)\rho (p_2)}},`$ (1)
where $`\rho (p_1,p_2)=(1/\sigma )(\mathrm{d}^8\sigma /\mathrm{d}^4p_1\mathrm{d}^4p_2)`$ is the two-particle phase space density subject to BE symmetry, and $`\rho (p)`$ is the corresponding single particle quantity for a particle with four-momentum $`p`$. The correlation function can be studied as a function of the four-momentum difference of the pair, $`Q`$, where $`Q^2`$ = $`(p_1p_2)^2`$ = $`M^24m_{\mathrm{b}oson}^2`$, and $`M`$ is the invariant mass of the pair of bosons each of mass $`m_{\mathrm{b}oson}`$.
From the study of the correlation in pairs one can determine the geometrical and dynamical properties of the emitting source. For a static sphere of emitters with Gaussian density, the correlation function is parametrised with the Goldhaber function as:
$`C(Q)=1+\lambda e^{(RQ)^2},`$ (2)
where $`\lambda `$, the strength of the correlation, is 0 for a completely coherent source and 1 for a completely incoherent one. The parameter $`R`$ in $`\mathrm{G}eV^1`$ is related to the radius of the source, $`R_0`$, through the relation $`R_0=R\mathrm{}c`$.
The two-particle phase space density $`\rho (p_1,p_2)`$ is obtained from a sample of pairs of identical bosons. In this analysis this sample, called the like-sign sample, is formed by pairs of kaons with the same charge. The denominator of Equation 1 $`\rho (p_1)\rho (p_2)`$ is in practice replaced by a reference distribution $`\rho _0(p_1,p_2)`$, which resembles $`\rho (p_1,p_2)`$ in all aspects except in the BE symmetry. A perfect reference sample should have the following properties: absence of BE correlations, and presence of the same correlations as in the the like-sign sample, arising from energy-momentum and charge conservation, the topology and global properties of the events, and resonance or long-lived particle decays.
The principal difficulty for measurements of BE correlations is in the definition of the reference sample from which $`\rho _0(p_1,p_2)`$ is obtained. When correlations amongst like-sign boson pairs are measured, the obvious reference sample is provided by unlike-sign boson pairs. Unfortunately, the $`Q`$ distribution of the unlike-sign pairs includes prominent peaks due to neutral meson resonances (e.g. $`\varphi \mathrm{K}^+K^{}`$). An alternative reference sample can be derived from Monte Carlo simulations in which BE correlations are not included. This relies on a correct simulation of the physics in the complete absence of any BE effect, and a correct modeling of detector effects. A third type of reference sample can be obtained using the methods of event- or hemisphere-mixing, where particles from different events or hemispheres are combined. The last, least model-dependent method was used in the present work.
Each event was divided into two hemispheres separated by the plane perpendicular to the thrust axis and containing the interaction point. A charged kaon track in one hemisphere was combined with a kaon track of the same charge found in the opposite hemisphere after reflecting the momentum of this second track through the origin. To ensure that the like-sign and reference samples were independent, the two kaon tracks forming a pair in the like-sign sample were required to be in the same hemisphere. If all the events have two back-to-back jets, the reference sample of hemisphere-mixed pairs will be similar to the like-sign pair sample apart from the lack of BE correlations. Monte Carlo studies have shown that the hemisphere-mixing technique only works in symmetric topologies such as back-to-back jets. Therefore, two-jet events were selected by requiring a high value of the thrust. The event- and hemisphere-mixing techniques have already been used in other experiments .
## 3 Event and track selection
A detailed description of the OPAL detector can be found in reference . This analysis is primarily based on information from the central tracking detectors, consisting of a silicon microvertex detector, a vertex drift chamber, a jet chamber and $`z`$-chambers<sup>1</sup><sup>1</sup>1A right handed coordinate system is used, with positive $`z`$ along the $`\mathrm{e}^{}`$ beam direction and $`x`$ pointing towards the centre of the LEP ring. The polar and azimuthal angles are denoted by $`\theta `$ and $`\varphi `$, and the origin is taken to be the centre of the detector., all of which lie within an axial magnetic field of 0.435 T. The jet chamber, which has an outer radius of 185 cm, provides up to 159 measurements of space points per track with a resolution in the $`r`$$`\varphi `$ plane of about 135 $`\mu `$m and a transverse momentum<sup>2</sup><sup>2</sup>2The projection onto the plane perpendicular to the beam axis. resolution of $`\sigma _{p_T}/p_T`$ = $`\sqrt{(0.02)^2+(0.0015p_T)^2}`$, with $`p_T`$ in GeV/$`c`$. Particle identification in the jet chamber is possible using the measurement of the specific energy loss $`\mathrm{d}E/\mathrm{d}x`$ with a resolution of approximately of $`3\%`$ for high-momentum tracks in hadronic decays . Since the identification of charged kaons using $`\mathrm{d}E/\mathrm{d}x`$ is crucial to this analysis, the calibration of the energy loss over the many years of data taking was checked and improved when necessary. Control samples of particles identified by techniques other than the energy loss were used to remove year-to-year variations in the measured $`\mathrm{d}E/\mathrm{d}x`$ and to resolve discrepancies between the measured $`\mathrm{d}E/\mathrm{d}x`$ in the data and the theoretical $`\mathrm{d}E/\mathrm{d}x`$ in the Monte Carlo.
This analysis used a sample of about 3.9 million hadronic $`\mathrm{Z}^0`$ decays recorded at LEP between the years 1992 and 1995. A sample of 6.75 million Monte Carlo hadronic events generated with JETSET 7.4 and tuned to reproduce the global features of the events was also used. The generated events were processed through a detailed simulation of the experiment and subjected to the same event and track selection as the data. A detailed description of the selection of hadronic events is given in . Events with two clear back-to-back jets, necessary for the proper functioning of the hemisphere-mixing technique, were selected by requiring that the thrust value was larger than 0.95. About 30$`\%`$ of the events passed the thrust cut.
Charged tracks were required to have a minimum transverse momentum of 150 MeV/$`c`$, a maximum reconstructed momentum of 60 GeV/$`c`$, a distance of closest approach to the interaction point less than 0.5 cm in the plane orthogonal to the beam direction and the corresponding distance along the beam direction of less than 40 cm. The first measured point had to be within a radius of 75 cm from the interaction vertex. The cuts on the transverse and longitudinal distances of closest approach help to remove particles from long-lived decays. About 50$`\%`$ of the reconstructed tracks passed these selections.
Kaons were identified using the $`\mathrm{d}E/\mathrm{d}x`$ measurement. Only tracks with at least 20 hits available for the measurement of $`\mathrm{d}E/\mathrm{d}x`$ and with a polar angle satisfying $`|\mathrm{cos}\theta |<\mathrm{\hspace{0.17em}0.9}`$ were considered. For each track, a $`\chi ^2`$ probability was formed for each stable particle hypothesis: e, $`\mu `$, $`\pi `$, $`\mathrm{K}`$ and p. A track was identified as a kaon candidate if it had a probability of at least 10$`\%`$ of being a kaon and if the probability of being a kaon was larger than the probability of being any other of the above particle species. In addition, in order to reject pions, each track was required to have a pion probability less than 5$`\%`$. Electrons from photon conversions were rejected using a neural network algorithm as described in . According to the Monte Carlo, approximately 34$`\%`$ of the true kaons passed these requirements and the estimated kaon purity of the track sample was about 72$`\%`$ on average, with variations between 50$`\%`$ and 97$`\%`$ depending on the momentum of the track. The lowest purity corresponded to the momentum range between 0.9 and 1.5 GeV/$`c`$ where the pion and kaon bands overlap in $`\mathrm{d}E/\mathrm{d}x`$ . In the selected kaon track sample the fraction of pions was estimated as 17$`\%`$, that of protons as 11$`\%`$, and the contribution of muons and electrons was negligible.
## 4 Data analysis
The $`Q`$ distributions of the like-sign kaon candidate pairs and the hemisphere-mixed kaon candidate pairs were determined using tracks that passed the selections described in Section 3. In the data, 76063 like-sign and 98558 hemisphere-mixed kaon candidate pairs with $`Q<\mathrm{\hspace{0.17em}2.0}\text{GeV}`$ were selected. In the Monte Carlo, the corresponding numbers of kaon candidate pairs were 109601 and 136785 respectively.
Since the BE correlation manifests itself only for identical particles, it was necessary to correct for impurities. The Monte Carlo predicted that the main contamination in the sample of like-sign kaon candidate pairs was from $`\mathrm{K}\pi `$ and $`\mathrm{K}p`$ pairs. In this sample, and for pairs with $`Q<\mathrm{\hspace{0.17em}2.0}\text{GeV}`$, the estimated fraction of $`\mathrm{KK}`$ pairs was about 48$`\%`$, the $`\mathrm{K}\pi `$ fraction about 27$`\%`$ and the $`\mathrm{K}p`$ fraction about 13$`\%`$. The contamination from pairs susceptible to BE correlations or Fermi-Dirac anti-correlations (i.e. $`\pi \pi `$ and $`\mathrm{pp}`$) was negligible, of the order of 3$`\%`$ for $`\pi \pi `$ and less than 1$`\%`$ for $`\mathrm{pp}`$ pairs. The fraction of $`\mathrm{KK}`$ pairs was constant over the whole $`Q`$ range. The sample of hemisphere-mixed kaon candidate pairs had an almost identical composition.
The Monte Carlo non-$`\mathrm{KK}`$ $`Q`$ distribution was subtracted from the data distribution using the fraction given by the simulation and normalised to the number of data pairs. Both like-sign and hemisphere-mixed pair distributions were corrected using this technique. Figure 1 shows these corrected $`Q`$ distributions for the data and the Monte Carlo in which BE correlations were not simulated. The hemisphere-mixed $`Q`$ distributions were normalised to the like-sign $`Q`$ distributions in the region $`0.6<Q<\mathrm{\hspace{0.17em}2.0}\text{GeV}`$ where no BE correlations are expected. Both data and Monte Carlo distributions show a similar behaviour at high values of $`Q`$: the hemisphere-mixed distribution is above the like-sign one in the region $`0.7<Q<\mathrm{\hspace{0.17em}1.2}\text{GeV}`$; there is a cross-over of both distributions at $`Q`$ about 1.2 GeV; and the like-sign distribution is above the hemisphere-mixed one at $`Q>1.2\text{GeV}`$.
Figure 2 shows the ratio $`N_{++}(Q)/N_{mix}(Q)`$ in both the data and the Monte Carlo, where $`N_{++}(Q)`$ is the number of like-sign pairs and $`N_{mix}(Q)`$ is the number of hemisphere-mixed pairs as functions of $`Q`$. The first two bins of the distribution shown in the figure were combined due to the small statistics. The data distribution shows a clear enhancement in the region $`Q<`$ 0.3 GeV. There is also a rise of the correlation at high values of $`Q`$ normally attributed to long-range correlations. Indeed, the slope of the correlations at high $`Q`$ is well reproduced by the Monte Carlo. The ratio $`N_{++}(Q)/N_{mix}(Q)`$ in the Monte Carlo deviates slightly from unity at high $`Q`$ and falls slowly with decreasing $`Q`$, probably due to features of string fragmentation and local conservation of charge and strangeness. These effects can be taken into account if both like-sign and hemisphere-mixed data distributions are divided by the corresponding Monte Carlo distributions.
The correlation function was therefore defined as the double-ratio:
$`C_{mix}(Q)={\displaystyle \frac{N_{++}^{data}(Q)}{N_{mix}^{data}(Q)}}/{\displaystyle \frac{N_{++}^{MC}(Q)}{N_{mix}^{MC}(Q)}},`$ (3)
and was parametrised using a modified version of Equation 2:
$`C(Q)=N(1+\lambda e^{(RQ)^2})(1+\delta Q+ϵQ^2),`$ (4)
where $`N`$ is a normalisation factor, and the empirical term $`(1+\delta Q+ϵQ^2)`$ accounts for the behaviour of the correlation function at high $`Q`$ values due to any remaining long-range correlations. Figure 3 shows the correlation $`C_{mix}(Q)`$ with the result of the fit. The fitted parameters and the correlation coefficients between them are given in Table 1. The fit has a $`\chi ^2`$ = 43 for 35 degrees of freedom.
## 5 Systematic effects
Systematic effects arising from the event and track selections, the modeling of $`\mathrm{d}E/\mathrm{d}x`$ in the Monte Carlo, the parametrisation of the correlation function and the choice of the reference sample are considered. In each case, the result of the fit to the correlation function is given in Table 2, with row (a) giving the result of the basic fit discussed in the previous section.
The overall systematic uncertainties in the parameters $`\lambda `$ and $`R_0`$ were calculated as the largest single deviations between the parameters of the fits from rows (b) to (l), and the parameters of the basic fit in row (a). The final values of the parameters are
$`\lambda `$ $`=`$ $`\text{0.82}\text{ }\pm \text{ }\text{0.22}_{0.12}^{+0.17}`$
$`R_0`$ $`=`$ $`\text{0.56}\text{ }\pm \text{ }\text{0.08}_{0.06}^{+0.08}\mathrm{f}m.`$
The mismodeling of the kaon momentum spectrum, residual BE correlations in the reference sample, the origins of the kaons and final-state interaction corrections are also discussed in this section.
### 5.1 Event and track selection
* As discussed in Section 2, the hemisphere-mixing technique only works when two-jet events are selected. The analysis was repeated for events selected using a cone jet finding algorithm instead of the cut in thrust (row (b)). The value of the thrust cut was also changed from 0.95 to 0.93 and the analysis repeated (row (c)).
* To obtain a purer sample of kaons, tracks with momenta in the pion-kaon $`\mathrm{d}E/\mathrm{d}x`$ overlap region, $`0.9<p<\mathrm{\hspace{0.17em}1.5}`$ GeV/$`c`$, were rejected (row (d)).
* The minimum number of $`\mathrm{d}E/\mathrm{d}x`$ hits required for each track was increased from 20 to 40 (row (e)).
### 5.2 Parametrisation of $`\mathrm{d}E/\mathrm{d}x`$
Since the analysis relies to a large extent on the separation of pions from kaons, it is especially important to understand the $`\mathrm{d}E/\mathrm{d}x`$ measurements of the copiously-produced pions. A sample of pions was identified in $`\mathrm{K}_s^0\pi ^+\pi ^{}`$ decays and was used to estimate the mismodeling of the normalised ionisation energy loss $`N_{\mathrm{d}E/\mathrm{d}x}^\sigma `$ . The normalised $`\mathrm{d}E/\mathrm{d}x`$ is defined as $`N_{\mathrm{d}E/\mathrm{d}x}^\sigma =(\mathrm{d}E/\mathrm{d}x(\mathrm{d}E/\mathrm{d}x)_0)/(\sigma (\mathrm{d}E/\mathrm{d}x))_0`$. Here, $`\mathrm{d}E/\mathrm{d}x`$ is the measured value, while $`(\mathrm{d}E/\mathrm{d}x)_0`$ and $`(\sigma (\mathrm{d}E/\mathrm{d}x))_0`$ are the expected value and the expected error assuming the track to be a pion. This analysis showed that the mean value and the width of the normalised $`\mathrm{d}E/\mathrm{d}x`$ distribution were both known to within $`\pm 10\%`$ of $`\sigma `$. By changing in the Monte Carlo the normalised $`\mathrm{d}E/\mathrm{d}x`$ of both kaons and pions by their known uncertainties, the fraction of non-$`\mathrm{KK}`$ pairs was found to vary by up to $`\pm 10\%`$. The BE analysis was consequently repeated with the assumed impurity value set to 57$`\%`$ and 47$`\%`$ (rows (f) and (f)).
### 5.3 Fit of the correlation function
* The binning of the $`Q`$ distributions was modified from 50 MeV to 20 MeV (row (g)).
* The $`Q`$ distribution normalisation range was changed from $`0.6<Q<\mathrm{\hspace{0.17em}2.0}\text{GeV}`$ to $`0.8<Q<\mathrm{\hspace{0.17em}2.0}\text{GeV}`$ (row (h)).
* The fit was repeated with the first bin $`Q<\mathrm{\hspace{0.17em}0.05}\text{GeV}`$, excluded. This was done to test the effect of potential problems at low $`Q`$ because of the limited resolution (row (i)).
* If the use of the double-ratio removes all correlations other than BE, such as long-range correlations at high $`Q`$, then the empirical term $`(1+\delta Q+ϵQ^2)`$ of Equation 4 would not be necessary. $`C_{mix}(Q)`$ was parametrised using the simplified function:
$`C(Q)=N(1+\lambda e^{(RQ)^2}).`$ (5)
The parameters of the fit are given in row (j). The fit has $`\chi ^2`$ = 49 for 37 degrees of freedom . This fit was also repeated with the range limited to $`Q<1.5`$ GeV to reduce possible long-range correlations (row (k)).
### 5.4 The reference sample
The BE correlation was also measured using Monte Carlo like-sign pairs as the reference sample in the sub-sample of two-jet events (row (l)). The correlation function in this case was defined as:
$`C_{MC}(Q)={\displaystyle \frac{N_{++}^{data}(Q)}{N_{++}^{MC}(Q)}}.`$ (6)
Figure 4 shows the correlation function $`C_{MC}(Q)`$ together with the results of a fit using Equation 4. The values of the parameters are: $`N`$ = 0.88 $`\pm `$ 0.05; $`\lambda `$ = 0.92 $`\pm `$ 0.17; $`R_0`$ = 0.59 $`\pm `$ 0.06 fm; $`\delta `$ = 0.04 $`\pm `$ 0.08 $`\mathrm{G}eV^1`$; $`ϵ`$ = 0.05 $`\pm `$ 0.03 $`\mathrm{G}eV^2`$, where the errors are statistical only. The fit has $`\chi ^2`$ = 42 for 35 degrees of freedom.
Although there are known imperfections in the simulation, particularly at low momenta in the kaon momentum spectrum (see section 5.5), these results were taken as an indication of systematic differences between the hemisphere-mixing method and the simple use of a Monte Carlo reference sample.
### 5.5 Mismodeling of the kaon momentum spectrum
If the simulation were perfect, one would expect to get the same results by measuring the correlation function with $`C_{mix}(Q)`$ and with $`C_{MC}(Q)`$. The double-ratio
$`C_{mix}(Q)={\displaystyle \frac{N_{++}^{data}(Q)}{N_{mix}^{data}(Q)}}/{\displaystyle \frac{N_{++}^{MC}(Q)}{N_{mix}^{MC}(Q)}}{\displaystyle \frac{N_{++}^{data}(Q)}{N_{++}^{MC}(Q)}}/{\displaystyle \frac{N_{mix}^{data}(Q)}{N_{mix}^{MC}(Q)}}`$ (7)
is equivalent to $`C_{MC}(Q)`$ only if the hemisphere-mixed sample is perfectly modelled, which implies $`N_{mix}^{data}(Q)/N_{mix}^{MC}(Q)=1`$. Figure 5 shows this ratio; the general agreement is good, although at $`Q<\mathrm{\hspace{0.17em}0.8}\text{GeV}`$ there is some indication that the ratio is below unity.
Any important mismodeling of the Monte Carlo would indicate that the results obtained with $`C_{MC}(Q)`$ and with $`C_{mix}(Q)`$ (since the Monte Carlo is used for normalising this correlation function) were not reliable. Therefore, the distribution of Monte Carlo tracks in a sample free from BE correlations was compared to the same distribution in the data. Figure 6 shows the momentum spectrum of kaon candidate tracks of the hemisphere-mixed sample in the data and in the Monte Carlo normalised to the same total number of pairs. At low momentum, the Monte Carlo does not describe the data spectrum well and differences are seen at the $`\pm 15\%`$ level. Studies of the differential cross-section of kaons in hadronic events in both the data and the Monte Carlo events, generated with JETSET 7.4 and tuned according to OPAL data, showed that the simulation predicted a kaon spectrum which is too soft.
As a check of the stability of the results obtained with $`C_{mix}(Q)`$, the $`Q`$ spectra of both like-sign and hemisphere-mixed pairs were reweighted. Each pair of kaons in the Monte Carlo was reweighted by the product of the weights of each kaon in the pair, where the weight was obtained in bins of momentum by dividing the data momentum spectrum by the Monte Carlo spectrum. The final measurement of the correlation function did not change significantly after reweighting, $`\lambda `$ varied by $`+`$0.03 and $`R_0`$ by $``$0.01 fm. The same exercise was done to check the results obtained with $`C_{MC}(Q)`$. In this case, $`\lambda `$ varied more significantly, by $`+`$0.12, and $`R_0`$ by $``$0.01 fm. Thus, by the use of a double-ratio technique, the correlation function $`C_{mix}(Q)`$ was found to be less sensitive to the Monte Carlo mismodeling than the correlation function $`C_{MC}(Q)`$.
### 5.6 Residual BE correlations in the reference sample
As suggested in , residual BE correlations could be a source of imperfection in the reference sample. To check that the hemisphere-mixed sample was free of effects induced by the BE correlations, a Monte Carlo study was done using the generator JETSET 7.4. Hadronic events were generated with and without BE correlations, with the BE correlations simulated assuming the Goldhaber parametrisation described in Section 2. The shape of the $`Q`$ distribution of the hemisphere-mixed pairs remained unchanged after including the BE correlations in the generation. The ratio $`N_{mix}^{withBE}(Q)/N_{mix}^{withoutBE}(Q)`$ was consistent with unity in the full range of $`Q`$, demonstrating that the reference sample was free from effects due to residual BE correlations.
### 5.7 Sources of kaons
In a substantial fraction of the kaon pairs, one or more of the kaons result from a long-lived particle decay — in such cases the kaon pairs cannot exhibit BE correlations. It is therefore useful to separate the various sources of kaons and to classify the parent particles to estimate the maximum possible value of $`\lambda `$, as suggested in . On the other hand, some studies have suggested that the correlation function is narrowed by the contribution of decay products of long-lived sources, and also that resonance decays induce a pair-momentum dependence of the radius.
The kaon sources as predicted by the JETSET Monte Carlo simulation are given in Table 3. These have been classified as in into two main groups: long-lived sources with life-time $`c\tau >`$ 10 fm, and short-lived sources with life-time $`c\tau <`$ 10 fm. The table shows that the fraction of kaon pairs at low $`Q`$ ($`<0.6`$ GeV) with at least one kaon from a short-lived source is about 81$`\%`$, so that the fraction of pairs in which both kaons arise from short-lived sources is about 66$`\%`$.
The pairs from short-lived sources cannot be identified in the data, and so the final results of this analysis were not corrected for the effect of short-lived sources because such a correction would be based on a Monte Carlo model with its inherent uncertainties. However, to illustrate the magnitude of the effect, a correction was applied to the correlation function using the estimated fraction of short-lived sources (66$`\%`$), assuming that kaons from long-lived sources do not contribute to the correlations. The fitted parameters of the corrected correlation function are: $`\lambda `$ = 1.27 $`\pm `$ 0.31 and $`R_0`$ = 0.55 $`\pm `$ 0.07 fm, where the errors are statistical only.
### 5.8 Final-state interactions
Charged kaons are subject to both the Coulomb and the strong interactions. In principle, every pair of like-sign kaons from short-lived sources should be corrected for these interactions in the data (but not in the Monte Carlo, where they were not simulated). To apply a correction to all pairs would result in an overestimate of the value of the strength parameter . As in Section 5.7, the final results of this analysis were not corrected for the electromagnetic and strong interactions because of the model-dependence of such corrections.
As a check of the possible magnitude of any correction, the electromagnetic repulsion of like-sign pairs was corrected for in the data. The like-sign kaon pair $`Q`$ spectrum was corrected using the Gamow factor $`G(\eta )=2\pi \eta /(e^{2\pi \eta }1)`$, where $`\eta =\alpha _{em}m_\mathrm{K}/Q`$, $`\alpha _{em}`$ is the electromagnetic coupling constant and $`m_\mathrm{K}`$ the kaon mass. On the assumption that all pairs are from short-lived sources, the fitted strength and radius of the Coulomb corrected correlation function are: $`\lambda `$ = 0.92 $`\pm `$ 0.25 and $`R_0`$ = 0.61 $`\pm `$ 0.17 fm, where the errors are statistical only. The correlation function was also corrected for both the effect of short-lived sources as in Section 5.7 and the Coulomb effect, resulting in the fitted parameters: $`\lambda `$ = 1.36 $`\pm `$ 0.55 and $`R_0`$ = 0.58 $`\pm `$ 0.11 fm, where the errors are statistical only.
## 6 Discussion and conclusions
Bose-Einstein correlations have been measured in identified pairs of charged kaons in hadronic $`\mathrm{Z}^0`$ decays using the OPAL experiment at LEP. The analysis was performed in events with a clear two-jet topology, a requirement which was necessary to obtain a suitable reference sample using a hemisphere-mixing technique. Monte Carlo simulation was used to correct for imperfections in the reference sample by use of a double-ratio for the correlation function. The enhancement was parametrised using a Gaussian formula, resulting in a strength
$`\lambda `$ $`=`$ $`\text{0.82}\text{ }\pm \text{ }\text{0.22}_{0.12}^{+0.17}`$
and a kaon emitter radius
$`R_0`$ $`=`$ $`\text{0.56}\text{ }\pm \text{ }\text{0.08}_{0.06}^{+0.08}\mathrm{f}m.`$
A definite conclusion from the present analysis is a confirmation of the results of reference , that there are indeed BE correlations in $`\mathrm{K}^\pm \mathrm{K}^\pm `$ pairs from hadronic $`\mathrm{Z}^0`$ decays. This implies that there should be such correlations in $`\mathrm{K}_s^0\mathrm{K}_s^0`$ pairs ; therefore it is unlikely that the previously observed threshold enhancements can be attributed entirely to the f$`{}_{0}{}^{}(980)`$ decay into kaons.
Values of $`\lambda `$ and $`R_0`$, as measured in hadronic $`\mathrm{Z}^0`$ decays at LEP for various particle types, are listed for comparison in Table 4. Since there is evidence that $`\lambda `$ and $`R_0`$ may depend on the event topology, the table gives the type of event used in the various measurements. The reference sample types used in each analysis are also listed: these may be event- or hemisphere-mixed, unlike-sign or Monte Carlo pairs.
In all events and for correlations measured using the unlike-sign reference sample, the radius of charged pion emitters varies between 0.8 and 1.0 fm while the radius of the charged kaon emitters is $`0.48\pm 0.08`$ fm. This gives the relation $`R_0(\pi ^\pm \pi ^\pm )>R_0(\mathrm{K}^\pm \mathrm{K}^\pm )`$ — a mass dependence of the emitting source, as already pointed out in . In two-jet events, the measured radius of charged pion emitters is seen to vary between 0.4 and 0.9 fm, inconsistent within the quoted errors. This large variation is usually attributed to the choice of the reference sample: in the case of the unlike-sign reference sample the radius is about 0.8–0.9 fm; in the case of the mixed reference sample the radius is about 0.4–0.5 fm. The comparison of results obtained using the event- or hemisphere-mixing techniques and a double-ratio for the correlation function shows that the present measurement of the radius of the kaon source, $`0.56\pm 0.11`$ fm, is compatible with the previous measurements of the radius of pion sources and does not support a strong mass dependence of the emitting source. However, both measurements of $`R_0`$ for kaon pairs are considerably larger than that obtained for $`\mathrm{\Lambda }\mathrm{\Lambda }`$ pairs, 0.14$`{}_{0.03}{}^{}{}_{}{}^{+0.07}`$ fm .
Acknowledgements:
We particularly wish to thank the SL Division for the efficient operation of the LEP accelerator at all energies and for their continuing close cooperation with our experimental group. We thank our colleagues from CEA, DAPNIA/SPP, CE-Saclay for their efforts over the years on the time-of-flight and trigger systems which we continue to use. In addition to the support staff at our own institutions we are pleased to acknowledge the
Department of Energy, USA,
National Science Foundation, USA,
Particle Physics and Astronomy Research Council, UK,
Natural Sciences and Engineering Research Council, Canada,
Israel Science Foundation, administered by the Israel Academy of Science and Humanities,
Minerva Gesellschaft,
Benoziyo Center for High Energy Physics,
Japanese Ministry of Education, Science and Culture (the Monbusho) and a grant under the Monbusho International Science Research Program,
Japanese Society for the Promotion of Science (JSPS),
German Israeli Bi-national Science Foundation (GIF),
Bundesministerium für Bildung, Wissenschaft, Forschung und Technologie, Germany,
National Research Council of Canada,
Research Corporation, USA,
Hungarian Foundation for Scientific Research, OTKA T-029328, T023793 and OTKA F-023259. |
warning/0001/cond-mat0001124.html | ar5iv | text | # A scalable solid-state quantum computer based on quantum dot pillar structures
## I Introduction
The possibility that a computer with exceptional properties could be built employing the laws of quantum physics has stimulated considerable interest in searching for useful algorithms and realizable physical implementations. Two useful algorithms, exhaustive search and factorization, have been discovered; others, including the suggestion that quantum computers will prove useful to model quantum systems, are being sought. Meanwhile, various physical implementations are being explored, including trapped ions, cavity quantum electrodynamics, ensemble nuclear magnetic resonance, small Josephson junctions, optical devices incorporating beam splitters and phase shifters and a number of solid state systems based on quantum dots. Although the advantages of quantum computing are enormous for particular key applications, the requirements for their implementation are extremely stringent, perhaps especially rigorous for solid-state systems. Nevertheless solid-state quantum computers are very appealing relative to other possible implementation schemes because of the well-known ability to customize the design through the use of artificially structured materials and the probable scalability of the resulting design. For example, integrated circuit manufacturing technology would be immediately applicable to quantum computers of the proper implementation; and such designs would not only be scalable to smaller dimensions along the ”semiconductor learning curve” but also large ensembles of ”identical” quantum computers could be manufactured, that could be individually fine-tuned electrically. To date, no solid-state implementation of quantum computing has been demonstrated.
In this paper, we investigate a solid-state quantum computer implementation that is amenable to manufacturing with integrated circuit technology. We develop a three-dimensional (3D) device model and self-consistently solve coupled Schrödinger and Poisson equations to generate a quantum computer design for a three-qubit quantum register that is based on pair wise coupled asymmetric III-V quantum dots. The design is optimized for a long coherence time and a rapid computation rate. Our results indicate that this structure may provide a realistic scalable candidate for quantum computing in solid-state systems.
## II Proposed Structure
The proposed quantum dot quantum computer (see Fig. 1) consists of a pillar structure composed of a chain of asymmetric quantum dots separated by intervening layers of higher bandgap composition fabricated in a GaAs/AlGaAs technology by means of a sequence of planar MBE growth steps and subsequent etching to form the pillar. A sheath of similar AlGaAs composition is then grown surrounding the pillar and a wrap-around gate electrode deposited. A drain (source) is formed at the top (bottom), the series of asymmetric quantum dots are in the center region, and the gate surrounds the region of the pillar containing the quantum dots. Tarucha et al. have reported similar n-type single electron transistor (SET) structures. Electron confinement along the pillar axis is produced by the band gap discontinuity of the dot structure. Encasing the quantum dot structure in the pillar core by the cylindrical sheath and the gate electrode provides confinement in the radial direction. By applying a negative bias that depletes carriers near the surface, an additional parabolic electrostatic potential is formed that allows for tuning of the radial confinement and localization of one electron per dot. The simultaneous insertion of a single electron per dot is accomplished by lining up the quantum dot ground state levels so that they lie close to the Fermi level; a single electron is confined in each dot over a finite range of the gate voltage due to shell filling effects.
Thus, the pillar consists of a vertical stack of coupled asymmetric GaAs/AlGaAs quantum dots of differing size and composition so that each dot possesses a distinct energy structure. Qubit registers, $`|0`$ and $`|1`$, are based on the ground and first excited state of the single electron within each quantum dot. Overall, parameters of the structure can be chosen to produce a well-resolved spectrum of distinguishable qubits. The asymmetric dots produce large built-in electrostatic dipole moments between the ground and first excited state, and electrons in adjacent dots are coupled through the electric dipole-dipole interaction, while coupling between non-adjacent dots is significantly weaker. This produces the desired quantum computer consisting of a linear array of binary states (qubits) with pair wise pillar-axis coupling between adjacent qubits. In addition to energy tuning, the asymmetry of each quantum dot can be designed so that dephasing due to electron-phonon scattering and spontaneous emission is minimized. The combination of strong dipole-dipole coupling and long dephasing times make it possible to perform many computational steps before loss of coherence, in fact, it is believed possible to design this device so that error correction substantially prohibits coherence loss.
Quantum computations are performed by means of a series of coherent optical pulses in the far infrared. Final readout of the amplitude and phase of the qubit states can be achieved through quantum state holography. Amplitude and phase information are extracted through mixing the final state with a reference state generated in the same system by an additional delayed laser pulse and detecting the total time- and frequency-integrated fluorescence as a function of the delay. Extracting the final state information using quantum state holography requires multiple experiments, one for each delay, as described in Ref. . Thus, the computation must be performed several times before an answer is arrived at. This is no real problem since the number of repetitions needed is only on the order of 40 or so, independent of the number of computational steps in a given quantum algorithm. Through the use of integrated circuit manufacturing technology, it is possible to simultaneously fabricate a large array of ”identical” pillar quantum dot quantum computers, that is, on the order of $`10^{10}`$ per wafer. Each of these quantum registers could be electrically connected through deposited interconnect in such a manner so that each could be individually tunable to produce an array of identical units. In general, inhomogeneity among the quantum dots will result in slightly different energy levels. Sherwin et al. have recently pointed out that one can perform accurate qubit operations in an inhomogeneous population of quantum dots arising from quenched disorder due to static charged defects, for example, provided that each SET is independently calibrated. This calibration can done by performing simple gate operations and tuning the gate electrodes appropriately. Efficient optical coupling to the resulting ensemble can be achieved through optical light guiding as suggested in Ref. . By this means direct observation of fluorescence is possible. Quantum computations are performed by means of a series of coherent optical pulses in the far infrared, and may be carried out in complete analogy with the operation of an NMR quantum computer. It should be remarked that while our scheme resembles NMR ensemble quantum computation in the use of a series of optical pulses to perform quantum logic gates, it differs from NMR quantum computation in that our use of a collection of single electron transistors is done to enable a stronger signal to noise ratio in the readout phase. In principle, the quantum computation could be done with only a single SET transistor structure if the readout measurements were sufficiently sensitive.
## III Device Model
In the context of studies of the Coulomb blockade in self-organized quantum dots and planar single-electron transistors, self-consistent calculations of electronic structure, shell filling effects, electron-electron interaction, Coulomb degeneracy, and Coulomb oscillation amplitudes have been carried out for various quantum dot structures. Our quantum register can be analyzed using methods similar to those used to study the self-consistent electronic structure in single-electron transistors. The problem we address is similar to those addressed by other authors who are interested in obtaining current-voltage characteristics and studying Coulomb oscillations in single-electron transistors over a wide range of gate biasing and shell filling conditions.
In our case, we are interested in obtaining the self-consistent electrostatic potential and electronic eigenstates in an equilibrium configuration in which the source and drain are grounded and the gate electrode is negatively biased. The electrostatic potential, $`V(r)`$, is obtained by solving the Poisson equation for n-doped semiconductors
$$^2V(\stackrel{}{r})=\frac{4\pi }{\epsilon }q\left[n(\stackrel{}{r})+N_D^+(\stackrel{}{r})\right]$$
(1)
In the Poisson equation, $`q`$ is the absolute value of the electron charge, $`\epsilon `$ is the static dielectric constant, $`n(\stackrel{}{r})`$ is the electron concentration and $`N_D^+(\stackrel{}{r})`$ is the known concentration of ionized donors in the structure. For the dielectric constant, we adopt the GaAs value $`\epsilon =12`$. The Poisson equation is solved subject to boundary conditions on the electrostatic potential, $`V(\stackrel{}{r})`$. At the interfaces between the semiconductor and the source, drain and gate electrodes, $`V(\stackrel{}{r})`$ is equal to the applied gate voltage while at the semiconductor-vacuum interfaces, the normal derivative of $`V(\stackrel{}{r})`$ vanishes.
Following Ref. , the global electron concentration, $`n(\stackrel{}{r})`$, in the device is obtained by partitioning the pillar structure into ”bulk” and ”quantum” regions. In the ”bulk” regions far from the quantum wells i.e. the source and drain regions, electrons are treated in the Thomas-Fermi approximation and the electron concentration is given by
$$n(\stackrel{}{r})=\{\begin{array}{cc}\frac{1}{3\pi ^2}\left[\frac{2m_e^{}}{\mathrm{}^2}(\mu U(\stackrel{}{r}))\right]^{3/2}\hfill & \text{ if }U(\stackrel{}{r})<\mu \hfill \\ 0\hfill & \text{ otherwise}\hfill \end{array}$$
(2)
where $`\mu `$ is the chemical potential and $`U(\stackrel{}{r})`$ is the effective electron potential. The chemical potential, $`\mu `$, is determined through the requirement that overall charge neutrality be maintained in the bulk regions, i.e. the chemical potential is adjusted until
$$\left(n(\stackrel{}{r})N_D^+(\stackrel{}{r})\right)𝑑\stackrel{}{r}=0$$
(3)
where the integration is carried out over the bulk source and drain regions.
The effective potential, $`U(\stackrel{}{r})`$, in the bulk regions includes the Hartree potential, $`U_H=qV(\stackrel{}{r})`$, and the conduction band offset, $`\mathrm{\Delta }E_c`$, which depends on the local Al concentration, $`x`$. Thus,
$$U(\stackrel{}{r})=qV(\stackrel{}{r})+\mathrm{\Delta }E_c$$
(4)
where the conduction band offset, $`\mathrm{\Delta }E_c`$, is taken to be 60% of the difference between the $`Al_xGa_{1x}As`$ and $`GaAs`$ bandgaps. Using the bandgap variation of $`Al_xGa_{1x}As`$ determined by Lee et al., we obtain the following expression for the conduction band offset as a function of the local Al concentration, $`x`$:
$$\mathrm{\Delta }E_c=0.6(1155x+370x^2)\text{meV}$$
(5)
In the ”quantum” regions containing the quantum wells, the electron concentration, $`n(r)`$, is determined by the electron wavefunctions, $`\psi _i(r)`$, and energies, $`E_i`$, through the relation
$$n(\stackrel{}{r})=\underset{i}{}n_i|\psi _i(\stackrel{}{r})|^2$$
(6)
The electron occupancy in each level, $`n_i`$, is a function of the electron energy and the temperature.
The electron wavefunctions and energy levels, $`E_i`$, are obtained by solving the Schrödinger equation in the effective mass approximation
$$\left[\frac{\mathrm{}^2}{2m_e^{}}^2+U(\stackrel{}{r})E_i\right]\psi _i(\stackrel{}{r})=0$$
(7)
The electron potential, $`U(\stackrel{}{r})`$, in the quantum regions is given by
$$U(\stackrel{}{r})=qV(\stackrel{}{r})+\mathrm{\Delta }E_c+U_{xc}(\stackrel{}{r})$$
(8)
where $`U_{xc}(\stackrel{}{r})`$ is the exchange-correlation potential of Perdew and Zunger.
In the quantum register discussed in the next section, the gate voltage is negatively biased in such a way that a single electron is strongly localized in each electrostatically confined quantum dot. The radial confinement potential is strong enough that the lowest few electron wavefunctions are strongly localized near the center of the pillar and die away far from the semiconductor-electrode interface. In our design, the quantum wells are wide enough and the barriers between the quantum wells are thick enough so that the lowest few electron wavefunctions do not penetrate to the center of the barriers separating the quantum wells. Since all the wavefunctions of interest vanish at the center of these barriers, we can divide the quantum well region into several regions (one for each qubit). These regions are taken to be cylinders stacked along the pillar axis with top and bottom surfaces located at the centers of the barriers between adjacent wells. We solve the Schrödinger equation in each dot separtely subject to the boundary condition that all wavefunctions vanish at the region boundaries.
Due to the cylindrical symmetry of the structure, the 3D Schrödinger equation can be reduced to a 2D equation in cylindrical coordinates. One might try to solve the 2D Schrödinger equation by finite differencing the partial differential equation and solving the resulting matrix eigenvalue equation. The size of the matrix to be diagonalized is equal to the number of interior mesh points in the 2D grid and this is much too large to be handled easily. Other authors have taken this brute-force approach to solving the Schrödinger equation in self-consistent Poisson-Schrödinger problems with the result that solving the Schrödinger equation is the most time consuming part of the computation. We find that it is possible to do better. In solving the 2D Schrödinger equation, we first approximate $`U(\rho ,z)`$ in each quantum dot by a separable potential
$$U(\rho ,z)U_s(\rho ,z)U_r(\rho )+U_z(z)$$
(9)
where the axial potential is defined as
$$U_z(z)=\frac{2}{R^2}_0^RU(\rho ,z)\rho 𝑑\rho $$
(10)
and the radial potential is given by
$$U_r(\rho )=\frac{1}{}_0^{}\left(U(\rho ,z)U_z(z)\right)𝑑z$$
(11)
In these last two expressions, $`R`$ and $``$ are the radius and height of the cylindrical region over which $`U(\rho ,z)`$ is defined in each dot. With the separable potential approximation, the 2D Schrödinger equation can be separated into two 1D equations which can be cast as finite difference eigenvalue equations and solved numerically for the electron energies and wavefunctions. The resulting 2D wavefunctions are the best product wavefunctions that approximate the solution of the 2D Schrödinger equation in each qubit. The electronic states in the separable potential approximation in our cylindrical pillar are labeled by three quantum numbers $`(n_\rho `$, $`n_\varphi `$, $`n_z)`$ which specify the number of nodes in the product wavefunctions associated with cylindrical coordinates $`\rho `$, $`\varphi `$, and $`z`$. In this notation, the qubit state $`|0`$ is denoted $`(0,0,0)`$ while $`|1`$ is denoted $`(0,0,1)`$. We find that the separable wavefunctions are reasonable approximations to the true wavefunctions since we are starting with a separable potential which is already close to the true potential in some average sense. We next obtain the exact energies and wavefunctions of the original non-separable Schrödinger equation by treating the residual $`U(\rho ,z)U_s(\rho ,z)`$ as a perturbation and expanding the exact wavefunctions as a sum of separable wavefunctions. Our expansion of the true wavefunctions in terms of separable wavefunctions is rapidly converging and we find that the dominant terms in the expansion of the true wavefunctions are the separable wavefunctions of the same symmetry. Our approach to solving the 2D Schrödinger equation is fast and most of computing time is spent solving the Poisson equation.
To complete the specification of the electron charge density in the quantum dots, it is necessary to compute the electron occupation numbers, $`n_i`$. One might expect that $`n_i`$ would be given by the Fermi-Dirac distribution and indeed this would be the case if the electrons in the dots were delocalized and in tunneling contact with the leads. In this case, the qubits could exchange electrons with their environment and the total number of electrons in the dot $`N=_in_i`$ could take on non-integer values. But clearly this is not tolerable in a quantum computer and we must carefully arrange things so the dot wavefunctions exhibit a high degree of localization. In this situation, only an integer number of electrons can occupy the dot and this constraint gives rise to what is known as the Gibbs distribution. The number of electrons, $`N`$, in the dot is determined by minimizing the Gibbs free energy with respect to the integer number of electrons, $`N`$. The Gibbs free energy is $`F(N)=kT\mathrm{ln}[Z(N)]`$, where the grand canonical partition function, $`Z(N)`$, is given by
$$Z(N)=\underset{\{n_i\}}{}\mathrm{exp}\left[\frac{_in_iE_iE_H(N)\mu N}{kT}\right]$$
(12)
The lack of diffusive contact between the quantum dots and the rest of the device means that the chemical potential, $`\mu `$, is determined by electrons in the leads and contacts. The summation in $`Z(N)`$ is carried out over all electron configurations $`\{n_i\}`$ for which $`_in_i=N`$. Double counting the Coulombic interaction is avoided by subtracting the Hartree energy $`E_H(N)`$ for the $`N`$ electrons. The Hartree energy appearing in the partition function is
$$E_H(N)=\frac{1}{8\pi \epsilon }\frac{n_e(\stackrel{}{r})n_e(\stackrel{}{r}^{})}{|\stackrel{}{r}\stackrel{}{r}^{}|}𝑑\stackrel{}{r}𝑑\stackrel{}{r}^{}$$
(13)
where $`n_e(\stackrel{}{r})`$ is the charge in the quantum dot and the integration is restricted to the dot region. Directly solving for the Hartree energy by performing a double integral over the quantum dot charge density is too time consuming and impractical due to the presence of the singularity in the integrand at $`\stackrel{}{r}=\stackrel{}{r}^{}`$. An alternative method of calculating the Hartree energy is to use the equivalent expression
$$E_H(N)=\frac{1}{8\pi \epsilon }V_e(\stackrel{}{r})n_e(\stackrel{}{r})𝑑\stackrel{}{r}$$
(14)
where the potential $`V_e(\stackrel{}{r})`$ is obtained by solving the Poisson equation in the pillar using the charge density, $`n_e(\stackrel{}{r})`$, in the quantum dot.
The boundary condition on $`V_e(\stackrel{}{r})`$ at the surface of the pillar is determined by asymptotically expanding $`V_e(\stackrel{}{r})`$ in a multipole expansion in the quantum dot charge density up through quadrupole terms and using this expansion to specify $`V_e`$ at the surface. This is a good approximation since the pillar boundaries are far from the localized quantum dot charge.
To obtain a self-consistent solution to the coupled Poisson and Schrödinger equations, we first specify the device structure including the $`Al_xGa_{1x}As`$ alloy composition, the doping concentration, and the arrangement of the electrodes. In all our runs, the source and drain are assumed to be grounded and the gate is assumed to be negatively biased. We initially assume complete depletion in the structure and solve the Poisson equation to obtain an initial guess for the electrostatic potential, $`V(\rho ,z)`$. With this electrostatic potential and the quantum well band offset potentials, we solve the Schrödinger equation for the unoccupied quantum dot energies and wavefunctions. The chemical potential in the depleted structure is set to the minimum of the Thomas-Fermi electron potential, $`U=qV(\rho ,z)+\mathrm{\Delta }E_c`$, in the source and drain regions. Starting with these initial guesses for the chemical potential in the leads and the solutions of the Poisson and Schrödinger equations, we obtain self-consistent solutions through the following relaxation procedure.
First, electron densities in the leads and the quantum dots are obtained from the chemical potential, the temperature, and the quantum dot electronic states. The global charge density, including the given doping charge, is then obtained. Next the Poisson equation is solved for $`V(\rho ,z)`$. The Hartree potential and exchange-correlation potentials are then obtained from $`V(\rho ,z)`$ and the electron charge density. With the electron potentials in hand, the Schröodinger equation is solved in each quantum dot region. The procedure is then repeated until convergence is achieved. In updating the electrostatic potential and electron charge density, the new solutions are mixed with the old to obtain the updated solutions. For the electrostatic potential
$$V(\rho ,z)\lambda V_{new}(\rho ,z)+(1\lambda )V_{old}(\rho ,z)$$
(15)
where $`\lambda <1`$ is a relaxation parameter which is dynamically adjusted to accelerate convergence. A similar scheme is used to update the electron charge density. The above procedure is iterated until the chemical potential, electrostatic potential, electron charge density, and quantum dot energy levels all change by less than some small relative tolerance between successive iterations at which point convergence is achieved. Typically about 400 iterations are required to achieve convergence to within one part in $`10^4`$.
## IV A Three qubit quantum register: 1D analysis
We can use the device modeling program described in the last section to obtain a design for a three-qubit quantum register. We could, in principle, do a full 3D analysis of the device and obtain suitable design parameters (i.e., pillar dimensions, doping concentrations, asymmetric well shapes, electrode placement and biasing, etc.) based on our computationally intensive 3D model. Clearly this would be prohibitively time consuming due to the size of the parameter space that would need to be investigated as well as the time required to perform each run. To narrow down the design parameters, we can take advantage of the fact that our quantum computer is operated in the extreme depletion limit and do a simple 1D analysis to gain some useful insight.
Let’s assume that inside the core of stacked quantum wells (radius $`R_c`$) we have complete depletion and uniform doping. In this limit, the quantum dot electron potential, $`U(r)`$, can be expressed in cylindrical coordinates as $`U(\stackrel{}{r})=U(z)+U(\rho )`$, where $`U(\rho )`$ is a radial potential arising from the uniform donor density and $`U(z)`$ is the conduction band offset potential along the growth direction. This separable potential assumption is a good approximation in the strong depletion regime where only a single electron resides in each dot. The assumption of a separable potential is commonly used in the study of quantum dot structures and enables us to consider the $`z`$ and $`\rho `$ motions separately. The z-directional potential $`U(z)`$, shown schematically in the inset of Fig. 2, is a step potential formed by a layer of $`Al_xGa_{1x}As`$ of thickness $`B`$ ($`0<z<B`$) and a layer of $`GaAs`$ of thickness $`LB`$ ($`B<z<L`$). The resulting asymmetric quantum dot/well is confined by $`Al_yGa_{1y}As`$ barriers with $`y>x`$ and the asymmetry is parameterized by the ratio $`B/L`$ where $`0<B/L<1`$.
In the effective mass approximation, the qubit wavefunctions are $`|i=R(\rho )\psi _i(z)u_s(\stackrel{}{r})`$ ($`i=0,1`$). Here $`R(\rho )`$ is the ground state of the radial envelope function, $`\psi _i(z)`$ is the envelope function along $`z`$, and $`u_s(\stackrel{}{r})`$ is the $`s`$-like zone center Bloch function including electron spin. For simplicity, we assume complete confinement by the $`Al_yGa_{1y}As`$ barriers along the z direction. Then, the envelope function $`\psi _i(z)`$ is obtained by solving the time-independent Schrödinger equation subject to the boundary conditions $`\psi _i(0)=\psi _i(L)=0`$. The energies of the qubit wavefunctions are given by $`E=E_\rho +E_i`$ where $`E_\rho `$ is the energy associated with $`R(\rho )`$ and $`E_i`$ is the energy associated with $`\psi _i(z)`$.
Figure 2 shows the probability density, $`|\psi _i(z)|^2`$, as a function of position, $`z`$, for the two qubit states $`|0`$ and $`|1`$ in a $`20nm`$ $`GaAs/Al_{0.3}Ga_{0.7}As`$ asymmetric quantum dot. The barrier thickness $`B=15nm`$ and the overall length of the dot is $`L=20nm`$. By choosing $`B/L=0.75`$ and $`x=0.3`$, it is found that the ground state wavefunction $`|0`$ is strongly localized in the $`GaAs`$ region while the $`|1`$ wavefunction is strongly localized in the $`Al_{0.3}Ga_{0.7}As`$ barrier. By appropriately choosing the asymmetric quantum dot parameters, the qubit wavefunctions can be spatially separated and a large difference in the electrostatic dipole moments can be achieved.
The transition energy $`\mathrm{\Delta }E_0=E_1E_0`$ between $`|1`$ and $`|0`$ is shown in Fig. 3 as a function of $`B/L`$ in a $`20nm`$ $`GaAs/Al_xGa_{1x}As`$ asymmetric quantum dot ($`L=20nm`$). Several values of Al concentration $`x`$ are considered. In Fig. 4, we fix the Al concentration at $`x=0.2`$ and plot $`\mathrm{\Delta }E_0`$ as a function of $`L`$ for several values of $`B/L`$. The continuous curves are based on our 1D analysis and the squares are the qubit energy gaps for a three qubit self-consistent quantum register calculation as described in the next section. It is clear that the transition energy can be tailored substantially by varying the asymmetry parameter. With three parameters available for adjustment ($`B`$, $`L`$, and $`x`$), we can make $`\mathrm{\Delta }E_0`$ unique for each dot in the register. In this way, we can address a given dot by using laser light with the correct photon energy.
It is desirable that the $`|1`$ state be the first excited level of the quantum dot. Thus, the lowest lying radial state $`(0,1,0)`$ should lie above the $`|1`$ state. The radial energy gap, $`\mathrm{\Delta }E_1`$, between the ground state, $`|0`$, and the first radial excited state, $`(0,1,0)`$, is found by solving a 2D Schrödinger equation for an electron in the radial potential, $`U_r(\rho )`$. If we take the barrier in the sheath to be infinite, then in the extreme depletion limit, we have
$$U_r(\rho )=\{\begin{array}{cc}qV(\rho )\hfill & \text{ if }\rho <R_c\hfill \\ \mathrm{}\hfill & \text{ otherwise}\hfill \end{array}$$
(16)
where $`V(\rho )`$ is the radial electrostatic potential. For complete depletion and uniform doping, the Poisson equation for $`V(\rho )`$ can be solved analytically. Thus,
$$V(\rho )=\frac{\pi N_D^+}{\epsilon }(R_c^2\rho ^2)$$
(17)
where $`R_c`$ is the sheath radius and $`N_D^+`$ is the doping density. Numerically solving the 2D Schrödinger equation for an electron in the potential, $`U_r(\rho )`$, is straightforward. Figure 5 shows the radial energy gap, $`\mathrm{\Delta }E_1`$, between the $`|0`$ and the lowest lying radial state, $`(0,1,0)`$, as a function of doping concentration, $`N_D^+`$, for several values of $`R_c`$. For narrow pillars with low doping concentrations, the radial energy gap is determined by size confinement. For large pillars with high doping concentrations the radial energy gap is determined by electrostatic confinement. From Fig. 4, we see that the qubit energy gaps reach a minimum near $`\mathrm{\Delta }E_070meV`$ for quantum wells with $`L20nm`$. The results of Fig. 5 suggest that radial gaps in the range of $`\mathrm{\Delta }E_1100meV`$ with strong size confinement can be achieved with doping densities in the range of $`N_D^+10^{17}cm^3`$ if $`R_c70\AA `$.
The electric field from an electron in one dot shifts the energy levels of electrons in adjacent dots through electrostatic dipole-dipole coupling. By appropriate choice of coordinate systems, the dipole moments associated with $`|0`$ and $`|1`$ equal in magnitude but oppositely directed. The dipole-dipole coupling energy is then defined as
$$V_{dd}=2\frac{|d_1||d_2|}{ϵ_rR_{12}^3},$$
(18)
where $`d_1`$ and $`d_2`$ are the ground state dipole moments in the two dots, $`ϵ_r=12.9`$ is the dielectric constant for $`GaAs`$, and $`R_{12}`$ is the distance between the dots.
Figure 6 shows the dipole-dipole coupling energy, $`V_{dd}`$, between two asymmetric $`GaAs/Al_xGa_{1x}As`$ quantum dots of widths $`L1=19nm`$ and $`L2=21nm`$ separated by a $`10nm`$ $`Al_yGa_{1y}As`$ barrier. The coupling energy is plotted as a function of $`B/L`$ for several values of $`x`$ where $`B/L`$ and $`x`$ are taken to be the same in both dots. The dipole-dipole coupling energies are a strongly peaked function of the asymmetry parameter, $`B/L`$. From the figure, we see that values of $`V_{dd}0.15meV`$ can be achieved.
Quantum dot electrons can interact with the environment through the phonon field, particularly the longitudinal-optical (LO) and acoustic (LA) phonons. The LO phonon energy, $`\mathrm{}\omega _{LO}`$, lies in a narrow band around $`36.2meV`$. As long as the quantum dot energy level spacings lie outside this band, LO phonon scattering is strongly suppressed by the phonon bottleneck effect. Acoustic phonon energies are much smaller than the energy difference, $`\mathrm{\Delta }E`$, between qubit states. Thus, acoustic phonon scattering requires multiple emission processes which are also very slow. Theoretical studies on phonon bottleneck effects in GaAs quantum dots indicate that LO and LA phonon scattering rates including multiple phonon processes could be slower than the spontaneous emission rate provided that the quantum dot energy level spacing is greater than $`1meV`$ and, at the same time, avoids a narrow window around the LO phonon energy. In Ref. , Inoshita and Sakaki compute multi-phonon relaxation rates in spherical single-electron $`GaAs`$ quantum dots due to one- and two-phonon scattering by LO and LA phonons at $`T=0K`$ and $`T=300K`$. Using the results of this calculation, we estimate that multi-phonon scattering dominates the spontaneous emission only if the qubit energy level spacing is within $`4meV`$ of the LO phonon energy. Likewise, multi-phonon LA scattering becomes important if the qubit energy gaps are smaller than $`1meV`$
While dephasing via interactions with the phonon field can be strongly suppressed by proper designing of the structure, quantum dot electrons are still coupled to the environment through spontaneous emission and this is the dominant dephasing mechanism. Decoherence resulting from spontaneous emission ultimately limits the total time available for a quantum computation. Thus, it is important that the spontaneous emission lifetime be large. The excited state lifetime, $`T_d`$, against spontaneous emission is
$$T_d=\frac{3\mathrm{}(\mathrm{}c)^3}{4e^2D^2\mathrm{\Delta }E^3},$$
(19)
where $`D=0|z|1`$ is the dipole matrix element between $`|0`$ and $`|1`$.
Figure 7 shows the spontaneous emission lifetime of an electron in qubit state $`|1`$ for a $`20nm`$ $`GaAs/Al_xGa_{1x}As`$ quantum dot as a function of asymmetry parameter, $`B/L`$, for several values of Al concentration, $`x`$. It is immediately obvious from Fig. 7 that the lifetime depends strongly on $`B/L`$. Depending on the value of $`x`$ chosen, the computed lifetime can achieve a maximum of between 4000 $`ns`$ and 6000 $`ns`$. In general, the maximum lifetime increases with $`x`$. In Eq. (19), the lifetime is inversely proportional to $`\mathrm{\Delta }E^3`$ and $`D^2`$, but the sharp peak seen in Fig. 7 is due primarily to a pronounced minimum in $`D`$.
Based on these results, we can estimate parameters for a solid state quantum register containing a stack of several asymmetric $`GaAs/Al_{0.3}Ga_{0.7}As`$ quantum dots in the $`L20nm`$ range separated by $`10nm`$ $`Al_yGa_{1y}As`$ barriers ($`y>0.4`$). An important design goal is obtaining a large spontaneous emission lifetime and a large dipole-dipole coupling energy. From Figs. 6 and 7, we see that both can be achieved by selecting an asymmetry parameter, $`B/L=0.8`$. This gives us a spontaneous emission lifetime $`T_d=3100ns`$ and a dipole-dipole coupling energy $`V_{dd}=0.14meV`$. The transition energy between the qubit states is on the order of $`100meV`$ ($`\lambda =12.4\mu m`$). In a quantum computation, the quantum register is optically driven by a laser as described in Ref. . In our example, we require a tunable infra-red laser in the mid-$`10\mu m`$ range so we can individually address various transitions between coupled qubit states.
## V A Three qubit quantum register: 3D analysis
Using the results of our simple 1D model as a starting point, we designed a three qubit quantum register by using the self-consistent device model described in Section III. Several criteria have to be met for a viable quantum register design and the structure we obtained through trial-and-error involved tradeoffs between several design goals.
For a self-consistent quantum register calculation, we assume the parameters of the free-standing quantum dot pillar structure (shown in Fig. 1) as follows: The height of the pillar is taken to be $`=1000nm`$ while the radii of the core and sheath are taken to be $`R_c=7nm`$ and $`R=50nm`$. The drain and source contacts at the top and bottom of the pillar are grounded and a cylindrical gate with a height of $`400nm`$ is placed around the center of the pillar. Near the source and drain contacts, layers of intrinsic semiconductor serve to inhibit gate-to-source and gate-to-drain currents. The central $`600nm`$ of the pillar is uniformly n-doped with a doping concentration of $`N_D=5\times 10^{17}cm^3`$.
The cylindrical sheath surrounding the core region is composed of high band gap $`Al_{0.45}Ga_{0.55}As`$ and serves to confine electrons to the core region. The three qubits in the core are defined by the composition profile of $`Al_xGa_{1x}As`$ along the pillar axis. In our structure, the Al concentration, $`x`$, in the core region is uniform in the radial direction. The composition profile along the pillar axis in the core region is shown in Fig. 8. The ground and first excited electronic states are the qubit states $`|0`$ and $`|1`$ and the electron charge densities for these states are shown schematically in the figure. We find that in thermal equilibrium the electrons reside entirely in the ground state $`|0`$ for temperatures as high as $`77K`$ since the energy gap between $`|0`$ and $`|1`$ is much greater than $`kT`$. This is indicated schematically by the solid circles in the diagram. The widths $`L`$ of the asymmetric quantum wells/dots defining qubits $`1`$ through $`3`$ are $`19.0`$, $`20.5`$ and $`22.0nm`$ while the $`B/L`$ ratios are $`0.670`$, $`0.683`$ and $`0.675`$ respectively. Our 1D analysis suggests that asymmetry parameters in this range will result in long spontaneous emission lifetimes and strong dipole-dipole coupling between neighboring qubits. The asymmetric quantum wells are composed of $`GaAs`$ and $`Al_{0.2}Ga_{0.8}As`$ layers and the barriers between the asymmetric dots/wells are composed of $`Al_{0.45}Ga_{0.55}As`$.
With a properly chosen reverse gate bias, $`V_g`$, the doping region in the center of the pillar is depleted and the equilibrium Fermi level lines up so that there is exactly one electron in each dot. Single electron occupancy in the dots is necessary in order for there to be a well defined qubit Hilbert space. Due to shell filling effects, single electron occupancy in all three dots holds over a finite range of the gate voltage. By running our device model for several values of $`V_g`$, we find that single electron occupancy is obtained over the range $`1.56VV_g1.48V`$. Thus, the requirement for single electron occupancy in the quantum dots is maintained in the presence of gate voltage fluctuations on the order of $`\mathrm{\Delta }V_g0.08V`$. For $`V_g=1.5V`$, the self-consistent electron potential along the pillar axis, (i.e. $`\rho =0`$) is shown in Fig. 9 as a function of position along the pillar axis. The position along the pillar axis is measured from the drain contact at $`z=0nm`$ to the source contact at $`z=1000nm`$. Figure 9 is centered on the active region of the register containing the three quantum dots and the origin of the energy scale is chosen to be the equilibrium Fermi level. The total electron potential is approximately the sum of the self-consistent electrostatic Hartree potential and the conduction band offset potential, the self-consistent exchange-correlation potential being negligible.
The self-consistent electron levels are obtained by solving the Schrödinger equation in the self-consistent potential shown in Fig. 9. In our structure, the $`|0`$ ground states have $`(n_\rho ,n_\varphi ,n_z)=(0,0,0)`$ symmetry and the $`|1`$ states (the first excited level) in all three qubits are $`(n_\rho ,n_\varphi ,n_z)=(0,0,1)`$ states. The self-consistent qubit energy gap, $`\mathrm{\Delta }E_0`$, between the $`|0`$ and $`|1`$ states, the radial energy gap, $`\mathrm{\Delta }E_1`$, and the spontaneous emission lifetime of the $`|1`$ state, $`\tau _s`$, and dipole moment, $`d`$, for the three qubits are listed in Table I. From Table I, we see that the radial energy gaps are larger than the qubit energy gaps. Another thing to note is that the qubit energy gaps are large compared to $`kT`$ at $`T=77K`$. Thus, in thermal equilibrium the electrons reside entirely in the $`|0`$ level at $`77K`$. This means that the initial state of the quantum register is characterized by a pure state density matrix $`\widehat{\rho }_0=|0,0,00,0,0|`$. Consequently, there is no need for initial state preparation in our quantum register. In Fig. 10, the self-consistent electron probability densities in the three quantum dots are plotted as a function of position along the pillar axis. Each dot traps one electron and the probability densities in the ground and first excited states are shown as solid and dot-dashed lines, respectively. The barriers are thick enough so that electron wavefunctions in adjacent dots do not overlap.
The energy levels for the three qubit quantum computer are shown in Table II with and without the inclusion of dipole-dipole coupling between the qubits. From Table II we see that a different energy corresponds to each three-electron state $`|i_1,i_2,i_3`$ of the register where $`i_n=(0,1)`$ labels the state of the $`n`$-th qubit. Transition energies between the states $`|0`$ and $`|1`$ for a given qubit are obtained by taking differences between the appropriate entries in Table II. For the first qubit, we take differences between all three-particle states $`|0,i_2,i_3`$ and $`|1,i_2,i_3`$. In general, the transition energy between $`|0`$ and $`|1`$ for an electron in the first qubit will depend on the states, $`i_2`$ and $`i_3`$, occupied by the second and third qubits, and there can be as many as four such conditional transitions. In the absence of dipole-dipole coupling between qubits, all four conditional transition energies between $`|0`$ and $`|1`$ for a given qubit are degenerate. When dipole-dipole interactions between the qubits are considered, the four-fold degenerate conditional transition energies split into multiplets depending on which states are occupied by the electrons in neighboring qubits.
The conditional transition energies between $`|0`$ and $`|1`$ states for our three qubit register are shown in Fig. 11 as a function of photon energy. In the absence of dipole-dipole coupling, the transition energies for the three qubits are $`40.86meV`$, $`47.14meV`$, and $`52.31meV`$, respectively. When dipole-dipole interactions between qubits are taken into account, the conditional transition energies split into multiplets as shown in this figure. Each transition in the spectrum is labeled by the neighboring electron states which give rise to it. By performing optical $`\pi `$-pulses at selected conditional transition frequencies, quantum logic operations can be performed. For example, a $`\pi `$ pulse performed on the lowest energy transition in Fig. 11 performs a bit flip on the first qubit provided the second qubit is in state $`|1`$. This operation is just a Controlled-Not gate with qubit 2 as the control bit and qubit 1 as the target bit.
The need to selectively perform $`\pi `$-pulses at the conditional transition frequencies allows us to make some preliminary estimates on the parameters of the laser system needed to drive a quantum computation. If we want to selectively drive a given transition without exciting neighboring transitions, then the bandwidth of the $`\pi `$-pulse needs to be less than the splitting between the two most closely spaced lines in the spectrum. From Fig. 11, the two most closely spaced lines are spaced $`\mathrm{\Delta }\mathrm{}\omega 0.0776meV`$ apart. If we require that the $`\pi `$-pulse bandwidth is $`\mathrm{\Delta }E_\pi 0.01meV`$, then the pulse length can be estimated from Heisenberg’s uncertainty principle, $`\mathrm{\Delta }E_\pi \mathrm{\Delta }T_\pi \mathrm{}/2`$, to be $`T_\pi 33ps`$. If we assume a square $`\pi `$-pulse, the magnitude of the optical electric field is given by
$$E_0=\frac{\pi \mathrm{}}{qdT_\pi }$$
(20)
where $`d`$ is the optical dipole from Table I and the average Poynting vector during the pulse is
$$S_{av}=\frac{1}{2}cϵ_0E_0^2$$
(21)
For $`d10\AA `$ and $`T_\pi 33ps`$, we obtain $`E_00.627kV/cm`$ and $`S_{av}522W/cm^2`$.
## VI Summary
In this paper, we have studied a solid state implementation of quantum computing based on coupled quantum dots. Our quantum register consists of a free standing n-type pillar with grounding leads at the top and bottom of the structure. Asymmetric quantum wells confine electrons along the pillar axis and a high bandgap $`AlGaAs`$ sheath wrapped around the center of the pillar allows for confinement in the radial direction. The ground and first excited electronic states of the quantum dots act as qubit states $`|0`$ and $`|1`$, respectively. We have developed a 3D device model for a general SET structure containing several quantum dots. We self-consistently solve coupled Schrödinger and Poisson equations for the device and develop a design for a three qubit quantum register with asymmetric quantum dots tailored for long dephasing time and large dipole-dipole coupling between the dots. Our results indicate that a single gate electrode can be used to localize a single electron in each of the quantum dots. Adjacent dots are strongly coupled by electric dipole-dipole interactions arising from the dot asymmetry thus enabling rapid computation rates.
###### Acknowledgements.
This study was supported, in part, by the Defense Advanced Research Project Agency and the Office of Naval Research. |
warning/0001/math0001102.html | ar5iv | text | # Random almost holomorphic sections of ample line bundles on symplectic manifolds
## Introduction
This paper is concerned with asymptotically holomorphic sections of ample line bundles over almost-complex symplectic manifolds $`(M,J,\omega ).`$ Such line bundles and sections are symplectic analogues of the usual objects in complex algebraic geometry. Interest in their properties has grown in recent years because of their use by Donaldson \[Don1, Don2\], Auroux \[Aur1, Aur2\] and others \[AuKa, Sik\] in proving symplectic analogues of standard results in complex geometry. These results involve properties of asymptotically holomorphic sections of high powers of the bundle, particularly those involving their zero sets and the maps they define to projective space.
We take up the study of asymptotically holomorphic sections from the viewpoint of the microlocal analysis of the $`\overline{}`$ operator on a symplectic almost-complex manifold, and define the class of ‘almost holomorphic sections’ by a method due to Boutet de Monvel and Guillemin \[Bout, BoGu\]. Sections of powers of a complex line bundle $`L^NM`$ over $`M`$ are identified with equivariant functions $`s`$ on the associated $`S^1`$-bundle $`X`$, and the $`\overline{}`$ operator is identified with the $`\overline{}_b`$ operator on $`X`$. In the non-integrable almost-complex symplectic case there are in general no solutions of $`\overline{}_bs=0`$. To define an ‘almost holomorphic section’ $`s`$, Boutet de Monvel and Guillemin define a certain (pseudodifferential) $`\overline{D}_j`$-complex over $`X`$ \[BoGu\] \[BoSj\]. The space $`H_J^0(M,L^N)`$ of almost holomorphic sections is then defined as the space of sections corresponding to solutions of $`\overline{D}_0s=0`$. The operator $`\overline{D}_0`$ is not uniquely or even canonically defined, and it is difficult to explicitly write down these almost holomorphic sections. The importance, and we hope usefulness, of these sections lies in the fact that they typically have the properties of asymptotically holomorphic sections as defined by Donaldson and Auroux, as we describe below. We use the term ‘almost holomorphic’ to emphasize that a priori, they are distinct from ‘asymptotically holomorphic’ sections.
Our main results involve the ‘typical’ behavior of almost holomorphic sections in a probabilistic sense, as in our work with Bleher \[BSZ1, BSZ2\] and our prior work \[ShZe\] on holomorphic sections. A wide variety of measures could be envisioned here, and much of what we do is independent of the precise choice of measure. However, the simplest measures are the Haar measures on the unit spheres in the spaces $`H_J^0(M,L^N)`$. To be precise, we use a hermitian metric $`h`$ on $`L`$ and the volume form $`dV=\frac{\omega ^m}{m!}`$ on $`M`$ to endow $`H_J^0(M,L^N)`$ with an $`^2`$ inner product. We denote by $`SH_J^0(M,L^N)`$ the elements of unit norm in $`H_J^0(M,L^N)`$ and by $`\nu _N`$ the Haar probability measure on the sphere $`SH_J^0(M,L^N)`$. We also consider the essentially equivalent Gaussian measures on $`H_J^0(M,L^N)`$. The theme of our work is to obtain results about almost holomorphic sections by calculating asymptotically (as $`N\mathrm{}`$) the probabilities that sections $`s_NH_J^0(M,L^N)`$ do various things. This theme has a variety of potential applications in geometry, which we hope to pursue in the future.
In this article we focus on two applications. Our first main result gives estimates on various norms of a typical sequence of almost holomorphic sections of growing degree. Let us recall that a sequence of sections $`s_N`$ is called asymptotically holomorphic by Donaldson and Auroux \[Don1, Aur1\] if
$$s_N_{\mathrm{}}+\overline{}s_N_{\mathrm{}}=O(1),s_N_{\mathrm{}}+\overline{}s_N|_{\mathrm{}}=O(\sqrt{N}),s_N_{\mathrm{}}=O(N).$$
We will prove that almost every sequence $`\{s_N\}`$ of $`^2`$-normalized ($`s_N_2=1`$) almost holomorphic sections in the Boutet de Monvel -Guillemin sense is close to being asymptotically holomorphic in the Donaldson sense. We also let $`:𝒞^{\mathrm{}}(M,L^N(T^{}M)^k)𝒞^{\mathrm{}}(M,L^N(T^{}M)^{(k+1)})`$ denote the connection, and we write $`^k=\mathrm{}:𝒞^{\mathrm{}}(M,L^N)𝒞^{\mathrm{}}(M,L^N(T^{}M)^k)`$. We also have the decomposition $`=+\overline{}`$; note that here $`\overline{}`$ depends on the choice of connection.
###### Theorem 0.1.
Endow the infinite product $`\mathrm{\Pi }_{N=1}^{\mathrm{}}SH_J^0(M,L^N)`$ with the product spherical measure $`\nu _{\mathrm{}}:=\mathrm{\Pi }_{N=1}^{\mathrm{}}\nu _N.`$ Then $`\nu _{\mathrm{}}`$-almost every sequence $`\{s_N\}`$ of sections satisfies the following estimates:
$$\begin{array}{cc}s_N_{\mathrm{}}=O(\sqrt{\mathrm{log}N}),\hfill & ^ks_N_{\mathrm{}}=O(N^{k/2}\sqrt{\mathrm{log}N}),\hfill \\ & \\ \overline{}s_N_{\mathrm{}}=O(\sqrt{\mathrm{log}N}),\hfill & ^k\overline{}s_N_{\mathrm{}}=O(N^{k/2}\sqrt{\mathrm{log}N}),(k1).\hfill \end{array}$$
Our second main result concerns the joint probability distribution
$$𝐃_{(z^1,\mathrm{},z^n)}^N=D^N(x,\xi ;z^1,\mathrm{},z^n)dxd\xi $$
of the random variables
$$x^p=s_N(z^p),\xi ^p=N^{\frac{1}{2}}s_N(z^p)(1pn)$$
(1)
on $`SH^0(M,L^N)`$, for $`n`$ distinct points $`z^1,\mathrm{},z^nM`$. We prove that upon rescaling, this joint probability distribution has a universal limit which agrees with that of the holomorphic case determined in \[BSZ2\].
###### Theorem 0.2.
Let $`L`$ be a pre-quantum line bundle over a $`2m`$-dimensional compact integral symplectic manifold $`(M,\omega )`$. Let $`P_0M`$ and choose complex local coordinates $`\{z_j\}`$ centered at $`P_0`$ so that $`\omega |_{P_0}`$ and $`g|_{P_0}`$ are the usual Euclidean Kähler form and metric respectively. Then
$$𝐃_{(z^1/\sqrt{N},\mathrm{},z^n/\sqrt{N})}^N𝐃_{(z^1,\mathrm{},n^n)}^{\mathrm{}}$$
where $`𝐃_{(z^1,\mathrm{},z^n)}^{\mathrm{}}`$ is a universal Gaussian measure supported on the holomorphic 1-jets.
A technically interesting novelty in the proof is the role of the $`\overline{}`$ operator. In the holomorphic case, $`𝐃_{(z^1,\mathrm{},z^n)}^N`$ is supported on the subspace of sections satisfying $`\overline{}s=0`$. In the almost complex case, sections do not satisfy this equation, so $`𝐃_{(z^1,\mathrm{},z^n)}^N`$ is a measure on a higher-dimensional space of jets. However, Theorem 0.2 says that the mass in the ‘$`\overline{}`$-directions’ shrinks to zero as $`N\mathrm{}`$.
An alternate statement of Theorem 0.2 involves equipping $`H^0(M,L^N)`$ with a Gaussian measure, and letting $`\stackrel{~}{𝐃}_{(z^1,\mathrm{},z^n)}^N`$ be the corresponding joint probability distribution on $`H_J^0(M,L^N)`$, which is a Gaussian measure on the complex vector space of 1-jets of sections. We show (Theorem 5.4) that these Gaussian measures $`\stackrel{~}{𝐃}^N`$ also have the same scaling limit $`𝐃^{\mathrm{}}`$, so that asymptotically the probabilities are the same as in the holomorphic case, as established in \[BSZ2\]. To be more precise, recall that a Gaussian measure on $`^n`$ is a measure of the form
$$\gamma _\mathrm{\Delta }=\frac{e^{\frac{1}{2}\mathrm{\Delta }^1x,x}}{(2\pi )^{n/2}\sqrt{det\mathrm{\Delta }}}dx,$$
where $`\mathrm{\Delta }`$ is a positive definite symmetric $`n\times n`$ matrix. It is then easy to see that $`\stackrel{~}{𝐃}_{(z^1,\mathrm{},z^n)}^N=\gamma _{\mathrm{\Delta }^N}`$ where $`\mathrm{\Delta }^N`$ is the covariance matrix of the random variables in (1). To deal with singular measures, we introduce in §5 generalized Gaussians whose covariance matrices are only semi-positive definite. A generalized Gaussian is simply a Gaussian supported on the subspace corresponding to the positive eigenvalues of the covariance matrix. The main step in the proof is to show that the covariance matrices $`\mathrm{\Delta }^N`$ underlying $`\stackrel{~}{𝐃}^N`$ tend in the scaling limit to a semi-positive matrix $`\mathrm{\Delta }^{\mathrm{}}`$. It follows that the scaled distributions $`\stackrel{~}{𝐃}^N`$ tend to a generalized Gaussian $`\gamma _\mathrm{\Delta }^{\mathrm{}}`$ ‘vanishing in the $`\overline{}`$-directions.’
In joint work with Bleher \[BSZ3\], we use this result to prove universality of scaling limits of correlations of zeros in the setting of almost holomorphic sections over almost-complex symplectic manifolds. The analysis underlying Theorem 0.2 should also be useful for calculating many other kinds of probabilities in the setting of asymptotically holomorphic sections. For instance, we believe it should be useful for proving existence results for asymptotically holomorphic sections satisfying transversality conditions.
These results are based on two essential analytical results which have an independent interest and which we believe will have future applications. The first is the scaling asymptotics of the Szegö kernels $`\mathrm{\Pi }_N(z,w)`$, i.e. the orthogonal projections onto $`H_J^0(M,L^N)`$. To be more precise, we lift the Szegö kernels to $`X`$ and the asymptotics are as follows:
> Choose local coordinates $`\{z_j\}`$ centered at a point $`P_0M`$ as in Theorem 0.2 and choose a ‘preferred’ local frame for $`L`$, which together with the coordinates on $`M`$ give us ‘Heisenberg coordinates’ on $`X`$ (see §1.2). We then have
>
> $$\begin{array}{c}N^m\mathrm{\Pi }_N(P_0+\frac{u}{\sqrt{N}},\frac{\theta }{N};P_0+\frac{v}{\sqrt{N}},\frac{\phi }{N})\hfill \\ \\ =\frac{1}{\pi ^m}e^{i(\theta \phi )+u\overline{v}\frac{1}{2}(|u|^2+|v|^2)}\left[1+_{r=1}^KN^{\frac{r}{2}}b_r(P_0,u,v)+N^{\frac{K+1}{2}}R_K(P_0,u,v,N)\right],\hfill \end{array}$$
> (2)
> where $`R_K(P_0,u,v,N)_{𝒞^j(\{|u|+|v|\rho \})}C_{K,j,\rho }`$ for $`j=1,2,3,`$.
A more precise statement will be given in Theorem 2.3. As more or less immediate corollaries of these scaling asymptotics, we prove symplectic analogues of the holomorphic Kodaira embedding theorem and Tian almost-isometry theorem \[Tian\]; these two results have previously been proved by Borthwick-Uribe \[BoUr1, BoUr2\] using a related microlocal approach. The Borthwick-Uribe proof of the almost-complex Tian theorem was in turn motivated by a similar proof in the holomorphic case in \[Zel\].
The proof of (2) is based on our second analytic result: the construction of explicit parametrices for $`\mathrm{\Pi }`$ and its Fourier coefficients $`\mathrm{\Pi }_N`$. These parametrices closely resemble those of Boutet de Monvel - Sjöstrand \[BoSj\] in the holomorphic case. The construction is new but closely follows the work of Menikoff and Sjöstrand \[MenSj, Sjö\] and of Boutet de Monvel and Guillemin \[Bout, BoGu\]. For the sake of completeness, we will give a fairly detailed exposition of the construction of the zeroth term of the $`\overline{D}_j`$ complex and of the Szegö kernel.
## Guide for the reader
For the readers’ convenience, we provide here a brief outline of the paper. We begin in §1 by first describing some terminology from symplectic geometry and then giving an outline of Boutet de Monvel and Guillemin’s construction \[Bout, BoGu\] of a complex of pseudodifferential operators, which replaces the $`\overline{}_b`$ complex in the symplectic setting. The zeroth term of this complex is used to define sequences of almost holomorphic sections and Szegö projectors analogous to the integrable complex case (§1.3). In §2, we show that the Szegö projectors $`\mathrm{\Pi }_N`$ are complex Fourier integral operators of the same type as in the holomorphic case, and we use this formulation to obtain the scaling asymptotics of $`\mathrm{\Pi }_N(P_0+\frac{u}{\sqrt{N}},\frac{\theta }{N};P_0+\frac{v}{\sqrt{N}},\frac{\phi }{N})`$. Section 3 gives two applications of these asymptotics: a proof of a ‘Kodaira embedding theorem’ (using global almost holomorphic sections) for integral symplectic manifolds, and a generalization of the asymptotic expansion theorem of \[Zel\] to symplectic manifolds. Section 4 uses the scaling asymptotics to prove that sequences of almost holomorphic sections are almost surely (in the probabilistic sense) asymptotically close to holomorphic (Theorem 0.1). Finally, in §5, we determine the joint probability distributions $`𝐃^N,\stackrel{~}{𝐃}^N`$ and again apply the scaling asymptotics to prove Theorem 0.2.
The following chart shows the interdependencies of the sections:
$$\begin{array}{cccccc}& & \mathrm{\S }\text{1}& & & \\ & & & & & \\ \mathrm{\S }\text{3}& & \mathrm{\S }\text{2}& & \mathrm{\S }\text{4}& \\ & & & & & \\ & & \mathrm{\S }\text{5}& & & \end{array}$$
We advise the reader who wishes to proceed quickly to the applications in §§35 that these sections depend only on the scaling asymptotics of the Szegö kernel stated in Theorem 2.3 and the notation and terminology given in §§1.11.3.
## 1. Circle bundles and almost CR geometry
We denote by $`(M,\omega )`$ a compact symplectic manifold such that $`[\frac{1}{\pi }\omega ]`$ is an integral cohomology class. As is well known (cf. \[Woo, Prop. 8.3.1\]; see also \[GuSt\]), there exists a hermitian line bundle $`(L,h)M`$ and a metric connection $``$ on $`L`$ whose curvature $`\mathrm{\Theta }_L`$ satisfies $`\frac{i}{2}\mathrm{\Theta }_L=\omega `$. We denote by $`L^N`$ the $`N^{\mathrm{th}}`$ tensor power of $`L`$. The ‘quantization’ of $`(M,\omega )`$ at Planck constant $`1/N`$ should be a Hilbert space of polarized sections of $`L^N`$ (\[GuSt, p. 266\]). In the complex case, polarized sections are simply holomorphic sections. The notion of polarized sections is problematic in the non-complex symplectic setting, since the Lagrangean subbundle $`T^{1,0}M`$ defining the complex polarization is not integrable and there usually are no ‘holomorphic’ sections. A subtle but compelling replacement for the notion of polarized section has been proposed by Boutet de Monvel and Guillemin \[Bout, BoGu\], and it is this notion which we describe in this section. For the asymptotic analysis, it is best to view sections of $`L^N`$ as functions on the unit circle bundle $`XL^{}`$; we shall describe the ‘almost CR geometry’ of $`X`$ in §1.2 below.
### 1.1. Almost complex symplectic manifolds
We begin by reviewing some terminology from almost complex symplectic geometry. An almost complex symplectic manifold is a symplectic manifold $`(M,\omega )`$ together with an almost complex structure $`J`$ satisfying the compatibility condition $`\omega (Jv,Jw)=\omega (v,w)`$ and the positivity condition. $`\omega (v,Jv)>0`$. We give $`M`$ the Riemannian metric $`g(v,w)=\omega (v,Jw)`$. We denote by $`T^{1,0}M,`$ resp. $`T^{0,1}M`$, the holomorphic, resp. anti-holomorphic, sub-bundle of the complex tangent bundle $`TM`$; i.e., $`J=i`$ on $`T^{1,0}M`$ and $`J=i`$ on $`T^{0,1}M`$. We give $`M`$ local coordinates $`(x_1,y_1,\mathrm{},x_m,y_m)`$, and we write $`z_j=x_j+iy_j`$. As in the integrable (i.e., holomorphic) case, we let $`\{\frac{}{z_j},\frac{}{\overline{z}_j}\}`$ denote the dual frame to $`\{dz_j,d\overline{z}_j\}`$. Although in our case, the coordinates $`z_j`$ are not holomorphic and consequently $`\frac{}{z_j}`$ is generally not in $`T^{1,0}M`$, we nonetheless have
$$\frac{}{z_j}=\frac{1}{2}\frac{}{x_j}\frac{i}{2}\frac{}{y_j},\frac{}{\overline{z}_j}=\frac{1}{2}\frac{}{x_j}+\frac{i}{2}\frac{}{y_j}.$$
At any point $`P_0M`$, we can choose a local frame $`\{\overline{Z}_1^M,\mathrm{},\overline{Z}_m^M\}`$ for $`T^{0,1}M`$ near $`P_0`$ and coordinates about $`P_0`$ so that
$$\overline{Z}_j^M=\frac{}{\overline{z}_j}+\underset{k=1}{\overset{m}{}}B_{jk}(z)\frac{}{z_k},B_{jk}(P_0)=0,$$
(3)
and hence $`/z_j|_{P_0}T^{1,0}(M)`$. This is one of the properties of our ‘preferred coordinates’ defined below.
Definition: Let $`P_0M`$. A coordinate system $`(z_1,\mathrm{},z_m)`$ on a neighborhood $`U`$ of $`P_0`$ is preferred at $`P_0`$ if
$$\underset{j=1}{\overset{m}{}}dz_jd\overline{z}_j=(gi\omega )|_{P_0}.$$
In fact, the coordinates $`(z_1,\mathrm{},z_m)`$ are preferred at $`P_0`$ if an only if any two of the following conditions (and hence all three) are satisfied:
1. $`/z_j|_{P_0}T^{1,0}(M)`$, for $`1jm`$,
2. $`\omega (P_0)=\omega _0`$,
3. $`g(P_0)=g_0`$,
where $`\omega _0`$ is the standard symplectic form and $`g_0`$ is the Euclidean metric:
$$\omega _0=\frac{i}{2}\underset{j=1}{\overset{m}{}}dz_jd\overline{z}_j=\underset{j=1}{\overset{m}{}}(dx_jdy_jdy_jdx_j),g_0=\underset{j=1}{\overset{m}{}}(dx_jdx_j+dy_jdy_j).$$
(To verify this statement, note that condition (i) is equivalent to $`J(dx_j)=dy_j`$ at $`P_0`$, and use $`g(v,w)=\omega (v,Jw)`$.) Note that by the Darboux theorem, we can choose the coordinates so that condition (ii) is satisfied on a neighborhood of $`P_0`$, but this is not necessary for our scaling results.
### 1.2. The circle bundle and Heisenberg coordinates
We now let $`(M,\omega ,J)`$ be a compact, almost complex symplectic manifold such that $`[\frac{1}{\pi }\omega ]`$ is an integral cohomology class, and we choose a hermitian line bundle $`(L,h)M`$ and a metric connection $``$ on $`L`$ with $`\frac{i}{2}\mathrm{\Theta }_L=\omega `$. In order to simultaneously analyze sections of all positive powers $`L^N`$ of the line bundle $`L`$, we work on the associated principal $`S^1`$ bundle $`XM`$, which is defined as follows: let $`\pi :L^{}M`$ denote the dual line bundle to $`L`$ with dual metric $`h^{}`$, and put $`X=\{vL^{}:v_h^{}=1\}`$. We let $`\alpha `$ be the the connection 1-form on $`X`$ given by $``$; we then have $`d\alpha =\pi ^{}\omega `$, and thus $`\alpha `$ is a contact form on $`X`$, i.e., $`\alpha (d\alpha )^m`$ is a volume form on $`X`$.
We let $`r_\theta x=e^{i\theta }x`$ ($`xX`$) denote the $`S^1`$ action on $`X`$ and denote its infinitesimal generator by $`\frac{}{\theta }`$. A section $`s`$ of $`L`$ determines an equivariant function $`\widehat{s}`$ on $`L^{}`$ by the rule $`\widehat{s}(\lambda )=(\lambda ,s(z))`$ ($`\lambda L_z^{},zM`$). It is clear that if $`\tau `$ then $`\widehat{s}(z,\tau \lambda )=\tau \widehat{s}`$. We henceforth restrict $`\widehat{s}`$ to $`X`$ and then the equivariance property takes the form $`\widehat{s}(r_\theta x)=e^{i\theta }\widehat{s}(x)`$. Similarly, a section $`s_N`$ of $`L^N`$ determines an equivariant function $`\widehat{s}_N`$ on $`X`$: put
$$\widehat{s}_N(\lambda )=(\lambda ^N,s_N(z)),\lambda X_z,$$
(4)
where $`\lambda ^N=\lambda \mathrm{}\lambda `$; then $`\widehat{s}_N(r_\theta x)=e^{iN\theta }\widehat{s}_N(x)`$. We denote by $`_N^2(X)`$ the space of such equivariant functions transforming by the $`N^{\mathrm{th}}`$ character.
In the complex case, $`X`$ is a CR manifold. In the general almost-complex symplectic case it is an almost CR manifold. The almost CR structure is defined as follows: The kernel of $`\alpha `$ defines a horizontal hyperplane bundle $`HTX`$. Using the projection $`\pi :XM`$, we may lift the splitting $`TM=T^{1,0}MT^{0,1}M`$ to a splitting $`H=H^{1,0}H^{0,1}`$. The almost CR structure on $`X`$ is defined to be the splitting $`TX=H^{1,0}H^{0,1}\frac{}{\theta }`$. We also consider a local orthonormal frame $`Z_1,\mathrm{},Z_n`$ of $`H^{1,0}`$ , resp. $`\overline{Z}_1,\mathrm{},\overline{Z}_m`$ of $`H^{0,1}`$, and dual orthonormal coframes $`\vartheta _1,\mathrm{},\vartheta _m,`$ resp. $`\overline{\vartheta }_1,\mathrm{},\overline{\vartheta }_m`$. On the manifold $`X`$ we have $`d=_b+\overline{}_b+\frac{}{\theta }\alpha `$, where $`_b=_{j=1}^m\vartheta _jZ_j`$ and $`\overline{}_b=_{j=1}^m\overline{\vartheta }_j\overline{Z}_j`$. We define the almost-CR $`\overline{}_b`$ operator by $`\overline{}_b=df|_{H^{1,0}}`$. Note that for an $`^2`$ section $`s^N`$ of $`L^N`$, we have
$$(_{L^N}s^N)\widehat{}=d^h\widehat{s}^N,$$
(5)
where $`d^h=_b+\overline{}_b`$ is the horizontal derivative on $`X`$.
Our near-diagonal asymptotics of the Szegö kernel (§2.2) are given in terms of the Heisenberg dilations, using local ‘Heisenberg coordinates’ at a point $`x_0X`$. To describe these coordinates, we first need the concept of a ‘preferred frame’:
Definition: A preferred frame for $`LM`$ at a point $`P_0M`$ is a local frame $`e_L`$ in a neighborhood of $`P_0`$ such that
1. $`e_L_{P_0}=1`$;
2. $`e_L|_{P_0}=0`$;
3. $`^2e_L|_{P_0}=(g+i\omega )e_L|_{P_0}T_M^{}T_M^{}L`$.
(A preferred frame can be constructed by multiplying an arbitrary frame by a function with specified 2-jet at $`P_0`$; any two such frames agree to third order at $`P_0`$.) Once we have property (ii), property (iii) is independent of the choice of connection on $`T_M^{}`$ used to define $`:𝒞^{\mathrm{}}(M,LT_M^{})𝒞^{\mathrm{}}(M,LT_M^{}T_M^{})`$. In fact, property (iii) is a necessary condition for obtaining universal scaling asymptotics, because of the ‘parabolic’ scaling in the Heisenberg group. Note that if $`e_L`$ is a preferred frame at $`P_0`$ and if $`(z_1,\mathrm{},z_m)`$ are preferred coordinates at $`P_0`$, then we compute the Hessian of $`e_L`$:
$$(^2e_L_h)_{P_0}=\mathrm{}(^2e_L,e_L)_{P_0}=g(P_0);$$
thus if the preferred coordinates are ‘centered’ at $`P_0`$ (i.e., $`P_0=0`$), we have
$$e_L_h=1\frac{1}{2}|z|^2+O(|z|^3).$$
(6)
Remark: Recall (\[BSZ2, §1.3.2\]) that the Bargmann-Fock representation of the Heisenberg group acts on the space of holomorphic functions on $`(M,\omega )=(^m,\omega _0)`$ that are square integrable with respect to the weight $`h=e^{|z^2|}`$. We let $`L=^m\times `$ be the trivial bundle. Then the trivializing section $`e_L(z):=(z,1)`$ is a preferred frame at $`P_0=0`$ with respect to the Hermitian connection $``$ given by
$$e_L=\mathrm{log}he_L=\underset{j=1}{\overset{m}{}}\overline{z}_jdz_je_L.$$
Indeed, the above yields $`^2e_L|_0=d\overline{z}_jdz_je_L(0)=(g_0+i\omega _0)e_L(0)`$.
The preferred frame and preferred coordinates together give us ‘Heisenberg coordinates’:
Definition: A Heisenberg coordinate chart at a point $`x_0`$ in the principal bundle $`X`$ is a coordinate chart $`\rho :UV`$ with $`0U^m\times `$ and $`\rho (0)=x_0VX`$ of the form
$$\rho (z_1,\mathrm{},z_m,\theta )=e^{i\theta }a(z)^{\frac{1}{2}}e_L^{}(z),$$
(7)
where $`e_L`$ is a preferred local frame for $`LM`$ at $`P_0=\pi (x_0)`$, and $`(z_1,\mathrm{},z_m)`$ are preferred coordinates centered at $`P_0`$. (Note that $`P_0`$ has coordinates $`(0,\mathrm{},0)`$ and $`e_L^{}(P_0)=x_0`$.)
We now give some computations using local coordinates $`(z_1,\mathrm{},z_m,\theta )`$ of the form (7) for a local frame $`e_L`$. (For the moment, we do not assume they are Heisenberg coordinates.) We write
$`a(z)`$ $`=`$ $`e_L^{}(z)_h^{}^2=e_L(z)_h^2,`$
$`\alpha `$ $`=`$ $`d\theta +\beta ,\beta ={\displaystyle \underset{j=1}{\overset{m}{}}}(A_jdz_j+\overline{A}_jd\overline{z}_j),`$
$`e_L`$ $`=`$ $`\phi e_L,\text{hence}e_L^N=N\phi e_L^N.`$
We let $`\frac{^h}{z_j}H^{1,0}X`$ denote the horizontal lift of $`\frac{}{z_j}`$. The condition $`(\frac{^h}{z_j},\alpha )=0`$ yields
$$\frac{^h}{z_j}=\frac{}{z_j}A_j\frac{}{\theta },\frac{^h}{\overline{z}_j}=\frac{}{\overline{z}_j}\overline{A}_j\frac{}{\theta }.$$
(8)
Suppose $`s_N=fe_L^N`$ is a local section of $`L^N`$. Then by (4) and (7),
$$\widehat{s}_N(z,\theta )=f(z)a(z)^{\frac{1}{2}}e^{iN\theta }.$$
(9)
Differentiating (9) and using (5), we conclude that
$`\phi `$ $`=`$ $`{\displaystyle \frac{1}{2}}d\mathrm{log}ai\beta `$ (10)
$`=`$ $`{\displaystyle \underset{j=1}{\overset{m}{}}}\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{log}a}{z_j}}+iA_j\right)dz_j{\displaystyle \underset{j=1}{\overset{m}{}}}\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{log}a}{\overline{z}_j}}+i\overline{A}_j\right)d\overline{z}_j.`$
Now suppose that $`(z_1,\mathrm{},z_m,\theta )`$ are Heisenberg coordinates at $`P_0`$; i.e., $`e_L`$ is a preferred frame at $`P_0`$ and $`(z_1,\mathrm{},z_m)`$ are preferred coordinates centered at $`P_0`$ (with $`P_0=0`$). By property (ii) of preferred frames, we have $`\phi (0)=0`$, and hence by (10)
$$da|_0=d\mathrm{log}a|_0=0,$$
(11)
$$A_j(0)=0,(1jm).$$
(12)
By differentiating (10) and applying the properties of preferred coordinates and frames, we further obtain
$$\underset{j=1}{\overset{m}{}}d\overline{z}_jdz_j=\phi =\underset{j=1}{\overset{m}{}}d\left(\frac{1}{2}\frac{\mathrm{log}a}{z_j}+iA_j\right)dz_j+\underset{j=1}{\overset{m}{}}d\left(\frac{1}{2}\frac{\mathrm{log}a}{\overline{z}_j}+i\overline{A}_j\right)d\overline{z}_j\text{at}0.$$
Thus the following four equations are satisfied at $`P_0=0`$:
$$\begin{array}{cccccc}\hfill \frac{1}{2}\frac{^2\mathrm{log}a}{z_jz_k}+i\frac{A_j}{z_k}& =& 0,\hfill & \hfill \frac{1}{2}\frac{^2\mathrm{log}a}{z_j\overline{z}_k}+i\frac{A_j}{\overline{z}_k}& =& \delta _k^j,\hfill \\ \hfill \frac{1}{2}\frac{^2\mathrm{log}a}{\overline{z}_jz_k}+i\frac{\overline{A}_j}{z_k}& =& 0,\hfill & \hfill \frac{1}{2}\frac{^2\mathrm{log}a}{\overline{z}_j\overline{z}_k}+i\frac{\overline{A}_j}{\overline{z}_k}& =& 0,\hfill \end{array}$$
(13)
at $`P_0`$. Solving (13) and recalling that $`a(0)=1,da|_0=0`$, we obtain
$$\frac{^2a}{z_jz_k}(0)=0,\frac{^2a}{z_j\overline{z}_k}(0)=\delta _k^j,$$
(14)
$$\frac{A_j}{z_k}(0)=0,\frac{A_j}{\overline{z}_k}=\frac{i}{2}\delta _k^j.$$
(15)
Hence $`A_j=\frac{i}{2}\overline{z}_j+O(|z|^2)`$ and
$$\frac{^h}{z_j}=\frac{}{z_j}+\left[\frac{i}{2}\overline{z}_j+O(|z|^2)\right]\frac{}{\theta },\frac{^h}{\overline{z}_j}=\frac{}{\overline{z}_j}\left[\frac{i}{2}z_j+O(|z|^2)\right]\frac{}{\theta }.$$
(16)
### 1.3. The $`\overline{D}`$ complex and Szegö kernels
In the complex case, a holomorphic section $`s`$ of $`L^N`$ lifts to a $`\widehat{s}_N^2(X)`$ which satisfying $`\overline{}_b\widehat{s}=0.`$ The operator $`\overline{}_b`$ extends to a complex satisfying $`\overline{}_b^2=0`$, which is a necessary and sufficient condition for having a maximal family of CR holomorphic coordinates. In the non-integrable case $`\overline{}_b^20`$, and there may be no solutions of $`\overline{}_bf=0.`$ To define polarized sections and their equivariant lifts, Boutet de Monvel \[Bout\] and Boutet de Monvel - Guillemin \[BoGu\] defined a complex $`\overline{D}_j`$, which is a good replacement for $`\overline{}_b`$ in the non-integrable case. Their main result is:
###### Theorem 1.1.
(see \[BoGu\], Lemma 14.11 and Theorem A 5.9) There exists an $`S^1`$-invariant complex of first order pseudodifferential operators $`\overline{D}_j`$ over $`X`$
$$0C^{\mathrm{}}(\mathrm{\Lambda }_b^{0,0})\stackrel{\overline{D}_0}{}C^{\mathrm{}}(\mathrm{\Lambda }_b^{0,1})\stackrel{\overline{D}_1}{}\mathrm{}\stackrel{\overline{D}_{m1}}{}C^{\mathrm{}}(\mathrm{\Lambda }_b^{0,m})0,$$
where $`\mathrm{\Lambda }_b^{0,j}=\mathrm{\Lambda }^j(H^{0,1}X)^{}`$, such that:
1. $`\sigma (\overline{D}_j)=\sigma (\overline{}_b)`$ to second order along $`\mathrm{\Sigma }:=\{(x,r\alpha _x):xX,r>0\}T^{}X`$;
2. The orthogonal projector $`\mathrm{\Pi }:^2(X)^2(X)`$ onto the kernel of $`\overline{D}_0`$ is a complex Fourier integral operator which is microlocally equivalent to the Cauchy-Szegö projector of the holomorphic case;
3. $`(\overline{D}_0,\frac{}{\theta })`$ is jointly elliptic.
The results stated here use only the $`\overline{D}_0`$ term of the complex; its kernel consists of the spaces of almost holomorphic sections of the powers $`L^N`$ of the line bundle $`L`$, as explained below. The complex $`\overline{D}_j`$ was used by Boutet de Monvel -Guillemin \[BoGu, Lemma 14.14\] to show that the dimension of $`H_J^0(M,L^N)`$ or $`_N^2(X)`$ is given by the Riemann-Roch formula (for $`N`$ sufficiently large). For our results, we need only the leading term of Riemann-Roch, which we obtain as a consequence of Theorem 3.1(a). (The reader should be warned that the symbol is described incorrectly in Lemma 14.11 of \[BoGu\]. However, it is correctly described in Theorem 5.9 of the Appendix to \[BoGu\] and also in \[GuUr\]).
We refer to the kernel $`^2(X)=\mathrm{ker}\overline{D}_0^2(X)`$ as the Hardy space of square-integrable ‘almost CR functions’ on $`X`$. The $`^2`$ norm is with respect to the inner product
$$F_1,F_2=\frac{1}{2\pi }_XF_1\overline{F_2}𝑑V_X,F_1,F_2^2(X),$$
(17)
where
$$dV_X=\frac{1}{m!}\alpha (d\alpha )^m=\alpha \pi ^{}dV_M.$$
(18)
The $`S^1`$ action on $`X`$ commutes with $`\overline{D}_0`$; hence $`^2(X)=_{N=0}^{\mathrm{}}_N^2(X)`$ where $`_N^2(X)=\{F^2(X):F(r_\theta x)=e^{iN\theta }F(x)\}`$. We denote by $`H_J^0(M,L^N)`$ the space of sections which corresponds to $`_N^2(X)`$ under the map $`s\widehat{s}`$. Elements of $`H_J^0(M,L^N)`$ are the almost holomorphic sections of $`L^N`$. (Note that products of almost holomorphic sections are not necessarily almost holomorphic.) We henceforth write $`\widehat{s}=s`$ and identify $`H_J^0(M,L^N)`$ with $`_N^2(X)`$. Since $`(\overline{D}_0,\frac{}{\theta })`$ is a jointly elliptic system, elements of $`H_J^0(M,L^N)`$ and $`_N^2(X)`$ are smooth. In many other respects, $`H_J^0(M,L^N)`$ is analogous to the space of holomorphic sections in the complex case. Subsequent results will bear this out.
We let $`\mathrm{\Pi }_N:^2(X)_N^2(X)`$ denote the orthogonal projection. The level $`N`$ Szegö kernel $`\mathrm{\Pi }_N(x,y)`$ is defined by
$$\mathrm{\Pi }_NF(x)=_X\mathrm{\Pi }_N(x,y)F(y)𝑑V_X(y),F^2(X).$$
(19)
It can be given as
$$\mathrm{\Pi }_N(x,y)=\underset{j=1}{\overset{d_N}{}}S_j^N(x)\overline{S_j^N(y)},$$
(20)
where $`S_1^N,\mathrm{},S_{d_N}^N`$ form an orthonormal basis of $`_N^2(X)`$.
### 1.4. Construction of the Szegö kernels
In this section, we will sketch the construction of the operator $`\overline{D}_0`$ of Theorem 1.1 in the special setting of almost complex manifolds, and in so doing we will describe the symbol of the complex in more detail. This will require the introduction of many objects from symplectic geometry and from the microlocal analysis of $`\overline{}_b`$. We will need this material later on in the construction of a parametrix for the Szegö kernel.
#### 1.4.1. The characteristic variety of $`\overline{}_b`$
In general, we denote by $`\sigma _A`$ the principal symbol of a pseudodifferential operator $`A`$. To describe the principal symbol of $`\overline{}_b`$, we introduce convenient local coordinates and frames. Recalling that $`HX=H^{1,0}XH^{0,1}X`$, we again consider local orthonormal frames $`Z_1,\mathrm{},Z_n`$ of $`H^{1,0}X`$, resp. $`\overline{Z}_1,\mathrm{},\overline{Z}_m`$ of $`H^{0,1}X`$, and dual orthonormal coframes $`\vartheta _1,\mathrm{},\vartheta _m,`$ resp. $`\overline{\vartheta }_1,\mathrm{},\overline{\vartheta }_m.`$ Then we have $`\overline{}_b=_{j=1}^m\overline{\vartheta }_j\overline{Z}_j`$. Let us define complex-valued functions on $`T^{}X`$ by:
$$p_j(x,\xi )=Z_j(x),\xi ),\overline{p}_j(x,\xi )=\overline{Z}_j(x),\xi .$$
Then
$$\sigma _{\overline{}_b}(x,\xi )=\underset{j=1}{\overset{m}{}}p_j(x,\xi )ϵ(\overline{\vartheta }_j)$$
where $`ϵ`$ denotes exterior multiplication. We note that $`\{\overline{p}_j,\overline{p}_k\}=[\overline{Z}_j,\overline{Z}_k],\xi `$.
To state results, it is convenient to introduce the operator $`\mathrm{}_b:=\overline{}_b^{}\overline{}_b=_{j=1}^m\overline{Z}_j^{}\overline{Z}_j`$ where $`\overline{Z}_j^{}`$ is the adjoint of the vector field regarded as a linear differential operator. To conform to the notation of \[BoGu\] we also put $`q=\sigma (\mathrm{}_b)=_{j=1}^m|\overline{p}_j|^2.`$ The characteristic variety $`\mathrm{\Sigma }=\{q=0\}`$ of $`\overline{}_b`$ is the same as that of $`\mathrm{}_b`$, namely the vertical sub-bundle of $`T^{}XM.`$ It is the conic submanifold of $`T^{}X`$ parametrized by the graph of the contact form, $`\mathrm{\Sigma }=\{(x,r\alpha _x):r>0\}X\times ^+`$. It follows that $`\mathrm{\Sigma }`$ is a symplectic submanifold. It is the dual (real) line bundle to the vertical subbundle $`VTX`$, since $`\alpha (X)=G(X,\frac{}{\theta }).`$
#### 1.4.2. The positive Lagrangean ideal $`I`$
To construct the $`\overline{D}_j`$-complex replacing the $`\overline{}_b`$-complex in the non-integrable case, and to construct the Szegö kernel, we will need to study a positive Lagrangean ideal $`I`$ whose generators will define the principal symbol of $`\overline{D}_0`$. For background on positive Lagrangean ideals, see \[Hör\].
###### Proposition 1.2.
There exists a unique positive Lagrangean ideal $`I`$ with respect to $`\mathrm{\Sigma }`$ containing $`q`$. That is, there exists a unique ideal $`II_\mathrm{\Sigma }`$ (where $`I_\mathrm{\Sigma }`$ is the ideal of functions vanishing on $`\mathrm{\Sigma }`$) satisfying:
* $`I`$ is closed under Poisson bracket;
* $`\mathrm{\Sigma }`$ is the set of common zeros of $`fI`$;
* There exist local generators $`\zeta _1,\mathrm{},\zeta _m`$ such that the matrix $`\left(\frac{1}{i}\{\zeta _j,\overline{\zeta }_k\}\right)`$ is positive definite on $`\mathrm{\Sigma }`$ and that $`q=_{j,k}\lambda _{j\overline{k}}\zeta _j\overline{\zeta }_k`$, where $`\{\lambda _{j\overline{k}}\}`$ is a hermitian positive definite matrix of functions.
###### Proof.
In the holomorphic case, $`I`$ is generated by the linear functions $`\zeta _j(x,\xi )=\xi ,\overline{Z}_j`$. In the general almost complex (or rather almost CR) setting, these functions do not Poisson commute and have to be modified. Since the deviation of an almost complex structure from being integrable (i.e. a true complex structure) is measured by the Nijenhuis bracket, it is not surprising that the generators $`\zeta _j`$ can be constructed from the linear functions $`\xi ,\overline{Z}_j`$ and from the Nijenhuis tensor. We now explain how to do this, basically following the method of \[BoGu\].
As a first approximation to the $`\zeta _j`$ we begin with the linear functions $`\zeta _j^{(1)}=\overline{p}_j`$ on $`T^{}X`$. As mentioned above, the $`\zeta _j^{(1)}`$ do not generate a Lagrangean ideal in the non-integrable almost complex case, indeed
$$\{\zeta _j^{(1)},\zeta _k^{(1)}\}=\xi ,[\overline{Z}_j(x),\overline{Z}_k(x)].$$
(21)
However we do have that
$$\{\zeta _j^{(1)},\zeta _k^{(1)}\}=\{\xi ,\overline{Z}_j(x),\xi ,\overline{Z}_k(x)\}=0\text{on}\mathrm{\Sigma }.$$
Indeed, for $`(x,\xi )\mathrm{\Sigma }`$, we have $`\xi =r\alpha _x`$ for some $`r>0`$ so that
$$\begin{array}{c}\{\xi ,\overline{Z}_j(x),\xi ,\overline{Z}_k(x)\}=r\alpha _x([\overline{Z}_j(x),\overline{Z}_k(x)])\hfill \\ =rd\alpha _x(\overline{Z}_j(x),\overline{Z}_k(x))=r\pi ^{}\omega (\overline{Z}_j(x),\overline{Z}_k(x))=0\hfill \end{array}$$
(22)
since $`\{\overline{Z}_j\}`$ forms a Lagrangean subspace for the horizontal symplectic form $`\pi ^{}\omega `$. Here, $`\pi :XM`$ is the natural projection. Moreover if we choose the local horizontal vector fields $`Z_j`$ to be orthonormal relative to $`\pi ^{}\omega `$, then we also have:
$$\begin{array}{c}\{\zeta _j^{(1)},\overline{\zeta }_k^{(1)}\}(x,\xi )=\xi ,[\overline{Z}_j(x),Z_k(x)]=r\pi ^{}\omega (\overline{Z}_j(x),Z_k(x))\hfill \\ =ir\delta _j^k=i\delta _j^kp_\theta (x,\xi ),(x,\xi )\mathrm{\Sigma }.\hfill \end{array}$$
(23)
Here, $`p_\theta (x,\xi )=\xi ,\frac{}{\theta }`$.
Finally, we have
$$q=\underset{j=1}{\overset{m}{}}|\xi ,Z_j|^2=\underset{j=1}{\overset{m}{}}|\zeta _j^{(1)}|^2.$$
Hence the second and third conditions on the $`\zeta _j`$ are satisfied by the functions $`\zeta _j^{(1)}`$. Furthermore, equation (21) tells us that the first condition is satisfied to zero-th order for the ideal $`I_1=(\zeta _1^{(1)},\mathrm{},\zeta _m^{(1)})`$. In fact, let us precisely describe the error. We consider the orthonormal (relative to $`\omega `$) vector fields $`Z_j^M=\pi _{}Z_j`$ of type (1,0) on $`M`$. Recall that the Nijenhuis tensor is given by
$$N(V,W)=\frac{1}{2}\left([JV,JW][V,W]J[V,JW]J[JV,W]\right).$$
Hence,
$$N(Z_j^M,Z_k^M)=(1iJ)[Z_j^M,Z_k^M]=2[Z_j^M,Z_k^M]_{(0,1)}\stackrel{\mathrm{def}}{=}\underset{p=1}{\overset{m}{}}N_{jk}^p\overline{Z}_p^M.$$
(24)
We note that by definition,
$$N_{jk}^p=N_{kj}^p.$$
(25)
Furthermore, by the Jacobi identity
$$\{\{\zeta _j,\zeta _k\},\zeta _p\}+\{\{\zeta _p,\zeta _j\},\zeta _k\}+\{\{\zeta _k,\zeta _p\},\zeta _j\}=0$$
applied to $`(x,\alpha _x)\mathrm{\Sigma }`$, we have
$$N_{jk}^p+N_{pj}^k+N_{kp}^j=0.$$
(26)
By (22) and (24), we have
$$\{\zeta _j^{(1)},\zeta _k^{(1)}\}=\underset{p=1}{\overset{m}{}}f_p^1\zeta _p^{(1)}+\underset{p=1}{\overset{m}{}}\overline{N}_{jk}^p\overline{\zeta }_p^{(1)}.$$
(27)
We now argue, following \[BoGu\], that these functions can be successively modified to satisfy the same conditions to infinite order on $`\mathrm{\Sigma }.`$ The next step is to modify the functions $`\zeta _j^{(1)}`$ by quadratic terms so that they satisfy the conditions $`\{\zeta _j,\zeta _k\}I`$ to first order and the condition $`q=_j|\zeta _j|^2`$ to order 3 on $`\mathrm{\Sigma }`$. So we try to construct new functions
$$\zeta _p^{(2)}=\zeta _p^{(1)}+R_p,R_p=\underset{j,k}{}\nu _p^{jk}\overline{\zeta }_j^{(1)}\overline{\zeta }_k^{(1)}$$
so that
$`\{\zeta _j^{(2)},\zeta _k^{(2)}\}`$ $`=`$ $`{\displaystyle \underset{p}{}}f_p^2\zeta _p^{(2)}+{\displaystyle \underset{\alpha _1,\alpha _2}{}}\mu _{jk}^{\alpha _1\alpha _2}\overline{\zeta }_{\alpha _1}^{(2)}\overline{\zeta }_{\alpha _2}^{(2)};`$ (28)
$`q`$ $`=`$ $`{\displaystyle \underset{p}{}}v_p^2\zeta _p^{(2)}+{\displaystyle \underset{\alpha }{}}\phi _p^\alpha \overline{\zeta }_{\alpha _1}^{(2)}\overline{\zeta }_{\alpha _2}^{(2)}\overline{\zeta }_{\alpha _3}^{(2)}\overline{\zeta }_{\alpha _4}^{(2)},(\alpha =(\alpha _1,\mathrm{},\alpha _4)).`$ (29)
Let us now solve (28)–(29) for the $`\nu _p^{jk}`$. First of all, we choose $`\nu _p^{jk}=\nu _p^{kj}`$. We have
$$\{\zeta _j^{(2)},\zeta _k^{(2)}\}=\underset{p=1}{\overset{m}{}}f_p^1\zeta _p^{(1)}+\underset{p=1}{\overset{m}{}}\overline{N}_{jk}^p\overline{\zeta }^{(1)}+\{\zeta _j^{(1)},R_k\}\{\zeta _k^{(1)},R_j\}modI_\mathrm{\Sigma }^2.$$
By (23), we have
$$\{\zeta _j^{(1)},\overline{\zeta }_k^{(1)}\}=i\delta _j^kp_\theta modI_\mathrm{\Sigma },$$
(30)
and thus
$$\{\zeta _j^{(1)},R_k\}=\underset{p=1}{\overset{m}{}}2i\nu _k^{pj}p_\theta \overline{\zeta }_p^{(1)}modI_\mathrm{\Sigma }^2.$$
Therefore,
$$\{\zeta _j^{(2)},\zeta _k^{(2)}\}=\underset{p=1}{\overset{m}{}}f_p^1\zeta _p^{(2)}+\underset{p=1}{\overset{m}{}}\left(\overline{N}_{jk}^p+2i(\nu _k^{pj}\nu _j^{pk})p_\theta \right)\overline{\zeta }_p^{(1)}modI_\mathrm{\Sigma }^2.$$
(31)
Hence
$$\overline{N}_{jk}^p=2i(\nu _j^{pk}\nu _k^{pj})p_\theta \text{on}\mathrm{\Sigma },$$
or equivalently,
$$\nu _j^{pk}\nu _k^{pj}=\frac{i}{2p_\theta }\overline{N}_{jk}^pmodI_\mathrm{\Sigma }.$$
(32)
On the other hand,
$`q`$ $`=`$ $`{\displaystyle \underset{p}{}}|\zeta _p^{(2)}R_p|^2={\displaystyle \underset{p}{}}v_p^2\zeta _p^{(2)}R_p\overline{\zeta }_p^{(2)}`$
$`=`$ $`{\displaystyle \underset{p}{}}v_p^2\zeta _p^{(2)}{\displaystyle \underset{j,k,p}{}}\nu _p^{jk}\overline{\zeta }_j^{(2)}\overline{\zeta }_k^{(2)}\overline{\zeta }_p^{(2)}+{\displaystyle \underset{\alpha }{}}\phi _p^\alpha \overline{\zeta }_{\alpha _1}^{(2)}\overline{\zeta }_{\alpha _2}^{(2)}\overline{\zeta }_{\alpha _3}^{(2)}\overline{\zeta }_{\alpha _4}^{(2)}.`$
Hence (29) is equivalent to
$$\nu _p^{jk}+\nu _k^{pj}+\nu _j^{kp}=0.$$
(33)
Using (25)–(26), we can solve the equations (30) and (32) to obtain
$$\nu _p^{jk}=\frac{i}{6p_\theta }\left(\overline{N}_{pj}^k+\overline{N}_{pk}^j\right).$$
(34)
Indeed, the solution (34) is unique (modulo $`I_\mathrm{\Sigma }`$) and hence the $`R_p`$ are unique modulo $`I_\mathrm{\Sigma }^3`$. In summary,
$$\zeta _p^{(2)}=\zeta _p^{(1)}+\frac{i}{3p_\theta }\underset{j,k}{}\overline{N}_{pj}^k\overline{\zeta }_j^{(1)}\overline{\zeta }_k^{(1)}.$$
(35)
The passage from the $`n^{\mathrm{th}}`$ to the $`(n+1)^{\mathrm{st}}`$ step is similar, and we refer to \[BoGu, pp. 147–149\]. ∎
Remark: Define $`p_\theta (x,\xi )=\xi ,\frac{}{\theta }.`$ Since the joint zero set of $`\{\zeta _1,\mathrm{},\zeta _m\}`$ equals $`\mathrm{\Sigma }`$ and since $`p_\theta 0`$ on $`\mathrm{\Sigma }0`$ it follows that $`\{\zeta _1,\mathrm{},\zeta _m,p_\theta \}`$ is an elliptic system of symbols.
#### 1.4.3. The complex canonical relation
Our eventual goal is to prove that $`\mathrm{\Pi }`$ is a complex Fourier integral operator and to construct a parametrix for it. As a preliminary step we need to construct and describe the complex canonical relation $`C`$ underlying $`\mathrm{\Pi }`$. As is typical with complex Fourier integral operators, $`C`$ does not live in $`T^{}X\times T^{}X`$ but rather in its almost analytic extension $`T^{}\stackrel{~}{X}\times T^{}\stackrel{~}{X}`$. Here, $`\stackrel{~}{N}`$ denotes the almost analytic extension of a $`C^{\mathrm{}}`$ manifold $`N`$. Although the language of almost analytic extensions may seem heavy, it is very helpful if one wishes to understand the full (complex) geometry of $`C`$. When $`N`$ is real analytic, $`\stackrel{~}{N}`$ is the usual complexification of $`X`$, i.e. a complex manifold in which $`N`$ sits as a totally real submanifold. The reader may find it simpler to make this extra assumption. For background on almost analytic extensions, we refer to \[MelSj, MenSj\].
Since $`\pi :XM`$ is an $`S^1`$ bundle over $`M`$, its complexification $`\stackrel{~}{\pi }:\stackrel{~}{X}\stackrel{~}{M}`$ defines a $`^{}`$ bundle over $`\stackrel{~}{M}`$. The connection form $`\alpha `$ has an (almost) analytic continuation to a connection $`\stackrel{~}{\alpha }`$ to this bundle and we may split $`T\stackrel{~}{X}=\stackrel{~}{H}\stackrel{~}{V}`$, where $`\stackrel{~}{V}T\stackrel{~}{M}`$ is the vertical subbundle of the fibration $`\stackrel{~}{X}\stackrel{~}{M}`$ and where $`\stackrel{~}{H}T\stackrel{~}{M}`$ is the kernel of $`\stackrel{~}{\alpha }.`$
The (almost) complexification of $`T^{}X`$ is of course $`T^{}(\stackrel{~}{X})`$. We denote the canonical symplectic form on $`T^{}X`$ by $`\sigma `$ and that on $`T^{}(\stackrel{~}{X}`$ by $`\stackrel{~}{\sigma }`$; the notation is consistent because it is the complexification of $`\sigma .`$ The symplectic cone $`\mathrm{\Sigma }`$ complexifies to $`\stackrel{~}{\mathrm{\Sigma }}`$ and it remains symplectic with respect to $`\stackrel{~}{\sigma }.`$ It is given by $`\{(\stackrel{~}{x},\stackrel{~}{\lambda }\stackrel{~}{\alpha }_{\stackrel{~}{x}}):\stackrel{~}{\lambda }^{}\}.`$ We have a natural identification $`L^{}\mathrm{\Sigma }`$ given by $`rx(x,r\alpha _x).`$ We further note that the $`^{}`$ bundle $`L^{}M`$ is the fiberwise complexification of the $`S^1`$ bundle $`XM`$, hence $`L^{}M`$ is the restriction of $`\stackrel{~}{\pi }`$ to $`\stackrel{~}{\pi }^1(M).`$ We will therefore view $`L^{}`$ as a submanifold of $`\stackrel{~}{X}.`$
#### 1.4.4. Definition of $`C`$
Let $`\stackrel{~}{\zeta }_j`$ be the almost analytic extensions of the functions $`\zeta _j`$. Then put
$$𝒥_+=\{(\stackrel{~}{x},\stackrel{~}{\xi })T^{}\stackrel{~}{X}:\stackrel{~}{\zeta }_j=0j\}.$$
(36)
It is an involutive manifold of $`T^{}(\stackrel{~}{X})`$ with the properties:
$$\begin{array}{cc}(i)\hfill & (𝒥_+)_{}=\mathrm{\Sigma }\hfill \\ & \\ (ii)\hfill & q|_{𝒥_+}0\hfill \\ & \\ (iii)\hfill & \frac{1}{i}\sigma (u,\overline{u})>0,uT(𝒥_+)^{}\hfill \\ & \\ (iv)\hfill & T_\rho (𝒥_+)=T_\rho \stackrel{~}{\mathrm{\Sigma }}\mathrm{\Lambda }_\rho ^+.\hfill \end{array}$$
(37)
Here, $`\mathrm{\Lambda }_\rho ^\pm `$ is the sum of the eigenspaces of $`F_\rho `$, the normal Hessian of $`q`$, corresponding to the eigenvalues $`\{\pm i\lambda _j\}.`$ The null foliation of $`𝒥_+`$ is given by the joint Hamilton flow of the $`\stackrel{~}{\zeta }_j`$’s.
The following proposition, proved in \[MenSj\] and in (\[BoGu\]),Appendix, Lemma 4.5) defines the complex canonical relation $`C`$:
###### Proposition 1.3.
There exists a unique strictly positive almost analytic canonical relation $`C`$ satisfying
$$\mathrm{diag}(\mathrm{\Sigma })C𝒥_+\times \overline{𝒥_+}.$$
(38)
Indeed,
$$C=\{(\stackrel{~}{x},\stackrel{~}{\xi },\stackrel{~}{y},\stackrel{~}{\eta })𝒥_+\times \overline{𝒥_+}:(\stackrel{~}{x},\stackrel{~}{\xi })(\stackrel{~}{y},\stackrel{~}{\eta })\},$$
(39)
where $``$ is the equivalence relation of ‘belonging to the same leaf of the null foliation of $`𝒥_+.`$ Thus, $`C`$ is the flow-out of its real points, $`\mathrm{diag}(\mathrm{\Sigma })`$, under the joint Hamilton flow of the $`\stackrel{~}{\zeta }_j`$’s. It is clear from the description that $`CC=C^{}=C,`$ i.e. that $`C`$ is an idempotent canonical relation. It follows that $`I^{}(X\times X,C)`$ is a $``$-algebra.
#### 1.4.5. Definition of the Szegö projector
Having constructed $`C`$, we define a Szegö projector $`\mathrm{\Pi }`$ associated to $`\mathrm{\Sigma }`$ and $`C`$ to be a self-adjoint projection $`\mathrm{\Pi }`$ in the Fourier integral operator class $`I^{}(X\times X,C)`$ with principal symbol $`1`$ (relative to the canonical 1/2-density of $`C`$).
It is simple to prove the existence of such a projection (see \[BoGu\], Appendix A.4): Since $`I^{}(X\times X,C)`$ is a $``$-algebra, there exists an element $`AI^{}(X\times X,C)`$ with $`\sigma _A=1`$ on $`\mathrm{diag}(\mathrm{\Sigma })`$, or more precisely with $`\sigma _A`$ equal to a projection onto a prescribed ’vacuum state’. The principal symbols of $`A^2A`$ and $`AA^{}`$ then vanish, so these operators are of negative order. It follows that the spectrum of $`A`$ is concentrated near $`\{0,1\}`$. Hence there exists an analytic function in a neighborhood of the spectrum such that $`F(A):=\mathrm{\Pi }`$ is a true projection. Since $`I^{}(X\times X,C)`$ is closed under functional calculus, this projection lies in that algebra.
We note that $`\mathrm{\Pi }`$ is far from unique; given any $`\mathrm{\Pi }`$ one could set $`\mathrm{\Pi }^{}=e^{iA}\mathrm{\Pi }e^{iA}`$ where $`A`$ is a pseudodifferential operator of order $`1`$. We just fix one choice in what follows.
Remark: In \[BoGu\] the term Szegö projector (or Toeplitz structure) is used for a projection operator with wave front set on $`\mathrm{\Sigma }`$ which is microlocally equivalent to the following model case on $`^{2m+2}\times ^{2m}`$ (\[Bout, Sec. 5\], \[BoGu\]).
Let us use coordinates $`y^{2m+2},t^{2m}`$, let $`\eta ,\tau `$ be the symplectically dual coordinates and consider the operators
$$A_j:=D_{y_j}+iy_j|D_t|,\overline{A}_j=D_{y_j}iy_j|D_t|.$$
Here, $`D_x=\frac{}{ix}`$ and $`|D_t|`$ is Fourier multiplication by $`|\tau |.`$ The operators $`A_j`$, resp. $`\overline{A}_j`$, are what are familiarly known as creation operators, resp. annihilation, operators in the representation theory of the Heisenberg group. The characteristic variety of the system $`\{\overline{A}_j\}`$ is given by $`\mathrm{\Sigma }^0=\{t=\tau =0\}^{2m+2}.`$ The Hardy space is given by $`^2=\{f:\overline{A}_jf=0,j\}`$ and the Szegö kernel is given by the complex Fourier integral kernel
$$\mathrm{\Pi }^0(t,y,t^{},y^{})=C_m_^me^{i\mathrm{\Phi }}|\tau |^m𝑑\tau ,\mathrm{\Phi }=tt^{},\tau +i|\tau |(|y|^2+|y^{}|^2).$$
The positive Lagrangean ideal $`I`$ is generated by the symbols $`\sigma (A_j)=\zeta _j=\eta _j+i|\tau |y_j`$.
#### 1.4.6. Construction of the complex
Having defined $`\mathrm{\Pi }`$, one first constructs $`\overline{D}_0`$ so that $`\overline{D}_0\mathrm{\Pi }=0`$. In terms of a local frame $`\overline{\vartheta }_j`$ of horizontal (0,1)-forms on $`X`$, we may write
$$\overline{D}_0f=\underset{j=1}{\overset{m}{}}\widehat{\zeta }_j(x,D)f\overline{\vartheta }_j.$$
(40)
The coefficient operators $`\widehat{\zeta }_j(x,D)`$ are first order pseudodifferential operators with principal symbols equal to $`\zeta _j`$ and satisfying
$$\widehat{\zeta }_j\mathrm{\Pi }0$$
modulo smoothing operators. That is, one ‘quantizes’ the $`\zeta _j`$’s as first order pseudodifferential operators which annihilate $`\mathrm{\Pi }.`$ Let us briefly summarize their construction (following \[BoGu, Appendix\]).
We begin with any $`S^1`$-equivariant symmetric first order pseudodifferential operator $`\overline{D}_0^{}`$ with principal symbol equal to $`_{j=1}^m\zeta _j\overline{\vartheta }_j.`$ Then $`\overline{D}_0^{}\mathrm{\Pi }`$ is of order $`0`$ so one may find a zeroth order pseudodifferential operator $`Q_0`$ such that $`\overline{D}_0^{}\mathrm{\Pi }Q_0\mathrm{\Pi }`$ (modulo smoothing operators). Then put: $`\overline{D}_0=(\overline{D}_0^{}Q_0)(\overline{D}_0^{}Q_0)\mathrm{\Pi }`$. Clearly, $`\overline{D}_0\mathrm{\Pi }=0`$ and $`\sigma (\overline{D}_0)=\sigma (\overline{D}_0^{})=_{j=1}^m\zeta _j\overline{\vartheta }_j.`$ The characteristic variety of $`\overline{D}_0`$ is then equal to $`\mathrm{\Sigma }`$. Since $`p_\theta `$ is the symbol of $`\frac{}{\theta }`$ and since the system $`\{\sigma _{\overline{D}_0},p_\theta \}`$ has no zeros in $`T^{}X0`$ it follows that $`\{\overline{D}_0,\frac{}{\theta }\}`$ is an elliptic system.
Remark: One can then construct the higher $`\overline{D}_j`$ recursively so that $`\overline{D}_j\overline{D}_{j1}=0`$. We refer to \[BoGu\], Appendix §5, for further details.
## 2. Parametrix for the Szegö projector
In \[BSZ2, Theorem 3.1\], we showed that for the complex case, the scaled Szegö kernel $`\mathrm{\Pi }_N`$ near the diagonal is asymptotic to the Szegö kernel $`\mathrm{\Pi }_1^𝐇`$ of level one for the reduced Heisenberg group, given by
$$\mathrm{\Pi }_1^𝐇(z,\theta ;w,\phi )=\frac{1}{\pi ^m}e^{i(\theta \phi )+i\mathrm{}(z\overline{w})\frac{1}{2}|zw|^2}=\frac{1}{\pi ^m}e^{i(\theta \phi )+z\overline{w}\frac{1}{2}(|z|^2+|w|^2)}.$$
(41)
The method was to apply the Boutet de Monvel-Sjöstrand oscillatory integral formula
$$\mathrm{\Pi }(x,y)=_0^{\mathrm{}}e^{it\psi (x,y)}s(x,y,t)𝑑t$$
(42)
arising from a parametrix construction (\[BoSj, Th. 1.5 and §2.c\]). Let us recall the construction of $`\psi (x,y)`$ in the integrable complex case. Fix a local holomorphic section $`e_L`$ of $`L`$ over $`UM`$ and define $`aC^{\mathrm{}}(U)`$ by $`a=|e_L|_h^2`$. Since $`L^{}|_UU\times `$ we can define local coordinates on $`L^{}`$ by $`(z,\lambda )\lambda e_L(z)`$. Then a defining function of $`XL^{}`$ is given by $`\rho (z,\lambda )=1|\lambda |^2a(z)`$. Define the function $`a(z,w)`$ as the almost analytic extension of $`a(z)`$, i.e. the solution of $`\overline{}_za=0=_wa,a(z,z)=a(z)`$ and put $`\psi (x,y)=i(1\lambda \overline{\mu }a(z,w)).`$ Then $`t\psi `$ is a phase for $`\mathrm{\Pi }.`$
The object of this section is to show that the universal asymptotic formula of \[BSZ2\] for the near-diagonal scaled Szegö kernel holds for the symplectic case (Theorem 2.3). To do this, we first show that the Boutet de Monvel-Sjöstrand construction can be extended to the symplectic almost-complex case. Indeed we will obtain (Theorem 2.1) an integral formula of the form (42) for the symplectic case. In fact, our local phase function $`\psi `$ will be shown to be of the form $`\psi (x,y)=i(1\lambda \overline{\mu }a(z,w))`$, where $`\overline{a(w,z)}=a(z,w)`$ and hence $`\psi (y,x)=\overline{\psi (x,y)}`$.
### 2.1. Oscillatory integral for $`\mathrm{\Pi }`$
In order to obtain our integral formula, we first recall the notion of parametrizing an almost analytic Lagrangean $`\mathrm{\Lambda }`$ by a phase function. We assume $`\phi (x,\theta )`$ is a regular phase function in the sense of (\[MelSj, Def. 3.5\]), i.e. that it has no critical points, is homogeneous of degree one in $`\theta `$, that the differentials $`d\frac{\phi }{\theta _j}`$ are linearly independent over $``$ on the set
$$C_\phi =\{(x,\theta ):d_\theta \phi =0\}$$
and such that $`\mathrm{}\phi 0`$. We then let $`\stackrel{~}{\phi }(\stackrel{~}{x},\stackrel{~}{\theta })`$ be an almost analytic extension, put
$$C_{\stackrel{~}{\phi }}=\{(\stackrel{~}{x},\stackrel{~}{\theta }):d_{\stackrel{~}{\theta }}\stackrel{~}{\phi }=0\}$$
and define the Lagrange immersion
$$\iota _{\stackrel{~}{\phi }}:(\stackrel{~}{x},\stackrel{~}{\theta })C_{\stackrel{~}{\phi }}(\stackrel{~}{x},d_{\stackrel{~}{x}}\stackrel{~}{\phi }(\stackrel{~}{x},\stackrel{~}{\theta })).$$
The phase $`\phi `$ parametrizes $`\mathrm{\Lambda }`$ if $`\mathrm{\Lambda }`$ is the image of this map.
The parametrix is an explicit construction of $`\mathrm{\Pi }(x,y)`$ as a complex Lagrangean kernel. What we wish to prove now is that $`C`$ can be parametrized, exactly as in the CR case, by a phase $`\lambda \psi (x,y)`$ defined on $`^+\times X\times X.`$ This is helpful in analyzing the scaling limit of $`\mathrm{\Pi }_N(x,y)`$. In the following we use local coordinates $`(z,\lambda )`$ on $`L^{}`$ coming from a choice of local coordinates $`z`$ on $`M`$ and a local frame $`e_L(z)`$ of $`L`$, and a corresponding local trivialization $`(\stackrel{~}{z},\lambda )`$ of $`\stackrel{~}{X}\stackrel{~}{M}`$. As before, we let $`a=e_L^{}^2`$.
###### Theorem 2.1.
Let $`\mathrm{\Pi }(x,y):^2(X)^2(X)`$ be the Szegö kernel. Then there exists a unique regular phase function $`it\psi (x,y)C^{\mathrm{}}(^+\times X\times X)`$ of positive type and a symbol $`sS^m(X\times X\times ^+)`$ of the type
$$s(x,y,t)\underset{k=0}{\overset{\mathrm{}}{}}t^{mk}s_k(x,y)$$
such that $`id_x\psi |_{x=y}=id_y\psi |_{x=y}=\alpha `$ and
$$\mathrm{\Pi }(x,y)=_0^{\mathrm{}}e^{it\psi (x,y)}s(x,y,t)𝑑t.$$
Furthermore, the almost analytic extension $`\stackrel{~}{\psi }C^{\mathrm{}}(\stackrel{~}{X}\times \stackrel{~}{X})`$ of $`\psi `$ has the form $`\stackrel{~}{\psi }(\stackrel{~}{x},\stackrel{~}{y})=i(1\lambda \overline{\mu }\stackrel{~}{a}(\stackrel{~}{z},\stackrel{~}{w}))`$ with $`\stackrel{~}{a}(z,z)=a(z)`$ and $`\stackrel{~}{a}(\stackrel{~}{z},\stackrel{~}{w})=\overline{\stackrel{~}{a}(\stackrel{~}{w},\stackrel{~}{z})}`$.
###### Proof.
We need to construct a function $`a(z,w)`$ so that $`it\psi `$ as above parametrizes the canonical relation $`C`$, i.e that $`C`$ is the image of the Lagrange immersion
$$\begin{array}{c}\iota _{\stackrel{~}{\psi }}:C_{t\stackrel{~}{\psi }}=^+\times \{\stackrel{~}{\psi }=0\}T^{}(\stackrel{~}{X}\times \stackrel{~}{X})\hfill \\ \\ (t,\stackrel{~}{x},\stackrel{~}{y})(\stackrel{~}{x},td_{\stackrel{~}{x}}\stackrel{~}{\psi };\stackrel{~}{y},td_{\stackrel{~}{y}}\stackrel{~}{\psi })\hfill \end{array}$$
(43)
Since $`C`$ is the unique canonical relation satisfying $`\mathrm{diag}(\mathrm{\Sigma })C𝒥_+\times \overline{𝒥_+}`$, the conditions that $`\stackrel{~}{\psi }`$ parametrize $`C`$ are the following:
1. $`\{(x,y)X\times X:\psi (x,y)=0\}=\mathrm{diag}(X)`$;
2. $`d_x\psi |_{x=y}=d_y\psi |_{x=y}=r\alpha `$ for $`x,yX`$ and for some function $`r(x)>0`$;
3. $`\stackrel{~}{\zeta }_j(\stackrel{~}{x},d_{\stackrel{~}{x}}\stackrel{~}{\psi })=0=\stackrel{~}{\zeta }_j(\stackrel{~}{y},d_{\stackrel{~}{y}}\stackrel{~}{\psi })`$ on $`\{\stackrel{~}{\psi }=0\}.`$
Such a $`\stackrel{~}{\psi }`$ is not unique, so we require that $`r1`$ in condition (ii), i.e.,
$$d_x\psi |_{x=y}=d_y\psi |_{x=y}=\alpha .$$
Suppose we have $`\stackrel{~}{\psi }(\stackrel{~}{x},\stackrel{~}{y})=i(1\lambda \overline{\mu }\stackrel{~}{a}(\stackrel{~}{z},\stackrel{~}{w}))`$. We observe that
$$\stackrel{~}{\psi }=0\stackrel{~}{a}(\stackrel{~}{z},\stackrel{~}{w})=(\lambda \overline{\mu })^1,$$
and hence
$$\begin{array}{c}id_{\stackrel{~}{x}}\stackrel{~}{\psi }=\overline{\mu }\stackrel{~}{a}(\stackrel{~}{z},\stackrel{~}{w})d\lambda +\lambda \overline{\mu }d_{\stackrel{~}{z}}\pi ^{}\stackrel{~}{a}(\stackrel{~}{z},\stackrel{~}{w})\hfill \\ \\ =\lambda ^1d\lambda +\stackrel{~}{a}^1d_{\stackrel{~}{z}}\stackrel{~}{a}(\stackrel{~}{z},\stackrel{~}{w})\stackrel{~}{\psi }=0.\hfill \end{array}$$
(44)
The conditions on $`a`$ are therefore:
$$\{\begin{array}{cc}a(z,w)\lambda \overline{\mu }=1(z,\lambda )=(w,\mu )X;\hfill & \\ & \\ (a^1d_za+\lambda ^1d\lambda )|_{\mathrm{diag}(X)}=(a^1d_wa+\lambda ^1d\lambda )|_{\mathrm{diag}(X)}=\alpha \hfill & \\ & \\ \stackrel{~}{\zeta }_j(\stackrel{~}{z},\lambda ,\lambda ^1d\lambda +\stackrel{~}{a}^1d_{\stackrel{~}{z}}\stackrel{~}{a}(\stackrel{~}{z},\stackrel{~}{w}))=0=\stackrel{~}{\zeta }_j(\stackrel{~}{w},\stackrel{~}{\mu },\mu ^1d\mu +\stackrel{~}{a}^1d_{\stackrel{~}{w}}\stackrel{~}{a}(\stackrel{~}{z},\stackrel{~}{w})),(z,w,\lambda ,\mu )\hfill & \end{array}$$
A solution $`a(z,w)`$ satisfying the first condition must satisfy $`a(z,z)|\lambda |^2=1`$ on $`X`$, so that $`a(z,z)|\lambda |^2`$ is the local hermitian metric on $`L^{}`$ with unit bundle $`X`$, i.e. $`a(z,z)=a(z)`$.
We now prove that these conditions have a unique solution near the diagonal. We do this by reducing the canonical relation $`C`$ by the natural $`S^1`$ symmetry. The reduced relation $`C_r`$ has a unique generating function $`\mathrm{log}a`$; the three conditions above on $`a`$ will follow automatically from this fact.
The $`S^1`$ action of $`X`$ lifts to $`T^{}X`$ as the Hamiltonian flow of the function $`p_\theta (x,\xi ):=\xi ,\frac{}{\theta }.`$ The $`\zeta _j`$ are invariant under this $`S^1`$ action, hence
$$\{p_\theta ,\zeta _j\}=0j.$$
(45)
Now consider the level set $`\{p_\theta =1\}T^{}X`$. Dual to the splitting $`TX=HV`$ we get a splitting $`T^{}X=H^{}V^{}`$, where
$$V^{}(X)=\alpha =H^o,H^{}(X)=V^o$$
where $`E^o`$ denotes the annihilator of a subspace $`E`$, i.e. the linear functionals which vanish on $`E`$. Thus, $`p_\theta =0`$ on the horizontal space $`H^{}(X)`$ and $`p_\theta (\alpha )=1`$. Since $`p_\theta `$ is linear on the fibers of $`T^{}X`$, the set $`\{p_\theta =1\}`$ has the form $`\{\alpha +h:hH^{}(X)\}`$. We also note that $`p_\theta (d\theta )=1`$ in the local coordinates $`(z,\theta )`$ on $`X`$ defined by $`\lambda =e^{i\theta }`$. Hence $`\{p_\theta =1\}`$ may also be identified with $`\{d\theta +h:hH^{}(X)\}`$.
Since $`\{p_\theta =1\}`$ is a hypersurface, its null-foliation is given by the orbits of the Hamiltonian flow of $`p_\theta `$, i.e. by the $`S^1`$ action. We use the term ‘reducing by the $`S^1`$-action’ to mean setting $`p_\theta =1`$ and then dividing by this action. The reduction of $`T^{}X`$ is thus defined by $`(T^{}X)_r=p_\theta ^1(1)/S^1`$. Since $`p_\theta ^1(1)`$ is an affine bundle over $`X`$ with fiber isomorphic to $`H^{}(X)T^{}M,`$ it is clear that $`(T^{}X)_rT^{}M`$ as vector bundles over $`M`$. We can obtain a symplectic equivalence using the local coordinates $`(z,\theta )`$ on $`X`$. Let $`(p_z,p_\theta )`$ be the corresponding symplectically dual coordinates, so that the natural symplectic form $`\sigma _{T^{}X}`$ on $`T^{}X`$ is given by $`\sigma _{T^{}X}=dzdp_z+d\theta dp_\theta `$. The notation $`p_\theta `$ is consistent with the above. Moreover, the natural symplectic form on $`T^{}M`$ is given locally by $`\sigma _{T^{}M}=dzdp_z.`$ Now define the projection
$$\chi :p_\theta ^1(1)T^{}M,\chi (z,p_z,1,p_\theta )=(z,p_z).$$
This map commutes with the $`S^1`$ action and hence descends to the quotient to define a local map over $`U`$, still denoted $`\chi `$, from $`(T^{}X)_rT^{}M`$. Clearly $`\chi `$ is symplectic.
We now reduce the canonical relation $`C`$. Thus we consider the $`^{}\times ^{}`$ action on $`T^{}\stackrel{~}{X}\times T^{}\stackrel{~}{X}0`$ generated by $`p_\theta (x,\xi ),p_\theta (y,\eta ).`$ The reduction of $`C`$ is given by
$$C_r=C(p_\theta \times p_\theta )^1(1,1)/^{}\times ^{}.$$
We then use $`\chi \times \chi `$ to identify $`C_r`$ with a (non-homogeneous) positive canonical relation in $`T^{}(\stackrel{~}{M}\times \stackrel{~}{M}).`$ Thus in coordinates,
$$C_r=\{(\stackrel{~}{z},\stackrel{~}{p_z},\stackrel{~}{w},\stackrel{~}{p_w})T^{}(\stackrel{~}{M}\times \stackrel{~}{M}):\lambda ,\mu ,(\stackrel{~}{z},\lambda ,\stackrel{~}{p_z},1;\stackrel{~}{w},\mu ,\stackrel{~}{p_w},1)C\}.$$
(46)
Since reduction preserves real points, it is clear that
$$\begin{array}{c}(C_r)_{}=C_{}(p_\theta \times p_\theta )^1(1,1)/^{}\times ^{}\hfill \\ \\ =\{(z,p_z,z,p_z)\mathrm{diag}(T^{}(M\times M)):\theta \text{such that}\alpha _{z,e^{i\theta }}=d\theta +p_z\}.\hfill \end{array}$$
Let us denote by $`\stackrel{~}{\zeta }_{jr}`$ the reductions of the functions $`\stackrel{~}{\zeta }_j`$ by the $`S^1`$ symmetry. Then $`\stackrel{~}{\zeta }_{jr}=0`$ on either pair of cotangent vectors in $`C_r`$. Moreover, by the uniqueness statement on $`C`$ it follows that $`C_r`$ is the unique canonical relation in $`T^{}(\stackrel{~}{M}\times \stackrel{~}{M})`$ with the given set of real points and in the zero set of the $`\stackrel{~}{\zeta }_{jr}`$’s.
We now observe that $`C_r`$ has, at least near the diagonal, a unique global generating function. This holds because the natural projection
$$C_rT^{}(\stackrel{~}{M}\times \stackrel{~}{M})\stackrel{~}{M}\times \stackrel{~}{M}$$
(47)
is a local diffeomorphism near the diagonal. Indeed, its derivative gives a natural isomorphism
$$T_{\rho ,\rho }C_rH^{}H^{}T(\stackrel{~}{M}\times \stackrel{~}{M}).$$
(48)
Therefore, there exists a global generating function $`\mathrm{log}\stackrel{~}{a}C^{\mathrm{}}(\stackrel{~}{M}\times \stackrel{~}{M})`$ i.e.
$$C_r=\{(\stackrel{~}{z},d_{\stackrel{~}{z}}\mathrm{log}\stackrel{~}{a},\stackrel{~}{w},d_{\stackrel{~}{w}}\mathrm{log}\stackrel{~}{a}),\stackrel{~}{z},\stackrel{~}{w}\stackrel{~}{M}\}.$$
(49)
Since $`C^{}=C`$ it follows that $`C_r^{}=C^r`$ and hence that $`a(w,z)=\overline{a(z,w)}.`$
Working backwards, we find that the function $`\stackrel{~}{\psi }(\stackrel{~}{x},\stackrel{~}{y})=i(1\lambda \overline{\mu }\stackrel{~}{a}(\stackrel{~}{z},\stackrel{~}{w}))`$ satisfies the equations $`\stackrel{~}{\zeta }_j(\stackrel{~}{x},d_{\stackrel{~}{x}}\stackrel{~}{\psi })=\stackrel{~}{\zeta }_j(\stackrel{~}{y},d_{\stackrel{~}{y}}\stackrel{~}{\psi })=0`$ on $`\stackrel{~}{\psi }=0.`$ Therefore the Lagrange immersion
$$\begin{array}{c}i_{\stackrel{~}{\psi }}:C_{t\stackrel{~}{\psi }}=^+\times \{\stackrel{~}{\psi }=0\}T^{}(\stackrel{~}{X}\times \stackrel{~}{X})\hfill \\ \\ (t,\stackrel{~}{x},\stackrel{~}{y})(\stackrel{~}{x},td_{\stackrel{~}{x}}\stackrel{~}{\psi };\stackrel{~}{y},td_{\stackrel{~}{y}}\stackrel{~}{\psi })\hfill \end{array}$$
(50)
takes its image inside $`𝒥_+\times \overline{𝒥_+}`$ and reduces to $`C_r`$ under the $`S^1`$-symmetry. To conclude the proof it is only necessary to show that the real points of the image of $`i_{\stackrel{~}{\psi }}`$ equal $`\mathrm{diag}(\mathrm{\Sigma }).`$ We know however that these real points reduce to $`(C_r)_{}`$ and hence that $`z=w`$ at real points. But we have
$$1=\lambda \overline{\mu }a(z,w)=e^{i(\theta \phi )}\frac{a(z,w)}{\sqrt{a(z)}\sqrt{a(w)}},\text{on}\{\stackrel{~}{\psi }=0\}$$
hence when $`z=w`$ we have $`e^{i(\theta \phi )}=1`$ and hence $`x=y`$. Since $`d_{\stackrel{~}{x}}\stackrel{~}{\psi }(x,y)|_{x=y}=\alpha _x`$, it follows that the real points indeed equal $`\mathrm{diag}(\mathrm{\Sigma }).`$ Therefore $`t\stackrel{~}{\psi }`$ parametrizes $`C`$ and hence there exists a classical symbol for which $`\mathrm{\Pi }(x,y)`$ has the stated oscillatory integral representation.
To show that the phase is of positive type, we need to describe the asymptotics of $`a(z,w)`$ near the diagonal. Note that in the almost-complex case, we cannot describe $`a(z,w)`$ as the almost analytic extension of $`a(z,z)`$. (Of course, $`\stackrel{~}{a}(\stackrel{~}{z},\stackrel{~}{w})`$ is the almost analytic extension of $`a(z,w)`$, by definition.) For our near-diagonal asymptotics in the nonintegrable case, we instead use the following second order expansion of $`a`$ at points on the diagonal:
###### Lemma 2.2.
Suppose that $`(z_1,\mathrm{},z_m)`$ are preferred coordinates and $`e_L`$ is a preferred frame at a point $`P_0M`$. Then the Taylor expansion of $`a(z,w)`$ at $`z=w=0`$ is
$$a(z,w)=1+z\overline{w}+\mathrm{}.$$
###### Proof.
To begin, we recall that $`a(0,0)=a(0)=e_L^{}(P_0)^2=1`$. To compute the first and second order terms, we return to the equation
$$\zeta _j(z,\lambda ,\frac{d\lambda }{\lambda }+d_z\mathrm{log}a(z,w))=0,(z,\lambda ;w)X\times M.$$
(51)
Let us write $`\zeta _j=\zeta _j^{(1)}+R_j^{(2)},`$ where $`R_j^{(2)}`$ vanishes to second order on $`\mathrm{\Sigma }`$ and we recall that $`\zeta _j^{(1)}(\xi )=(\overline{Z}_j,\xi )`$. Let us also Taylor expand $`\mathrm{log}a`$:
$$\mathrm{log}a=L(z,w)+Q(z,w)+\mathrm{},$$
where $`L`$ is linear and $`Q`$ is quadratic. Since $`e_L`$ is a preferred frame at $`P_0`$, it follows from (6) that $`a(z,z)=1+|z|^2+\mathrm{}`$ and hence
$$L(z,z)=0,Q(z,z)=|z|^2.$$
(52)
Since $`d_z\mathrm{log}a|_{z=w}+\frac{d\lambda }{\lambda }=\alpha \mathrm{\Sigma }`$, it follows from (51) that
$$\zeta _j^{(1)}(z,\lambda ,\frac{d\lambda }{\lambda }+d_z\mathrm{log}a)=R_j^{(2)}(z,\lambda ,\frac{d\lambda }{\lambda }+d_z\mathrm{log}a)=O(|zw|^2).$$
(53)
Since $`a(z,w)=\overline{a(w,z)}`$, we can write
$$L(z,w)=\underset{j=1}{\overset{m}{}}(b_jz_j+c_j\overline{z}_j+\overline{c}_jw_j+\overline{b}_j\overline{w}_j).$$
Since the $`z_j`$ are preferred coordinates and $`e_L`$ is a preferred frame at $`P_0`$, we can choose the $`\overline{Z}_j`$ so that $`\overline{Z}_j(0)=\frac{}{\overline{z}_j}`$ and hence by (53),
$$0=\zeta _j^{(1)}(z,\lambda ,\frac{d\lambda }{\lambda }+d_z\mathrm{log}a)|_{z=w=0,\lambda =1}=(\frac{}{\overline{z}_j},d_z\mathrm{log}a)|_{(0,0)}=c_jj.$$
Since $`L(z,z)=0`$, we have $`b_j+\overline{c}_j=0`$, and hence $`L=0`$.
To investigate the quadratic term $`Q`$ in (52), we write
$$(\frac{d\lambda }{\lambda }+d_z\mathrm{log}a)|_{(z,w)}=\alpha _z+\underset{j=1}{\overset{m}{}}\left[z_jU_j^{}+\overline{z}_jU_j^{\prime \prime }+w_jV_j^{}+\overline{w}_jV_j^{\prime \prime }\right]+O(|z|^2+|w|^2),$$
(54)
where
$$\begin{array}{cccccc}U_j^{}\hfill & =& _{k=1}^m\left(\frac{^2Q}{z_jz_k}dz_k+\frac{^2Q}{z_j\overline{z}_k}d\overline{z}_k\right),\hfill & U_j^{\prime \prime }\hfill & =& _{k=1}^m\left(\frac{^2Q}{\overline{z}_jz_k}dz_k+\frac{^2Q}{\overline{z}_j\overline{z}_k}d\overline{z}_k\right),\hfill \\ V_j^{}\hfill & =& _{k=1}^m\left(\frac{^2Q}{w_jz_k}dz_k+\frac{^2Q}{w_j\overline{z}_k}d\overline{z}_k\right),\hfill & V_j^{\prime \prime }\hfill & =& _{k=1}^m\left(\frac{^2Q}{\overline{w}_jz_k}dz_k+\frac{^2Q}{\overline{w}_j\overline{z}_k}d\overline{z}_k\right).\hfill \end{array}$$
Applying $`\zeta _k^{(1)}`$ to (54) and using (53) and the fact that $`\zeta _k^{(1)}(\alpha _z)=0`$, we have
$$\underset{j=1}{\overset{m}{}}\left[z_j(\overline{Z}_k|_z,U_j^{})+\overline{z}_j(\overline{Z}_k|_z,U_j^{\prime \prime })+w_j(\overline{Z}_k|_z,V_j^{})+\overline{w}_j(\overline{Z}_k|_z,V_j^{\prime \prime })\right]=O(|z|^2+|w|^2).$$
(55)
By (3) and (8),
$$\overline{Z}_k|_z=\frac{}{\overline{z}_k}+\underset{l=1}{\overset{m}{}}B_{kl}(z)\frac{}{z_l}+C_k(z)\frac{}{\theta },B_{kl}(0)=0.$$
Hence by (55),
$$\frac{^2Q}{z_j\overline{z}_k}=(\frac{}{\overline{z}_k},U_j^{})=(\overline{Z}_k|_0,U_j^{})=0.$$
Similarly, $`\frac{^2Q}{\overline{z}_j\overline{z}_k}=\frac{^2Q}{w_j\overline{z}_k}=\frac{^2Q}{\overline{w}_j\overline{z}_k}=0`$. Thus $`Q(z,w)`$ has no terms containing $`\overline{z}_k`$. Since $`Q(z,w)=\overline{Q(w,z)}`$, the quadratic function $`Q`$ also has no terms containing $`w_k`$, so we can write
$$Q(z,w)=B(z,z)+H(z,\overline{w})+\overline{B(w,w)},$$
where $`B`$, resp. $`H`$, is a bilinear, resp. hermitian, form on $`^m`$. Since $`Q(z,z)=|z|^2`$ (recall (52)), we conclude that $`B(z,z)=0`$ and hence $`Q(z,w)=H(z,\overline{w})=z\overline{w}`$. ∎
To complete the proof of Theorem 2.1, it remains to show that the phase is of positive type; i.e., $`\mathrm{}\psi 0`$ on some neighborhood of the diagonal in $`X\times X`$. Let $`xX`$ be arbitrary and choose Heisenberg coordinates $`(z,\theta )`$ at $`P_0=\pi (x)`$ (so that $`x`$ has coordinates $`(0,0)`$). Recalling that $`\lambda =a(z)^{\frac{1}{2}}e^{i\theta }`$ on $`X`$, we have by Lemma 2.2,
$$\frac{1}{i}\psi (0,0;z,\theta )=1\frac{a(0,z)}{\sqrt{a(z)}}e^{i\theta }=(1e^{i\theta })+e^{i\theta }\left[\frac{1}{2}|z|^2+O(|z|^3)\right].$$
Thus,
$$\mathrm{}\left[\frac{1}{i}\psi (0,0;z,\theta )\right]0\text{for }|\theta |<\frac{\pi }{2},|z|<\epsilon ,$$
where $`\epsilon `$ is independent of the point $`P_0M`$. ∎
### 2.2. Scaling limit of the Szegö kernel
The Szegö kernels $`\mathrm{\Pi }_N`$ are the Fourier coefficients of $`\mathrm{\Pi }`$ defined by:
$$\mathrm{\Pi }_N(x,y)=_0^{\mathrm{}}_0^{2\pi }e^{iN\theta }e^{it\psi (r_\theta x,y)}s(r_\theta x,y,t)𝑑\theta 𝑑t$$
(56)
where $`r_\theta `$ denotes the $`S^1`$ action on $`X`$. Changing variables $`tNt`$ gives
$$\mathrm{\Pi }_N(x,y)=N_0^{\mathrm{}}_0^{2\pi }e^{iN(\theta +t\psi (r_\theta x,y))}s(r_\theta x,y,Nt)𝑑\theta 𝑑t.$$
(57)
We now determine the scaling limit of the Szegö kernel by the argument of \[BSZ2\]. For the sake of completeness, we provide the details of the argument and add some new details on homogeneities, which are useful in applications. To describe the scaling limit at a point $`x_0X`$, we choose a Heisenberg chart $`\rho :U,0X,x_0`$ centered at $`P_0=\pi (x_0)M`$. Recall (§1.2) that choosing $`\rho `$ is equivalent to choosing preferred coordinates centered at $`P_0`$ and a preferred local frame $`e_L`$ at $`P_0`$. We then write the Szegö kernel $`\mathrm{\Pi }_N`$ in terms of these coordinates:
$$\mathrm{\Pi }_N^{P_0}(u,\theta ;v,\phi )=\mathrm{\Pi }_N(\rho (u,\theta ),\rho (v,\phi )),$$
where the superscript $`P_0`$ is a reminder that we are using coordinates centered at $`P_0`$. (We remark that the function $`\mathrm{\Pi }_N^{P_0}`$ depends also on the choice of preferred coordinates and preferred frame, which we omit from the notation.) The first term in our asymptotic formula below says that the $`N^{\mathrm{th}}`$ scaled Szegö kernel looks approximately like the Szegö kernel of level one for the reduced Heisenberg group (recall (41)):
$$\mathrm{\Pi }_N^{P_0}(\frac{u}{\sqrt{N}},\frac{\theta }{N};\frac{v}{\sqrt{N}},\frac{\phi }{N})\mathrm{\Pi }_1^𝐇(u,\theta ;v,\phi )=\frac{1}{\pi ^m}e^{i(\theta \phi )+i\mathrm{}(u\overline{v})\frac{1}{2}|uv|^2}.$$
In the following, we shall denote the Taylor series of a $`𝒞^{\mathrm{}}`$ function $`f`$ defined in a neighborhood of $`0^K`$ by $`ff_0+f_1+f_2+\mathrm{}`$ where $`f_j`$ is the homogeneous polynomial part of degree $`j`$. We also denote by $`R_n^ff_{n+1}+\mathrm{}`$ the remainder term in the Taylor expansion.
The following is our main result on the scaling asymptotics of the Szegö kernels near the diagonal. Since the result is of independent interest, we state our asymptotic formula in a more precise form than is needed for the applications in this paper.
###### Theorem 2.3.
Let $`P_0M`$ and choose a Heisenberg coordinate chart about $`P_0`$. Then
$$\begin{array}{c}N^m\mathrm{\Pi }_N^{P_0}(\frac{u}{\sqrt{N}},\frac{\theta }{N};\frac{v}{\sqrt{N}},\frac{\phi }{N})\hfill \\ \\ =\mathrm{\Pi }_1^𝐇(u,\theta ;v,\phi )\left[1+_{r=1}^KN^{r/2}b_r(P_0,u,v)+N^{(K+1)/2}R_K(P_0,u,v,N)\right],\hfill \end{array}$$
where:
* $`b_r=_{\alpha =0}^{2[r/2]}_{j=0}^{[3r/2]}(\psi _2)^\alpha Q_{r,\alpha ,3r2j},`$ where $`Q_{r,\alpha ,d}`$ is homogeneous of degree $`d`$ and
$$\psi _2(u,v)=u\overline{v}\frac{1}{2}(|u|^2+|v|^2);$$
in particular, $`b_r`$ has only even homogeneity if $`r`$ is even, and only odd homogeneity if $`r`$ is odd;
* $`R_K(P_0,u,v,N)_{𝒞^j(\{|u|\rho ,|v|\rho \}}C_{K,j,\rho }`$ for $`j0,\rho >0`$ and $`C_{K,j,\rho }`$ is independent of the point $`P_0`$ and choice of coordinates.
###### Proof.
We now fix $`P_0`$ and consider the asymptotics of
$$\begin{array}{c}\mathrm{\Pi }_N(\frac{u}{\sqrt{N}},0;\frac{v}{\sqrt{N}},0)\hfill \\ =N_0^{\mathrm{}}_0^{2\pi }e^{iN(\theta +t\psi (\frac{u}{\sqrt{N}},\theta ;\frac{v}{\sqrt{N}},0))}s(\frac{u}{\sqrt{N}},\theta ;\frac{v}{\sqrt{N}},0),Nt)d\theta dt,\hfill \end{array}$$
(58)
where $`\psi `$ and $`s`$ are the phase and symbol from Theorem 2.1 written in terms of the Heisenberg coordinates.
On $`X`$ we have $`\lambda =a(z)^{\frac{1}{2}}e^{i\phi }`$. So for $`(x,y)=(z,\phi ,w,\phi ^{})X\times X`$, we have by Theorem 2.1,
$$\psi (z,\phi ,w,\phi ^{})=i\left[1\frac{a(z,w)}{\sqrt{a(z)}\sqrt{a(w)}}e^{i(\phi \phi ^{})}\right].$$
(59)
It follows that
$$\begin{array}{c}\psi (\frac{u}{\sqrt{N}},\theta ;\frac{v}{\sqrt{N}},0)\hfill \\ =i\left[1\frac{a(\frac{u}{\sqrt{N}},\frac{v}{\sqrt{N}})}{\sqrt{a(\frac{u}{\sqrt{N}},\frac{u}{\sqrt{N}})}\sqrt{a(\frac{v}{\sqrt{N}},\frac{v}{\sqrt{N}})}}e^{i\theta }\right].\hfill \end{array}$$
(60)
We observe that the asymptotic expansion of a function $`f(\frac{u}{\sqrt{N}},\frac{v}{\sqrt{N}})`$ in powers of $`N^{\frac{1}{2}}`$ is just the Taylor expansion of $`f`$ at $`u=v=0`$. By Lemma 2.2 and the notational convention established above, we have
$$a(\frac{u}{\sqrt{N}},\frac{v}{\sqrt{N}})=1+\frac{1}{N}u\overline{v}+R_3^a(\frac{u}{\sqrt{N}},\frac{v}{\sqrt{N}}),R_3^a(\frac{u}{\sqrt{N}},\frac{v}{\sqrt{N}})=O(N^{3/2}).$$
(61)
The entire phase
$$\begin{array}{c}t\psi (\frac{u}{\sqrt{N}},\theta ;\frac{v}{\sqrt{N}},0)\theta =it\left[1\frac{a(\frac{u}{\sqrt{N}},\frac{v}{\sqrt{N}})}{a(\frac{u}{\sqrt{N}},\frac{u}{\sqrt{N}})^{\frac{1}{2}}a(\frac{v}{\sqrt{N}},\frac{v}{\sqrt{N}})^{\frac{1}{2}}}e^{i\theta }\right]\theta \hfill \end{array}$$
(62)
then has the asymptotic $`N`$-expansion
$$it[1e^{i\theta }]\theta \frac{it}{N}\psi _2(u,v)e^{i\theta }+tR_3^\psi (\frac{u}{\sqrt{N}},\frac{v}{\sqrt{N}})e^{i\theta }.$$
(63)
As in \[BSZ2\], we absorb $`(i\psi _2+NR_3^\psi )te^{i\theta }`$ into the amplitude, so that $`\mathrm{\Pi }_N^{P_0}(\frac{u}{\sqrt{N}},0;\frac{v}{\sqrt{N}},0)`$ is an oscillatory integral with phase
$$\mathrm{\Psi }(t,\theta ):=it(1e^{i\theta })\theta $$
and with amplitude
$$A(t,\theta ;P_0,u,v):=Ne^{te^{i\theta }\psi _2(u,v)+ite^{i\theta }NR_3^\psi (\frac{u}{\sqrt{N}},\frac{v}{\sqrt{N}})}\underset{k=0}{\overset{\mathrm{}}{}}N^{mk}t^{mk}s_k(\frac{u}{\sqrt{N}},\frac{v}{\sqrt{N}},\theta );$$
i.e.,
$$\mathrm{\Pi }_N(\frac{u}{\sqrt{N}},0;\frac{v}{\sqrt{N}},0)=_0^{\mathrm{}}_0^{2\pi }e^{iN\mathrm{\Psi }(t,\theta )}A(t,\theta ;P_0,u,v)𝑑\theta 𝑑t$$
(64)
Before proceeding, it is convenient to expand $`\mathrm{exp}\left[ite^{i\theta }NR_3^\psi (\frac{u}{\sqrt{N}},\frac{v}{\sqrt{N}})\right]`$ in powers of $`N^{\frac{1}{2}}`$ and to keep track of the homogeneity in $`(u,v)`$ of the coefficients. We simplify the notation by writing $`g(t,\theta ):=ite^{i\theta }`$. By definition,
$$R_3^\psi (\frac{u}{\sqrt{N}},\frac{v}{\sqrt{N}})N^{3/2}\psi _3(u,v)+N^2\psi _4(u,v)+\mathrm{}+N^{d/2}\psi _d(u,v)+\mathrm{}.$$
We then have
$$e^{NgR_3^\psi (\frac{u}{\sqrt{N}},\frac{v}{\sqrt{N}})}\underset{r=0}{\overset{\mathrm{}}{}}N^{r/2}c_r(u,v),$$
(65)
where
$$c_r=\underset{\lambda =1}{\overset{r}{}}c_{r,r+2\lambda }(u,v;t,\theta ),r1,c_0=c_{00}=1,$$
(66)
with $`c_{rd}`$ homogeneous of degree $`d`$ in $`u,v`$. (The explicit formula for $`c_{rd}`$ is:
$$c_{rd}=\{\frac{g^n}{n!}\mathrm{\Pi }_{j=1}^n\psi _{a_j}(u,v):n1,a_j3,_{j=1}^na_j=d,_{j=1}^n(a_j2)=r\},r1.$$
The range of $`d`$ is determined by the fact that $`d=_{j=1}^na_j=r+2n`$ with $`0nr`$.)
We further decompose the factor $`_{k=0}^{\mathrm{}}N^{mk}t^{mk}s_k(\frac{u}{\sqrt{N}},\frac{v}{\sqrt{N}},\theta )`$ into the homogeneous terms $`_{k,\mathrm{}=0}^{\mathrm{}}N^{mk\mathrm{}/2}t^{mk}s_k\mathrm{}(P_0,u,v)`$ where $`s_k\mathrm{}`$ is the homogeneous term of degree $`\mathrm{}`$ of $`s_k`$. Finally, we have
$`A`$ $``$ $`Ne^{gNR_3^\psi (\frac{u}{\sqrt{N}},\frac{v}{\sqrt{N}},\theta )}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}N^{mk}t^{mk}s_k({\displaystyle \frac{u}{\sqrt{N}}},{\displaystyle \frac{v}{\sqrt{N}}},\theta )`$
$`=`$ $`N^{m+1}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}N^{n/2}f_n(u,v;t,\theta ,P_0),`$
$$f_n=\underset{r+\mathrm{}+2k=n}{}c_rs_k\mathrm{}=\underset{k=0}{\overset{[n/2]}{}}t^{mk}\left(s_{k,n2k}+\underset{r=1}{\overset{n2k}{}}\underset{\lambda =1}{\overset{r}{}}c_{r,r+2\lambda }s_{k,n2kr}\right)=\underset{j=0}{\overset{[3n/2]}{}}f_{n,3n2j}$$
where $`f_{n,d}`$ is homogeneous of degree $`d`$ in $`(u,v)`$. (The asymptotic expansion holds in the sense of semiclassical symbols, i.e. the remainder after summing $`K`$ terms is a symbol of order $`mk(K+1)/2`$.)
We now evaluate the integral for $`\mathrm{\Pi }_N`$ by the method of stationary phase as in \[BSZ2\]. The phase is independent of the parameters $`(u,v)`$ and we have
$$\begin{array}{c}\frac{}{t}\mathrm{\Psi }=i(1e^{i\theta })\hfill \\ \frac{}{\theta }\mathrm{\Psi }=te^{i\theta }1\hfill \end{array}$$
(67)
so the critical set of the phase is the point $`\{t=1,\theta =0\}`$. The Hessian $`\mathrm{\Psi }^{\prime \prime }`$ on the critical set equals
$$\left(\begin{array}{cc}0\hfill & 1\hfill \\ 1\hfill & i\hfill \end{array}\right)$$
so the phase is non-degenerate and the Hessian operator $`L_\mathrm{\Psi }`$ is given by
$$L_\mathrm{\Psi }=\mathrm{\Psi }^{\prime \prime }(1,0)^1D,D=2\frac{^2}{t\theta }i\frac{^2}{t^2}.$$
We smoothly decompose the integral into one over $`|t1|<1`$ and one over $`|t1|>\frac{1}{2}`$. Since the only critical point of the phase occurs at $`t=1,\theta =0`$, the latter is rapidly decaying in $`N`$ and we may assume the integrand to be smoothly cut off to $`|t1|<1`$. It follows by the stationary phase method for complex oscillatory integrals (\[Hör\], Theorem 7.7.5) that
$$N^m\mathrm{\Pi }_N^{P_0}(\frac{u}{\sqrt{N}},\frac{v}{\sqrt{N}},\theta )=C\underset{j=0}{\overset{J}{}}\underset{n=0}{\overset{K}{}}N^{n/2j}L_j[e^{ig\psi _2}f_n]|_{t=1,\theta =0}+\widehat{R}_{JK}(P_0,u,v,N),$$
(68)
where
$$C=N\frac{1}{\sqrt{\mathrm{det}(\mathrm{N}\mathrm{\Psi }^{\prime \prime }(1,0)/2\pi \mathrm{i})}}=\sqrt{2\pi i}$$
and $`L_j`$ is the differential operator of order $`2j`$ in $`(t,\theta )`$ defined by
$$L_jf(t,\theta ;P_0,u,v)=\underset{\nu \mu =j}{}\underset{2\nu 3\mu }{}\frac{1}{2^\nu i^j\mu !\nu !}L_\mathrm{\Psi }^\nu [f(t,\theta ;P_0,u,v)(R_3^\mathrm{\Psi })^\mu (t,\theta )]$$
(69)
with $`R_3^\mathrm{\Psi }(t,\theta )`$ the third order remainder in the Taylor expansion of $`\mathrm{\Psi }`$ at $`(t,\theta )=(1,0).`$ Also, the remainder is estimated by
$$|\widehat{R}_{JK}(P_0,u,v,N)|C^{}N^{J\frac{K+1}{2}}\underset{n=0}{\overset{K}{}}\underset{|\alpha |2J+2}{}\underset{t,\theta }{sup}|D_{t,\theta }^\alpha e^{ig\psi _2}f_n|.$$
(70)
Since $`L_\mathrm{\Psi }`$ is a second order operator in $`(t,\theta )`$, we see that
$$L_j[e^{ig\psi _2}f_n]|_{t=1,\theta =0}=e^{\psi _2}\underset{\alpha 2j}{}(\psi _2)^\alpha F_{nj\alpha }.$$
(71)
Therefore
$`N^m\mathrm{\Pi }_N^{P_0}({\displaystyle \frac{u}{\sqrt{N}}},{\displaystyle \frac{v}{\sqrt{N}}},\theta )`$ $``$ $`e^{\psi _2}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\alpha =0}{\overset{2j}{}}}(\psi _2)^\alpha N^{\frac{n}{2}j}F_{nj\alpha }`$ (72)
$``$ $`e^{\psi _2}{\displaystyle \underset{r=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{j=0}{\overset{[r/2]}{}}}{\displaystyle \underset{\alpha =0}{\overset{2j}{}}}(\psi _2)^\alpha N^{r/2}F_{r2j,j,\alpha }`$
$``$ $`e^{\psi _2}{\displaystyle \underset{r=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\alpha =0}{\overset{2[r/2]}{}}}(\psi _2)^\alpha N^{r/2}Q_{r\alpha }.`$
Thus, as with $`f_n`$ we have the homogeneous expansion:
$$Q_{r\alpha }=\underset{j=0}{\overset{[3r/2]}{}}Q_{r,\alpha ,3r2j}.$$
(73)
Here, $`Q_{r,\alpha ,d}`$ is homogeneous of degree $`d`$ in $`(u,v)`$. (The term $`\psi _2`$ is distinguished by being ‘holomorphic’ in $`u`$ and ‘anti-holomorphic’ in $`v`$ in a sense to be elaborated below.) Thus we have the desired Taylor series. The estimate for the remainder follows from (70). ∎
## 3. Kodaira embedding and Tian almost isometry theorem
Definition: By the Kodaira maps we mean the maps $`\mathrm{\Phi }_N:MPH^0(M,L^N)^{}`$ defined by $`\mathrm{\Phi }_N(z)=\{s^N:s^N(z)=0\}`$. Equivalently, we can choose an orthonormal basis $`S_1^N,\mathrm{},S_{d_N}^N`$ of $`H^0(M,L^N)`$ and write
$$\mathrm{\Phi }_N:M^{d_N1},\mathrm{\Phi }_N(z)=(S_1^N(z):\mathrm{}:S_{d_N}^N(z)).$$
(74)
We also define the lifts of the Kodaira maps:
$$\stackrel{~}{\mathrm{\Phi }}_N:X^{d_N},\stackrel{~}{\mathrm{\Phi }}_N(x)=(S_1^N(x),\mathrm{},S_{d_N}^N(x)).$$
(75)
Note that
$$\mathrm{\Pi }_N(x,y)=\stackrel{~}{\mathrm{\Phi }}_N(x)\overline{\stackrel{~}{\mathrm{\Phi }}_N(y)};$$
(76)
in particular,
$$\mathrm{\Pi }_N(x,x)=\stackrel{~}{\mathrm{\Phi }}_N(x)^2.$$
(77)
We now prove the following generalization to the symplectic category of the asymptotic expansion theorem of \[Zel\] (also proved independently by \[Cat\] using the Bergman kernel in place of the Szegö kernel) and Tian’s approximate isometry theorem \[Tian\]:
###### Theorem 3.1.
Let $`L(M,\omega )`$ be the pre-quantum line bundle over a $`2m`$-dimensional symplectic manifold, and let $`\{\mathrm{\Phi }_N\}`$ be its Kodaira maps. Then:
(a) There exists a complete asymptotic expansion:
$$\mathrm{\Pi }_N(z,0;z,0)=a_0N^m+a_1(z)N^{m1}+a_2(z)N^{m2}+$$
for certain smooth coefficients $`a_j(z)`$ with $`a_0=\pi ^m`$. Hence, the maps $`\mathrm{\Phi }_N`$ are well-defined for $`N0`$.
(b) Let $`\omega _{FS}`$ denote the Fubini-Study form on $`^{d_N1}`$. Then
$$\frac{1}{N}\mathrm{\Phi }_N^{}(\omega _{FS})\omega _{𝒞^k}=O(\frac{1}{N})$$
for any $`k`$.
###### Proof.
(a) Using the expansion of Theorem 2.3 with $`u=v=0`$ and noting that $`b_r(z,0,0)=0`$ for $`r`$ odd, we obtain the above expansion of $`\mathrm{\Pi }_N(z,0;z,0)`$ with $`a_r(z)=b_{2r}(z,0,0)`$. (The expansion also follows by precisely the same proof as in \[Zel\].)
(b) In the holomorphic case, (b) followed by differentiating (a), using that $`\mathrm{\Phi }_N^{}(\overline{}\mathrm{log}|\xi |^2)=\overline{}\mathrm{log}|\mathrm{\Phi }_N|^2`$. In the almost complex case, $`\mathrm{\Phi }_N^{}`$ does not commute with the complex derivatives, so we need to modify the proof. To do so, we use the following notation: the exterior derivative on a product manifold $`Y_1\times Y_2`$ can be decomposed as $`d=d^1+d^2`$, where $`d^1`$ and $`d^2`$ denote exterior differentiation on the first and second factors, respectively. (This is formally analogous to the decomposition $`d=+\overline{}`$; e.g., $`d^1d^1=d^2d^2=d^1d^2+d^2d^1=0`$.)
Recall that the Fubini-Study form $`\omega _{FS}`$ on $`^{m1}`$ is induced by the 2-form $`\stackrel{~}{\omega }_m=\frac{i}{2}\overline{}\mathrm{log}|\xi |^2`$ on $`^m\{0\}`$. We consider the 2-form $`\mathrm{\Omega }`$ on $`(^m\{0\})\times (^m\{0\})`$ given by
$$\mathrm{\Omega }=\frac{i}{2}\overline{}\mathrm{log}\zeta \overline{\eta }=\frac{i}{2}d^1d^2\mathrm{log}\zeta \overline{\eta }.$$
Note that $`\mathrm{\Omega }`$ is smooth on a neighborhood of the diagonal $`\{\zeta =\eta \}`$, and
$$\mathrm{\Omega }|_{\zeta =\eta }=\stackrel{~}{\omega }_m$$
(where the restriction to $`\{\zeta =\eta \}`$ means the pull-back under the map $`\zeta (\zeta ,\zeta )`$).
It suffices to show that
$$\frac{1}{N}\stackrel{~}{\mathrm{\Phi }}_N^{}\omega _{d_N}\pi ^{}\omega ,\pi :XM.$$
To do this, we consider the maps
$$\mathrm{\Psi }_N=\stackrel{~}{\mathrm{\Phi }}_N\times \stackrel{~}{\mathrm{\Phi }}_N:X\times X^{d_N}\times ^{d_N},\mathrm{\Psi }_N(x,y)=(\stackrel{~}{\mathrm{\Phi }}_N(x),\stackrel{~}{\mathrm{\Phi }}_N(y)).$$
It is elementary to check that $`\mathrm{\Psi }_N^{}`$ commutes with $`d^1`$ and $`d^2`$. By (76), we have
$$\mathrm{\Psi }_N^{}(\mathrm{log}\zeta \overline{\eta })=(\mathrm{log}\zeta \overline{\eta })\mathrm{\Psi }_N=\mathrm{log}\mathrm{\Pi }_N.$$
Therefore,
$$\frac{1}{N}\stackrel{~}{\mathrm{\Psi }}_N^{}\mathrm{\Omega }_{d_N}=\frac{i}{2N}\mathrm{\Psi }_N^{}d^1d^2\mathrm{log}\zeta \overline{\eta }=\frac{i}{2N}d^1d^2\mathrm{\Psi }_N^{}\mathrm{log}\zeta \overline{\eta }=\frac{i}{2N}d^1d^2\mathrm{log}\mathrm{\Pi }_N.$$
(78)
Restricting (78) to the diagonal, we then have
$$\frac{1}{N}\stackrel{~}{\mathrm{\Phi }}_N^{}\omega _{d_N}=\frac{i}{2N}(d^1d^2\mathrm{log}\mathrm{\Pi }_N)|_{x=y}=\mathrm{diag}^{}(d^1d^2\mathrm{log}\mathrm{\Pi }_N),$$
where $`\mathrm{diag}:XX\times X`$ is the diagonal map $`\mathrm{diag}(x)=(x,x)`$.
Using Heisenberg coordinates as in Theorem 2.3, we have by the near-diagonal scaling asymptotics
$`{\displaystyle \frac{1}{N}}\stackrel{~}{\mathrm{\Phi }}_N^{}\omega _{d_N}|_{P_0}`$ $`=`$ $`{\displaystyle \frac{i}{2N}}\mathrm{diag}^{}d^1d^2\mathrm{log}\mathrm{\Pi }_N^{P_0}({\displaystyle \frac{u}{\sqrt{N}}},{\displaystyle \frac{\theta }{N}};{\displaystyle \frac{v}{\sqrt{N}}},{\displaystyle \frac{\phi }{N}})|_0`$
$`=`$ $`{\displaystyle \frac{i}{2N}}\mathrm{diag}^{}d^1d^2\mathrm{log}\mathrm{\Pi }_1^𝐇(u,\theta ;v,\phi )|_0+O(N^{\frac{1}{2}}).`$
Finally,
$`{\displaystyle \frac{i}{2N}}\mathrm{diag}^{}d^1d^2\mathrm{log}\mathrm{\Pi }_1^𝐇(u,\theta ;v,\phi )|_0`$ $`=`$ $`{\displaystyle \frac{i}{2N}}\mathrm{diag}^{}d^1d^2\left[i(\theta \phi )+u\overline{v}{\displaystyle \frac{1}{2}}(|u|^2+|v|^2)\right]`$ (79)
$`=`$ $`{\displaystyle \frac{i}{2N}}{\displaystyle \underset{q=1}{\overset{m}{}}}du_qd\overline{u}_q={\displaystyle \frac{i}{2}}{\displaystyle \underset{q=1}{\overset{m}{}}}dz_qd\overline{z}_q=\omega |_{P_0}.`$
Remark: A more explicit way to show (b) is to expand the Fubini-Study form:
$$\stackrel{~}{\omega }_m=\frac{i}{2}|\xi |^4\left[|\xi |^2\underset{j=1}{\overset{m}{}}d\xi _jd\overline{\xi }_j\underset{j,k=1}{\overset{m}{}}\overline{\xi }_j\xi _kd\xi _jd\overline{\xi }_k\right].$$
Then
$$\frac{1}{N}\stackrel{~}{\mathrm{\Phi }}_N^{}\omega _{d_N}=\frac{i}{2}\mathrm{\Pi }_N(x,x)^2\{(\mathrm{\Pi }_N(x,x)d^1d^2\mathrm{\Pi }_N(x,y)d^1\mathrm{\Pi }_N(x,y)d^2\mathrm{\Pi }_N(x,y)\}|_{x=y},$$
and (b) follows from a short computation using Theorem 2.3 as above.
It follows from Theorem 3.1(b) that $`\mathrm{\Phi }_N`$ is an immersion for $`N0`$. Using in part an idea of Bouche \[Bch\], we give a simple proof of the ‘Kodaira embedding theorem’ for symplectic manifolds:
###### Theorem 3.2.
For $`N`$ sufficiently large, $`\mathrm{\Phi }_N`$ is an embedding.
###### Proof.
Let $`\{P_N,Q_N\}`$ be any sequence of distinct points such that $`\mathrm{\Phi }_N(P_N)=\mathrm{\Phi }_N(Q_N)`$. By passing to a subsequence we may assume that one of the following two cases holds:
1. The distance $`r_N:=r(P_N,Q_N)`$ between $`P_N,Q_N`$ satisfies $`r_N\sqrt{N}\mathrm{};`$
2. There exists a constant $`C`$ independent of $`N`$ such that $`r_NC\sqrt{N}.`$
To prove that case (i) cannot occur, we observe that
$$_{B(P_N,r_N)}|N^m\mathrm{\Pi }_N^{P_N}|^2𝑑V1o(1)$$
where $`\mathrm{\Pi }_N^{P_N}(x)=\mathrm{\Pi }_N(,P_N)`$ is the ‘peak section’ at $`P_N`$. The same inequality holds for $`Q_N`$. If $`\mathrm{\Phi }_N(P_N)=\mathrm{\Phi }_N(Q_N)`$ then the total $`^2`$-norm of $`\mathrm{\Pi }_N(x,)`$ would have to be $`2N^m`$, contradicting the asymptotic $`N^m`$ from Theorem 3.1(a).
To prove that case (ii) cannot occur, we assume on the contrary that $`\mathrm{\Phi }_N(P_N)=\mathrm{\Phi }_N(Q_N)`$, where $`P_N=\rho _N(0)`$ and $`Q_N=\rho _N(\frac{v_N}{\sqrt{N}})`$, $`0|v_N|C`$, using a Heisenberg coordinate chart $`\rho _N`$ about $`P_N`$. We consider the function
$$f_N(t)=\frac{|\mathrm{\Pi }_N^{P_N}(0,\frac{tv_N}{\sqrt{N}})|^2}{\mathrm{\Pi }_N^{P_N}(0,0)\mathrm{\Pi }_N^{P_N}(\frac{tv_N}{\sqrt{N}},\frac{tv_N}{\sqrt{N}})}.$$
(80)
Recalling that
$$\mathrm{\Pi }_N(x,y)=\stackrel{~}{\mathrm{\Phi }}_N(x)\overline{\stackrel{~}{\mathrm{\Phi }}_N(y)},$$
we see that $`f_N(0)=1`$, which is a global and strict local maximum of $`f_N`$; furthermore, since $`\mathrm{\Phi }_N(P_N)=\mathrm{\Phi }_N(Q_N)`$, we also have $`f_N(1)=1`$. Thus for some value of $`t_N`$ in the open interval $`(0,1)`$, we have $`f_N^{\prime \prime }(t_N)=0`$. By Theorem 2.3,
$$f_N(t)=e^{|v_N|^2t^2}\left[1+N^{1/2}\stackrel{~}{R}_N(tv_N)\right],$$
(81)
where
$$\stackrel{~}{R}_N(v)=R_1(P_N;0,v,N)+R_1(P_N;v,0,N)R_1(P_N;v,v,N)R_1(P_N;0,0,N)+O(N^{1/2}).$$
The estimate for $`R_1`$ yields:
$$\stackrel{~}{R}_N_{𝒞^2\{|v|C\}}=O(1)$$
(82)
Since $`f_N(1)=1`$, it follows from (81)–(82) that $`|v_N|^2=O(N^{1/2})`$. (A more careful analysis shows that we can replace $`N^{1/2}`$ with $`N^1`$ in (81) and thus $`|v_N|=O(N^{1/2})`$.)
Write $`e^x=1+x+x^2\phi (x)`$. We then have
$$f_N(t)=1|v_N|^2t^2+|v_N|^4t^4\phi (|v_N|^2t^2)+N^{1/2}\stackrel{~}{R}_N(tv_N)\left[1|v_N|^2t^2+|v_N|^4t^4\phi (|v_N|^2t^2)\right].$$
Thus by (82),
$$f_N^{\prime \prime }(t)=2|v_N|^2+O(|v_N|^4)+O(N^{1/2}|v_N|^2),|t|1.$$
Since $`|v_N|=o(1)`$, it follows that
$$0=f_N^{\prime \prime }(t_N)=(2+o(1))|v_N|^2,$$
which contradicts the assumption that $`v_N0`$. ∎
## 4. Asymptotically holomorphic versus almost holomorphic sections
We now use the scaling asymptotics of Theorem 2.3 to prove Theorem 0.1, which states that $`\nu _{\mathrm{}}`$-almost every sequence $`\{s_N\}`$ of sections (with unit $`^2`$-norm) satisfies the sup-norm estimates
$`s_N_{\mathrm{}}+\overline{}s_N_{\mathrm{}}`$ $`=`$ $`O(\sqrt{\mathrm{log}N}),`$
$`^ks_N_{\mathrm{}}+^k\overline{}s_N_{\mathrm{}}`$ $`=`$ $`O(N^{\frac{k}{2}}\sqrt{\mathrm{log}N}),`$
for $`k=1,2,3,\mathrm{}`$.
The following elementary probability lemma is central to our arguments:
###### Lemma 4.1.
Let $`AS^{2d1}^d`$, and give $`S^{2d1}`$ Haar probability measure. Then the probability that a random point $`PS^{2d1}`$ satisfies the bound $`|P,A|>\lambda `$ is $`(1\lambda ^2)^{d1}`$.
###### Proof.
We can assume without loss of generality that $`A=(1,0,\mathrm{},0)`$. Let
$$V_\lambda =\mathrm{Vol}\left(\{PS^{2d1}:|P_1|>\lambda \}\right)(0\lambda <1),$$
where $`\mathrm{Vol}`$ denotes $`(2d1)`$-dimensional Euclidean volume. Our desired probability equals $`V_\lambda /V_0`$. Let $`\sigma _n=\mathrm{Vol}(S^{2n1})=\frac{2\pi ^n}{(n1)!}`$. We compute
$`V_\lambda `$ $`=`$ $`{\displaystyle _\lambda ^1}\sigma _{d1}(1r^2)^{\frac{2d3}{2}}{\displaystyle \frac{2\pi rdr}{\sqrt{1r^2}}}=2\pi \sigma _{d1}{\displaystyle _\lambda ^1}(1r^2)^{d2}r𝑑r`$
$`=`$ $`{\displaystyle \frac{\pi \sigma _{d1}}{d1}}(1\lambda ^2)^{d1}=\sigma _d(1\lambda ^2)^{d1}.`$
Therefore $`V_\lambda /V_0=(1\lambda ^2)^{d1}`$. ∎
### 4.1. Notation
For the readers’ convenience, we summarize here our notation for the various differential operators that we use in this section and elsewhere in the paper:
1. Derivatives on $`M`$:
* $`\frac{}{z_j}=\frac{1}{2}\frac{}{x_j}\frac{i}{2}\frac{}{y_j},\frac{}{\overline{z}_j}=\frac{1}{2}\frac{}{x_j}+\frac{i}{2}\frac{}{y_j}`$;
* $`Z_j^M=\frac{}{z_j}+\overline{B}_{jk}(z)\frac{}{\overline{z}_k},\overline{Z}_j^M=\frac{}{\overline{z}_j}+B_{jk}(z)\frac{}{z_k},B_{jk}(P_0)=0`$,
$`\{Z_1,\mathrm{},Z_m\}`$ is a local frame for $`T^{1,0}M`$.
2. Derivatives on $`X`$:
* $`\frac{^h}{z_j}=\frac{}{z_j}A_j(z)\frac{}{\theta }=`$ horizontal lift of $`\frac{}{z_j}`$, $`A_j(P_0)=0`$;
* $`Z_j=`$ horizontal lift of $`Z_j^M`$;
* $`d^h=_b+\overline{}_b`$ = horizontal exterior derivative on $`X`$.
3. Covariant derivatives on $`M`$:
* $`:𝒞^{\mathrm{}}(M,L^N(T^{}M)^k)𝒞^{\mathrm{}}(M,L^N(T^{}M)^{(k+1)})`$;
* $`^k=\mathrm{}:𝒞^{\mathrm{}}(M,L^N)𝒞^{\mathrm{}}(M,L^N(T^{}M)^k)`$;
* $`=+\overline{},\overline{}:𝒞^{\mathrm{}}(M,L^N)𝒞^{\mathrm{}}(M,L^NT^{0,1}M)`$.
4. Derivatives on $`X\times X`$:
* $`d_j^1,d_j^2`$: the operator $`\frac{^h}{z_j}`$ applied to the first and second factors, respectively;
* $`Z_j^1,Z_j^2`$: the operator $`Z_j`$ applied to the first and second factors, respectively.
### 4.2. The estimate $`s_N_{\mathrm{}}/s_N_2=O(\sqrt{\mathrm{log}N})`$ almost surely
Throughout this section we assume that $`s_N_^2=1`$. We begin the proof of Theorem 0.1 by showing that
$$\nu _N\{s^NSH_J^0(M,L^N):\underset{M}{sup}|s_N|>C\sqrt{\mathrm{log}N}\}<O\left(\frac{1}{N^2}\right),$$
(83)
for some constant $`C<+\mathrm{}`$. (In fact, for any $`k>0`$, we can bound the probabilities by $`O(N^k)`$ by choosing $`C`$ to be sufficiently large.) The estimate (83) immediately implies that
$$\underset{N\mathrm{}}{lim\; sup}\frac{sup_X|s^N|}{\sqrt{\mathrm{log}N}}C\text{almost surely},$$
which gives the first statement of Theorem 2.3.
We now show (83), following an approach inspired by Nonnenmacher and Voros \[NoVo\]. Recalling (75), we note that
$$\mathrm{\Pi }_N(x,y)=\underset{j=1}{\overset{d_N}{}}S_j^N(x)\overline{S_j^N(y)}=\stackrel{~}{\mathrm{\Phi }}_N(x),\stackrel{~}{\mathrm{\Phi }}_N(y).$$
(84)
Let $`s^N=_{j=1}^{d_N}c_jS_j^N`$ ($`|c_j|^2=1`$) denote a random element of $`SH^0(M,L^N)=S_N^2(X)`$, and write $`c=(c_1,\mathrm{},c_{d_N})`$. Recall that
$$s^N(x)=_X\mathrm{\Pi }_N(x,y)s^N(y)𝑑y=\underset{j=1}{\overset{d_N}{}}c_jS_j^N(x)=c\stackrel{~}{\mathrm{\Phi }}_N(x).$$
(85)
Thus
$$|s^N(x)|=\stackrel{~}{\mathrm{\Phi }}_N(x)\mathrm{cos}\theta _x,\text{where }\mathrm{cos}\theta _x=\frac{\left|c\stackrel{~}{\mathrm{\Phi }}_N(x)\right|}{\stackrel{~}{\mathrm{\Phi }}_N(x)}.$$
(86)
(Note that $`\theta _x`$ can be interpreted as the distance in $`^{d_N1}`$ between $`[\overline{c}]`$ and $`\stackrel{~}{\mathrm{\Phi }}_N(x)`$.)
We have by Theorem 3.1(a),
$$\stackrel{~}{\mathrm{\Phi }}_N(x)=\mathrm{\Pi }_N(x,x)^{\frac{1}{2}}=N^{m/2}+O(N^{m/21})=(1+\epsilon _N)N^{m/2},$$
(87)
where $`\epsilon _N`$ denotes a term satisfying the uniform estimate
$$\underset{xX}{sup}|\epsilon _N(x)|O\left(\frac{1}{N}\right).$$
(88)
Now fix a point $`xX`$. By Lemma 4.1,
$`\nu _N\{s^N:\mathrm{cos}\theta _xCN^{m/2}\sqrt{\mathrm{log}N}\}`$ $`=`$ $`\left(1{\displaystyle \frac{C^2\mathrm{log}N}{N^m}}\right)^{d_N1}`$ (89)
$``$ $`\left(e^{\frac{C^2\mathrm{log}N}{N^m}}\right)^{d_N1}=N^{C^2N^m(d_N1)}.`$
We can cover $`M`$ by a collection of $`k_N`$ balls $`B(z^j)`$ of radius
$$R_N:=\frac{1}{N^{\frac{m+1}{2}}}$$
(90)
centered at points $`z^1,\mathrm{},z^{k_N}`$, where
$$k_NO(R^{2m})O(N^{m(m+1)}).$$
By (89), we have
$$\nu _N\{s^NSH_J^0(M,L^N):\underset{j}{\mathrm{max}}\mathrm{cos}\theta _{x^j}CN^{m/2}\sqrt{\mathrm{log}N}\}k_NN^{C^2N^m(d_N1)},$$
(91)
where $`x^j`$ denotes a point in $`X`$ lying above $`z^j`$.
Equation (91) together with (86)–(87) implies (by the argument below) that the desired sup-norm estimate holds at the centers of the small balls with high probability. To complete the proof of (83), we first need to extend (91) to points within the balls. To do this, we consider an arbitrary point $`w^jB(z^j)`$, and choose points $`y^jX`$ lying above the points $`w^j`$. We must estimate the distance, which we denote by $`\delta _N^j`$, between $`\mathrm{\Phi }_N(z^j)`$ and $`\mathrm{\Phi }_N(w^j)`$ in $`^{d_N1}`$. Letting $`\gamma `$ denote the geodesic in $`M`$ from $`z^j`$ to $`w^j`$, we conclude by Theorem 3.1(b) that
$`\delta _N^j`$ $``$ $`{\displaystyle _{\mathrm{\Phi }_N\gamma }}\sqrt{\omega _{FS}}={\displaystyle _\gamma }\sqrt{\mathrm{\Phi }_N^{}\omega _{FS}}\sqrt{N}{\displaystyle _\gamma }(1+\epsilon _N))\sqrt{\omega }`$ (92)
$``$ $`(1+\epsilon _N)N^{\frac{1}{2}}R_N={\displaystyle \frac{1+\epsilon _N}{N^{m/2}}}.`$
By the triangle inequality in $`^{d_N1}`$, we have $`|\theta _{x^j}\theta _{y^j}|\delta _N^j`$. Therefore by (92),
$$\mathrm{cos}\theta _{x^j}\mathrm{cos}\theta _{y^j}\delta _N^j\mathrm{cos}\theta _{y^j}\frac{1+\epsilon _N}{N^{m/2}}.$$
(93)
By (93),
$$\mathrm{cos}\theta _{y^j}\frac{(C+1)\sqrt{\mathrm{log}N}}{N^{m/2}}\mathrm{cos}\theta _{x^j}\frac{(C+1)\sqrt{\mathrm{log}N}(1+\epsilon _N)}{N^{m/2}}\frac{C\sqrt{\mathrm{log}N}}{N^{m/2}}$$
and thus
$$\begin{array}{c}\{s^NSH_J^0(M,L^N):sup\mathrm{cos}\theta (C+1)N^{m/2}\sqrt{\mathrm{log}N}\}\hfill \\ \{s^NSH_J^0(M,L^N):\mathrm{max}_j\mathrm{cos}\theta _{x^j}CN^{m/2}\sqrt{\mathrm{log}N}\}.\hfill \end{array}$$
Hence by (91),
$$\nu _N\{s^NSH_J^0(M,L^N):sup\mathrm{cos}\theta (C+1)N^{m/2}\sqrt{\mathrm{log}N}\}k_NN^{C^2N^m(d_N1)}.$$
(94)
By the Riemann-Roch formula of Boutet de Monvel - Guillemin \[BoGu, §14\] (which is a consequence of Theorem 1.1), we have the estimate for the dimensions $`d_N`$:
$$d_N=\frac{c_1(L)^m}{m!}N^m+O(N^{m1}).$$
(95)
We can also obtain (95) from Theorem 3.1(a) as follows: We note first that
$$_X\mathrm{\Pi }_N(x,x)d\mathrm{Vol}_X=_X\underset{j=1}{\overset{d_N}{}}|S_j^N(x)|^2d\mathrm{Vol}_X=d_N.$$
On the other hand, by Theorem 3.1(a),
$$_X\mathrm{\Pi }_N(x,x)d\mathrm{Vol}_X=\left[\frac{1}{\pi ^m}N^m+O(N^{m1})\right]\mathrm{Vol}(X),$$
where
$$\mathrm{Vol}(X)=\mathrm{Vol}(M)=_M\frac{1}{m!}\omega ^m=\frac{\pi ^m}{m!}c_1(L)^m.$$
Equating the above computations of the integral yields (95).
It follows from (86), (87), (94) and (95) that
$`\nu _N\{s^NSH_J^0(M,L^N):\underset{M}{sup}|s^N|(C+2)\sqrt{\mathrm{log}N}\}`$
$`k_NN^{C^2N^m(d_N1)}O\left(N^{m(m+1)\frac{C^2}{m!+1}}\right).`$
Choosing $`C=(m+1)\sqrt{m!+1}`$, we obtain (83).
Remark: An alternate proof of this estimate, which does not depend on Tian’s theorem, is given by the case $`k=0`$ of the $`𝒞^k`$ estimate in the next section.
### 4.3. The estimate $`^ks_N_{\mathrm{}}/s_N_2=O(\sqrt{N^k\mathrm{log}N})`$ almost surely
The proof of this and the other assertions of Theorem 0.1 follow the pattern of the above sup-norm estimate. First we note a consequence (Lemma 4.3) of our near-diagonal asymptotics. Recall that a differential operator on $`X`$ is horizontal if it is generated by horizontal vector fields. In particular the operators $`^k:𝒞^{\mathrm{}}(M,L^N)𝒞^{\mathrm{}}(M,L^N(T^{}M)^k)`$ are given by (vector valued) horizontal differential operators (independent of $`N`$) on $`X`$. By definition, horizontal differential operators on $`X\times X`$ are generated by the horizontal differential operators on the first and second factors. We begin with the following estimate:
###### Lemma 4.2.
Let $`P_k`$ be a horizontal differential operator of order $`k`$ on $`X\times X`$. Then
$$P_k\mathrm{\Pi }_N(x,y)|_{x=y}=O(N^{m+k/2}).$$
###### Proof.
Let $`x_0=(P_0,0)`$ be an arbitrary point of $`X`$, and choose local real coordinates $`(x_1,\mathrm{},x_{2m},\theta )`$ about $`(P_0,0)`$ as in the hypothesis of Theorem 2.3 (with $`z_q=x_q+ix_{m+q}`$). We let $`\frac{^h}{x_q}`$ denote the horizontal lift of $`\frac{}{x_q}`$ to $`X`$:
$$\frac{^h}{x_q}=\frac{}{x_q}\stackrel{~}{A}_q(x)\frac{}{\theta },\stackrel{~}{A}_q=(\alpha ,\frac{}{x_q}).$$
Since $`\frac{}{x_q}|_{x_0}`$ is assumed to be horizontal, we have $`\stackrel{~}{A}_q(P_0)=0`$.
We let $`d_q^1,d_q^2`$ denote the operator $`\frac{^h}{x_q}`$ applied to the first and second factors, respectively, on $`X\times X`$. For this result, we need only the zeroth order estimate of Theorem 2.3:
$$\mathrm{\Pi }_N(\frac{u}{\sqrt{N}},\frac{s}{N};\frac{v}{\sqrt{N}},\frac{t}{N})=N^me^{i(st)+\psi _2(u,v)}(P_0,u,v,N),$$
(96)
where $`(P_0,u,v,N)`$ denotes a term satisfying the remainder estimate of Theorem 2.3: $`(P_0,u,v,N)_{𝒞^j(\{|u|\rho ,|v|\rho \}}C_{j,\rho }`$ for $`j0,\rho >0`$, where $`C_{j,\rho }`$ is independent of the point $`P_0`$ and choice of coordinates.
Differentiating (96) and noting that $`/x_q=\sqrt{N}/u_q`$, $`/\theta =N/s`$, we have
$`d_q^1\mathrm{\Pi }_N({\displaystyle \frac{u}{\sqrt{N}}},{\displaystyle \frac{s}{N}};{\displaystyle \frac{v}{\sqrt{N}}},{\displaystyle \frac{t}{N}})`$ (97)
$`=`$ $`\sqrt{N}\left({\displaystyle \frac{}{u_q}}\sqrt{N}\stackrel{~}{A}_q(P_0+\frac{u}{\sqrt{N}}){\displaystyle \frac{}{s}}\right)\left(N^me^{i(st)+\psi _2(u,v)}\right)`$
$`=`$ $`N^{m+1/2}e^{i(st)+\psi _2(u,v)}\left\{\left[L_q(u,v)i\sqrt{N}\stackrel{~}{A}_q(\frac{u}{\sqrt{N}})\right]+{\displaystyle \frac{}{u_q}}\right\}`$
$`=`$ $`N^{m+1/2}e^{i(st)+\psi _2(u,v)}\stackrel{~}{}=O(N^{m+1/2}),`$
where $`L_q:=\frac{\psi _2}{u_q}`$ is a linear function. The same estimate holds for $`d_q^2\mathrm{\Pi }_N`$. Indeed, the above computation yields:
$$d_q^je^{i(st)+\psi _2(u,v)}(P_0,u,v,N)=\sqrt{N}e^{i(st)+\psi _2(u,v)}\stackrel{~}{}(P_0,u,v,N),$$
(98)
for $`j=1,2,q=1,\mathrm{},2m`$. The desired estimate follows by iterating (98). ∎
Remark: The assumption that $`P_k`$ is horizontal in Lemma 4.2 is necessary, since the operator $`/\theta `$ multiplies the estimate by $`N`$ instead of $`\sqrt{N}`$.
###### Lemma 4.3.
Let $`P_k`$ be a horizontal differential operator of order $`k`$ on $`X`$. Then
$$\underset{X}{sup}P_k\stackrel{~}{\mathrm{\Phi }}_N=O(N^{\frac{m+k}{2}}).$$
###### Proof.
Let $`P_k^1,P_k^2`$ denote the operator $`P_k`$ applied to the first and second factors, respectively, on $`X\times X`$. Differentiating (84) and restricting to the diagonal, we obtain
$$P_k^1\overline{P}_k^2\mathrm{\Pi }_N(x,x)=P_k\stackrel{~}{\mathrm{\Phi }}_N(x)^2.$$
(99)
The conclusion follows from (99) and Lemma 4.2 applied to the horizontal differential operator (of order $`2k`$) $`P_k^1\overline{P}_k^2`$ on $`X\times X`$. ∎
We are now ready to use the small-ball method of the previous section to show that $`^ks_N_{\mathrm{}}/s_N_^2=O(\sqrt{N^k\mathrm{log}N})`$ almost surely. It is sufficient to show that
$$\nu _N\{s^NSH_J^0(M,L^N):\underset{M}{sup}|^ks_N|>C\sqrt{N^k\mathrm{log}N}\}<O\left(\frac{1}{N^2}\right),$$
(100)
for $`C`$ sufficiently large. To verify (100), we may regard $`s_N`$ as a function on $`X`$ and replace $`^k`$ by a horizontal $`r_\theta `$-invariant differential operator of order $`k`$ on $`X`$.
As before, we let $`s^N=c_js_j^N`$ denote a random element of $`S_N^2(X)`$. By (85), we have
$$P_ks^N(x)=_XP_k^1\mathrm{\Pi }_N(x,y)s^N(y)𝑑y=\underset{j=1}{\overset{d_N}{}}c_jP_kS_j^N(x)=cP_k\stackrel{~}{\mathrm{\Phi }}_N(x).$$
(101)
We then have
$$|P_ks^N(x)|=P_k\stackrel{~}{\mathrm{\Phi }}_N(x)\mathrm{cos}\theta _x,\text{where }\mathrm{cos}\theta _x=\frac{\left|cP_k\stackrel{~}{\mathrm{\Phi }}_N(x)\right|}{P_k\stackrel{~}{\mathrm{\Phi }}_N(x)}.$$
(102)
Now fix a point $`xX`$. As before, (89) holds, and hence by Lemma 4.3 we have
$$\nu _N\{s^NS_N^2:|P_ks^N(x)|C^{}\sqrt{N^k\mathrm{log}N}\}k_NN^{C^2N^m(d_N1)},$$
(103)
where $`C^{}=Csup_{N,x}N^{(m+k)/2}|P_k\stackrel{~}{\mathrm{\Phi }}_N(x)|`$.
We again cover $`M`$ by a collection of $`k_N`$ very small balls $`B(z^j)`$ of radius $`R_N=N^{\frac{m+1}{2}}`$ and first show that the probability of the required condition holding at the centers of all the balls is small. Choosing points $`x^jX`$ lying above the centers $`z^j`$ of the balls, we then have
$$\nu _N\{s^NS_N^2:\underset{j}{\mathrm{max}}|P_ks^N(x_j)|C^{}\sqrt{N^k\mathrm{log}N}\}k_NN^{C^2N^m(d_N1)}.$$
(104)
Now suppose that $`w^j`$ is an arbitrary point in $`B(z^j)`$, and let $`y^j`$ be the point of $`X`$ above $`w^j`$ such that the horizontal lift of the geodesic from $`z^j`$ to $`w^j`$ connects $`x^j`$ and $`y^j`$. Hence by Lemma 4.3, we have
$$P_k\stackrel{~}{\mathrm{\Phi }}_N(x^j)P_k\stackrel{~}{\mathrm{\Phi }}_N(y^j)\underset{M}{sup}d^h(P_k\stackrel{~}{\mathrm{\Phi }}_N)r_N=O(N^{\frac{m+k+1}{2}})r_N=O(N^{\frac{k}{2}}).$$
(105)
It follows as before from (104) and (105) that
$`\nu _N\{s^NS_N^2:\underset{X}{sup}|P_ks^N|(C^{}+1)\sqrt{N^k\mathrm{log}N}\}`$
$`k_NN^{C^2N^m(d_N1)}O\left(N^{m(m+1)\frac{C^2}{m!+1}}\right).`$
(Here, we used the fact that $`|P_ks^N|`$ is constant on the fibers of $`\pi :XM`$.) Thus, (100) holds with $`C`$ sufficiently large.∎
### 4.4. The estimate $`\overline{}s_N_{\mathrm{}}/s_N_2=O(\sqrt{\mathrm{log}N})`$ almost surely
The proof of the $`\overline{}s_N`$ estimate follows the pattern of the above estimate. However, there is one crucial difference: we must show the following upper bound for the modulus of $`\overline{}_b\stackrel{~}{\mathrm{\Phi }}_N`$. This estimate is a factor of $`\sqrt{N}`$ better than the one for $`d^h\stackrel{~}{\mathrm{\Phi }}_N`$ arising from Lemma 4.3; the proof depends on the precise second order approximation of Theorem 2.3.
###### Lemma 4.4.
$`sup_X\overline{}_b\stackrel{~}{\mathrm{\Phi }}_N(x)O(N^{m/2}).`$
###### Proof.
Let $`x_0=(P_0,0)`$ be an arbitrary point of $`X`$, and choose preferred local coordinates $`(z_1,\mathrm{},z_m,\theta )`$ about $`(P_0,0)`$ as in the hypothesis of Theorem 2.3. We lift a local frame $`\{\overline{Z}_q^M\}`$ of the form (3) to obtain the local frame $`\{\overline{Z}_1,\mathrm{},\overline{Z}_m\}`$ for $`H^{0,1}X`$ given by
$$\overline{Z}_q=\frac{^h}{\overline{z}_q}+\underset{r=1}{\overset{m}{}}B_{qr}(z)\frac{^h}{z_r},B_{qr}(P_0)=0.$$
(106)
It suffices to show that
$$N^{m/2}|\overline{Z}_q\stackrel{~}{\mathrm{\Phi }}_N(x_0)|C,$$
(107)
where $`C`$ is a constant independent of $`x_0`$.
By Theorem 2.3, we have
$`N^m\mathrm{\Pi }_N({\displaystyle \frac{u}{\sqrt{N}}},{\displaystyle \frac{s}{N}};{\displaystyle \frac{v}{\sqrt{N}}},{\displaystyle \frac{t}{N}})`$ (108)
$`=`$ $`{\displaystyle \frac{1}{\pi ^m}}\phi _0(u,s)\overline{\phi _0(v,t)}e^{u\overline{v}}\left[1+{\displaystyle \frac{1}{\sqrt{N}}}b_1(P_0,u,v)+{\displaystyle \frac{1}{N}}R_2(P_0,u,v,N)\right],`$
where
$$\phi _0(z,\theta )=e^{i\theta |z|^2/2}.$$
(The function $`\phi _0`$ is the ‘ground state’ for the ‘annihilation operators’ $`\overline{Z}_q`$ in the Heisenberg model; see the remark in §1.2 and \[BSZ2, §1.3.2\]). In our case, $`\overline{Z}_q\phi _0`$ does not vanish as in the model case, but instead satisfies the asymptotic bound (110) below.)
Recalling (8) and (16), we have
$$\frac{^h}{\overline{z}_q}=\frac{}{\overline{z}_q}+\left[\frac{i}{2}z_qR_1^{\overline{A}_q}(z)\right]\frac{}{\theta },$$
(109)
where $`R_1^{\overline{A}_q}(z)=O(|z|^2)`$. Recalling that $`z=u/\sqrt{N}`$, $`\theta =s/N`$, we note that $`\phi _0(u,s)=e^{iN\theta N|z|^2/2}=\phi _0(z,\theta )^N`$, and thus by (109),
$$\frac{^h}{\overline{z}_q}\phi _0(u,s)=\frac{^h}{\overline{z}_q}e^{iN\theta N|z|^2/2}=iNR_1^{\overline{A}_q}(\frac{u}{\sqrt{N}})\phi _0(u,s)=(P_0,u,N)\phi _0(u,s),$$
(110)
where as before $``$ denotes a term satisfying the remainder estimate of Theorem 2.3.
We let $`Z_q^1,Z_q^2`$ denote the operator $`Z_q`$ applied to the first and second factors, respectively, on $`X\times X`$; we similarly let $`d_q^1,d_q^2`$ denote the operator $`\frac{^h}{z_q}`$ applied to the factors of $`X\times X`$. Equation (99) tells us that
$$\overline{Z}_q\stackrel{~}{\mathrm{\Phi }}_N(x)^2=\overline{Z}_q^1Z_q^2\mathrm{\Pi }_N(x,x).$$
(111)
By (106),
$$\overline{Z}_q^1Z_q^2=\left(\overline{d_q^1}+\underset{r=1}{\overset{m}{}}B_{qr}(z)d_r^1\right)\left(d_q^2+\underset{\rho =1}{\overset{m}{}}\overline{B}_{q\rho }(w)\overline{d_\rho ^2}\right),$$
(112)
where we recall that $`B_{qr}(P_0)=0`$.
Differentiating (108), again noting that $`/z_q=\sqrt{N}/u_q`$, $`/w_q=\sqrt{N}/v_q`$ and using (110), we obtain
$`N^m\left(\overline{d_q^1}d_q^2\mathrm{\Pi }_N\right)({\displaystyle \frac{u}{\sqrt{N}}},{\displaystyle \frac{s}{N}};{\displaystyle \frac{v}{\sqrt{N}}},{\displaystyle \frac{t}{N}})`$ (113)
$`=`$ $`{\displaystyle \frac{1}{\pi ^m}}\phi _0(u,s)\overline{\phi _0(v,t)}e^{u\overline{v}}\left[\sqrt{N}{\displaystyle \frac{^2}{\overline{u}_qv_q}}b_1+\stackrel{~}{}\right].`$
Since $`b_1`$ has no terms that are quadratic in $`(u,\overline{u},v,\overline{v})`$, it follows from (112)–(113) that
$$N^m\left|\overline{Z}_q^1Z_q^2\mathrm{\Pi }_N(P_0,0;P_0,0)\right|=N^m\left|\overline{d_q^1}d_q^2\mathrm{\Pi }_N(P_0,0;P_0,0)\right|=\frac{1}{\pi ^m}\left|\stackrel{~}{}(P_0,0,0,N)\right|O(1).$$
(114)
The desired estimate (107) now follows immediately from (111) and (114). ∎
By covering $`M`$ with small balls and repeating the argument of the previous section, using Lemma 4.4, we conclude that
$$\nu _N\{s^NSH_J^0(M,L^N):\underset{M}{sup}|\overline{}s_N|>C\sqrt{\mathrm{log}N}\}<O\left(\frac{1}{N^2}\right).$$
(115)
Thus $`\overline{}s_N_{\mathrm{}}/s_N_2=O(\sqrt{\mathrm{log}N})`$ almost surely.
### 4.5. The estimate $`^k\overline{}s_N_{\mathrm{}}/s_N_2=O(\sqrt{N^k\mathrm{log}N})`$ almost surely
To obtain this final estimate of Theorem 0.1, it suffices to verify the probability estimate
$$\nu _N\{s^NSH_J^0(M,L^N):\underset{M}{sup}|^k\overline{}s_N|>C\sqrt{N^k\mathrm{log}N}\}<O\left(\frac{1}{N^2}\right).$$
(116)
Equation (116) follows by again repeating the argument of §4.3, using the following lemma.
###### Lemma 4.5.
Let $`P_k`$ be a horizontal differential operator of order $`k`$ on $`X`$ ($`k0`$). Then
$$\underset{X}{sup}|P_k\overline{}_b\stackrel{~}{\mathrm{\Phi }}_N|=O(N^{\frac{m+k}{2}}).$$
###### Proof.
It suffices to show that
$$\underset{U}{sup}|P_k\overline{Z}_q^k\stackrel{~}{\mathrm{\Phi }}_N|=O(N^{\frac{m+k}{2}})$$
(117)
for a local frame $`\{\overline{Z}_q\}`$ of $`T^{0,1}M`$ over $`U`$. As before, we have
$$P_k^1\overline{P}_k^2\overline{Z}_q^1Z_q^1\mathrm{\Pi }_N(x,x)=\left|P_k\overline{Z}_q^h\stackrel{~}{\mathrm{\Phi }}_N(x)\right|^2.$$
(118)
We claim that
$$N^m\overline{Z}_q^1Z_q^2\mathrm{\Pi }_N=\frac{1}{\pi ^m}e^{i(st)+\psi _2(u,v)}\left[\sqrt{N}\frac{^2}{\overline{u}_qv_q}b_1+(P_0,u,v,N)\right].$$
(119)
To obtain the estimate (119), we recall from (112) in the proof of Lemma 4.4 that
$$\overline{Z}_q^1Z_q^2=\overline{d_q^1}d_q^2+\underset{\rho =1}{\overset{m}{}}\overline{B}_{q\rho }(w)\overline{d_q^1}\overline{d_\rho ^2}+\underset{r=1}{\overset{m}{}}B_{qr}(z)d_r^1d_q^2+\underset{r,\rho }{}B_{qr}(z)\overline{B}_{q\rho }(w)d_r^1\overline{d_\rho ^2}.$$
(120)
Equation (113) says that the first term of $`N^m\overline{Z}_q^1Z_q^2\mathrm{\Pi }_N`$ coming from the expansion (120) satisfies the estimate of (119). To obtain the estimate for the second term, we compute:
$`N^m\overline{d_\rho ^2}\mathrm{\Pi }_N({\displaystyle \frac{u}{\sqrt{N}}},{\displaystyle \frac{s}{N}};{\displaystyle \frac{v}{\sqrt{N}}},{\displaystyle \frac{t}{N}})={\displaystyle \frac{\sqrt{N}}{\pi ^m}}e^{i(st)+\psi _2(u,v)}`$ (121)
$`\left(\left[{\displaystyle \frac{\psi _2}{\overline{v}_\rho }}+i\sqrt{N}A_\rho ({\displaystyle \frac{v}{\sqrt{N}}})\right]\left[1+{\displaystyle \frac{1}{\sqrt{N}}}b_1+{\displaystyle \frac{1}{N}}R_2\right]+{\displaystyle \frac{1}{\sqrt{N}}}{\displaystyle \frac{b_1}{\overline{v}_\rho }}+{\displaystyle \frac{1}{N}}{\displaystyle \frac{R_2}{\overline{v}_\rho }}\right)`$
$`=`$ $`{\displaystyle \frac{\sqrt{N}}{\pi ^m}}\phi _0(u,s)\overline{\phi _0(v,t)}e^{u\overline{v}}\left[{\displaystyle \frac{\psi _2}{\overline{v}_\rho }}+L_\rho (v)+{\displaystyle \frac{1}{\sqrt{N}}}\stackrel{~}{}\right],`$
where $`L_\rho `$ is a linear function. Since $`^2\psi _2/\overline{u}_q\overline{v}_\rho 0`$, it then follows that
$$N^m\overline{d_q^1}\overline{d_\rho ^2}\mathrm{\Pi }_N=\frac{\sqrt{N}}{\pi ^m}e^{i(st)+\psi _2(u,v)}\frac{}{\overline{u}_q}\stackrel{~}{}(P_0,u,v,N).$$
(122)
The estimate (119) for the second term follows from (122), using the fact that $`B_{q\rho }(\frac{v}{\sqrt{N}})=\frac{1}{\sqrt{N}}L_{q\rho }(v)+\mathrm{}`$. The proofs of the estimate for the third and fourth terms are similar.
The desired estimate (117) follows as before from (118), (119), and (98), using the fact that $`\frac{^2}{\overline{u}_qv_q}b_1`$ is linear. ∎
## 5. The joint probability distribution
In this section, we shall use Theorem 2.3 and the methods of of \[BSZ2\] to prove Theorem 0.2 and its analogue for Gaussian measures (Theorem 5.4), which say that the joint probability distributions on almost complex symplectic manifolds have the same universal scaling limit as in the complex case.
### 5.1. Generalized Gaussians
Recall that a Gaussian measure on $`^n`$ is a measure of the form
$$\gamma _\mathrm{\Delta }=\frac{e^{\frac{1}{2}\mathrm{\Delta }^1x,x}}{(2\pi )^{n/2}\sqrt{det\mathrm{\Delta }}}dx_1\mathrm{}dx_n,$$
where $`\mathrm{\Delta }`$ is a positive definite symmetric $`n\times n`$ matrix. The matrix $`\mathrm{\Delta }`$ gives the second moments of $`\gamma _\mathrm{\Delta }`$:
$$x_jx_k_{\gamma _\mathrm{\Delta }}=\mathrm{\Delta }_{jk}.$$
(123)
This Gaussian measure is also characterized by its Fourier transform
$$\widehat{\gamma _\mathrm{\Delta }}(t_1,\mathrm{},t_n)=e^{\frac{1}{2}{\scriptscriptstyle \mathrm{\Delta }_{jk}t_jt_k}}.$$
(124)
If we let $`\mathrm{\Delta }`$ be the $`n\times n`$ identity matrix, we obtain the standard Gaussian measure on $`^n`$,
$$\gamma _n:=\frac{1}{(2\pi )^{n/2}}e^{\frac{1}{2}|x|^2}dx_1\mathrm{}dx_n,$$
with the property that the $`x_j`$ are independent Gaussian variables with mean 0 and variance 1. Hence
$$x^2_{\gamma _n}=\underset{j=1}{\overset{n}{}}x_j^2_{\gamma _n}=n.$$
Since we wish to put Gaussian measures on the spaces $`H_J^0(M,L^N)`$ with rapidly growing dimensions, it is useful to consider the normalized standard Gaussians
$$\stackrel{~}{\gamma }_n:=k_ne^{\frac{n}{2}|x|^2}dx_1\mathrm{}dx_n,k_n=\left(\frac{n}{2\pi }\right)^{n/2},$$
which have the property that
$$x^2_{\stackrel{~}{\gamma }_n}=1.$$
The push-forward of a Gaussian measure by a surjective linear map is also Gaussian. In the next section, we shall push forward Gaussian measures (on the spaces $`H_J^0(M,L^N)`$) by linear maps that are sometimes not surjective. Since these non-surjective push-forwards are singular measures, we need to consider the case where $`\mathrm{\Delta }`$ is positive semi-definite. In this case, we use (124) to define a measure $`\gamma _\mathrm{\Delta }`$, which we call a generalized Gaussian. If $`\mathrm{\Delta }`$ has null eigenvalues, then the generalized Gaussian $`\gamma _\mathrm{\Delta }`$ is a Gaussian measure on the subspace $`\mathrm{\Lambda }_+^n`$ spanned by the positive eigenvectors. (Precisely, $`\gamma _\mathrm{\Delta }=\iota _{}\gamma _{\mathrm{\Delta }|\mathrm{\Lambda }_+}`$, where $`\iota :\mathrm{\Lambda }_+^n`$ is the inclusion. For the completely degenerate case $`\mathrm{\Delta }=0`$, we have $`\gamma _\mathrm{\Delta }=\delta _0`$.) Of course, (123) holds for semi-positive $`\mathrm{\Delta }`$. One useful property of generalized Gaussians is that the push-forward by a (not necessarily surjective) linear map $`T:^n^m`$ of a generalized Gaussian $`\gamma _\mathrm{\Delta }`$ on $`^n`$ is a generalized Gaussian on $`^m`$:
$$T_{}\gamma _\mathrm{\Delta }=\gamma _{T\mathrm{\Delta }T^{}}$$
(125)
Another useful property of generalized Gaussians is the following fact:
###### Lemma 5.1.
The map $`\mathrm{\Delta }\gamma _\mathrm{\Delta }`$ is a continuous map from the positive semi-definite matrices to the space of positive measures on $`^n`$ (with the weak topology).
###### Proof.
Suppose that $`\mathrm{\Delta }^N\mathrm{\Delta }^0`$. We must show that $`(\mathrm{\Delta }^N,\phi )(\mathrm{\Delta }^0,\phi )`$ for a compactly supported test function $`\phi `$. We can assume that $`\phi `$ is $`𝒞^{\mathrm{}}`$. It then follows from (124) that
$$(\gamma _{\mathrm{\Delta }^N},\phi )=(\widehat{\gamma _{\mathrm{\Delta }^N}},\widehat{\phi })(\widehat{\gamma _{\mathrm{\Delta }^0}},\widehat{\phi })=(\gamma _{\mathrm{\Delta }^0},\phi ).$$
We shall use the following general result relating spherical measures to Gaussian measures in order to prove Theorem 0.2 on asymptotics of the joint probability distributions for $`SH_J^0(M,L^N)`$.
###### Lemma 5.2.
Let $`T_N:^{d_N}R^k`$, $`N=1,2,`$, be a sequence of linear maps, where $`d_N\mathrm{}`$. Suppose that $`\frac{1}{d_N}T_NT_N^{}\mathrm{\Delta }`$. Then $`T_N\nu _{d_N}\gamma _\mathrm{\Delta }`$.
###### Proof.
Let $`V_N`$ be a $`k`$-dimensional subspace of $`^{d_N}`$ such that $`V_N^{}\mathrm{ker}T^N`$, and let $`p_N:^{d_N}V_N`$ denote the orthogonal projection. We decompose $`T_N=B_NA_N`$, where $`A_N=d_N^{1/2}p_N:^{d_N}V_N`$, and $`B_N=d_N^{1/2}T_N|_{V_N}:V_N\stackrel{}{}^k`$. Write
$$A_N\nu _{d_N}=\alpha _N,T_N\nu _{d_N}=B_N\alpha _N=\beta _N.$$
We easily see that (abbreviating $`d=d_N`$)
$$\alpha _N=A_N\nu _d=\psi _ddx,\psi _d=\{\begin{array}{cc}\frac{\sigma _{dk}}{\sigma _dd^k}[1\frac{1}{d}|x|^2]^{(dk2)/2}& \mathrm{for}|x|<\sqrt{d}\hfill \\ 0& \mathrm{otherwise}\hfill \end{array},$$
(126)
where $`dx`$ denotes Lebesgue measure on $`V_N`$, and $`\sigma _n=\mathrm{Vol}(S^{n1})=\frac{2\pi ^{n/2}}{\mathrm{\Gamma }(n/2)}`$. (The case $`k=1,d=3`$ of (126) is Archimedes’ formula \[Arc\].) Since $`[1|x|^2/d]^{(dk2)/2}e^{|x|^2/2}`$ uniformly on compacta and $`\frac{\sigma _{dk}}{\sigma _dd^k}\frac{1}{(2\pi )^{k/2}}`$, we conclude that $`\alpha _N\gamma _k`$. (This is the Poincaré-Borel Theorem; see Corollary 5.3 below.) Furthermore,
$$\left(1\frac{1}{d}|x|^2\right)^{(dk2)/2}\mathrm{exp}\left(\frac{dk2}{2d}|x|^2\right)e^{\frac{k+2}{2}}e^{\frac{1}{2}|x|^2}\mathrm{for}dk+2,|x|\sqrt{d},$$
and hence
$$\psi _{d_N}(x)C_ke^{|x|^2/2}.$$
(127)
Now let $`\phi `$ be a compactly supported continuous test function on $`^k`$. We must show that
$$\phi 𝑑\beta _N\phi 𝑑\gamma _\mathrm{\Delta }.$$
(128)
Suppose on the contrary that (128) does not hold. After passing to a subsequence, we may assume that $`\phi 𝑑\beta _Nc\phi 𝑑\gamma _\mathrm{\Delta }`$. Since the eigenvalues of $`B_N`$ are bounded, we can assume (again taking a subsequence) that $`B_NB_0`$, where
$$B_0B_0^{}=\underset{N\mathrm{}}{lim}B_NB_N^{}=\underset{N\mathrm{}}{lim}\frac{1}{d_N}T_NT_N^{}=\mathrm{\Delta }.$$
Hence,
$$_^k\phi 𝑑\beta _N=_{V_N}\phi (B_Nx)\psi _{d_N}(x)𝑑x_{V_N}\phi (B_0x)\frac{e^{|x|^2/2}}{(2\pi )^{k/2}}𝑑x=_{V_N}\phi (B_0x)𝑑\gamma _k(x),$$
where the limit holds by dominated convergence, using (127). By (125), we have $`B_0\gamma _k=\gamma _{B_0B_0^{}}=\gamma _\mathrm{\Delta }`$, and hence
$$_{V_N}\phi (B_0x)𝑑\gamma _k(x)=_^k\phi 𝑑\gamma _\mathrm{\Delta }.$$
Thus (128) holds for the subsequence, giving a contradiction. ∎
We note that the above proof began by establishing the Poincaré-Borel Theorem (which is a special case of the of Lemma 5.2):
###### Corollary 5.3.
(Poincaré-Borel) Let $`P_d:^d^k`$ be given by $`P_d(x)=\sqrt{d}(x_1,\mathrm{},x_k)`$. Then
$$P_d\nu _d\gamma _k.$$
By a generalized complex Gaussian measure on $`^n`$, we mean a generalized Gaussian measure $`\gamma _\mathrm{\Delta }^c`$ on $`^nR^{2n}`$ with moments
$$z_j_\gamma =0,z_jz_k_\gamma =0,z_j\overline{z}_k_\gamma =\mathrm{\Delta }_{jk},1j,kn,$$
where $`\mathrm{\Delta }=(\mathrm{\Delta }_{jk})`$ is an $`n\times n`$ positive semi-definite hermitian matrix; i.e. $`\gamma _\mathrm{\Delta }^c=\gamma _{\frac{1}{2}\mathrm{\Delta }^c}`$, where $`\mathrm{\Delta }^c`$ is the $`2n\times 2n`$ real symmetric matrix of the inner product on $`R^{2n}`$ induced by $`\mathrm{\Delta }`$. As we are interested here in complex Gaussians, we shall henceforth drop the ‘$`c`$’ and write $`\gamma _\mathrm{\Delta }^c=\gamma _\mathrm{\Delta }`$. In particular, if $`\mathrm{\Delta }`$ is a strictly positive hermitian matrix, then
$$\gamma _\mathrm{\Delta }=\frac{e^{\mathrm{\Delta }^1z,\overline{z}}}{\pi ^ndet\mathrm{\Delta }}d(z),$$
where $``$ denotes Lebesgue measure on $`^n`$.
### 5.2. Proof of Theorem 0.2
We return to our complex Hermitian line bundle $`(L,h)`$ on a compact almost complex $`2m`$-dimensional symplectic manifold $`M`$ with symplectic form $`\omega =\frac{i}{2}\mathrm{\Theta }_L`$, where $`\mathrm{\Theta }_L`$ is the curvature of $`L`$ with respect to a connection $``$. We now describe the $`n`$-point joint distribution arising from our probability space $`(SH_J^0(M,L^N),\nu _N)`$. Recalling (17), we have the Hermitian inner product on $`H_J^0(M,L^N)`$:
$$s_1,s_2=_Mh^N(s_1,s_2)\frac{1}{m!}\omega ^m(s_1,s_2H_J^0(M,L^N)),$$
and we write $`s_2=s,s^{1/2}`$. Recall that $`SH_J^0(M,L^N)`$ denotes the unit sphere $`\{s=1\}`$ in $`H_J^0(M,L^N)`$ and $`\nu _N`$ denotes its Haar probability measure.
We let $`J^1(M,L^N)`$ denote the space of 1-jets of sections of $`L^N`$. Recall that we have the exact sequence of vector bundles
$$0T_M^{}L^N\stackrel{\iota }{}J^1(M,L^N)\stackrel{\rho }{}L^N0.$$
(129)
We consider the jet maps
$$J_z^1:H_J^0(M,L^N)J^1(M,V)_z,J_z^1s=\text{the 1-jet of}s\text{at}z,\text{for}zM.$$
The covariant derivative $`:J^1(M,L^N)T_M^{}L^N`$ provides a splitting of (129) and an isomorphism
$$(\rho ,):J^1(M,L^N)\stackrel{}{}L^N(T_M^{}L^N).$$
(130)
Definition: The $`n`$-point joint probability distribution at points $`P^1,\mathrm{},P^n`$ of $`M`$ is the probability measure
$$𝐃_{(P^1,\mathrm{},P^n)}^N:=(J_{P^1}^1\mathrm{}J_{P^n}^1)_{}\nu _N$$
(131)
on the space $`J^1(M,L^N)_{P^1}\mathrm{}J^1(M,L^N)_{P^n}`$.
Since we are interested in the scaling limit of $`𝐃^N`$, we need to describe this measure more explicitly: Suppose that $`P^1,\mathrm{},P^n`$ lie in a coordinate neighborhood of a point $`P_0M`$ and choose preferred coordinates $`(z_1,\mathrm{},z_m)`$ and a preferred frame $`e_L`$ at $`P_0`$. We let $`z_1^p,\mathrm{},z_m^p`$ denote the coordinates of the point $`P^p`$ ($`1pn`$), and we write $`z^p=(z_1^p,\mathrm{},z_m^p)`$. (The coordinates of $`P_0`$ are $`0`$.) We consider the $`n(2m+1)`$ complex-valued random variables $`x^p,\xi _q^p`$ ($`1pn,1q2m`$) on $`S_N^2(X)SH_J^0(M,L^N)`$ given by
$$x^p(s)=s(z^p,0),$$
(132)
$$\xi _q^p(s)=\frac{1}{\sqrt{N}}\frac{^hs}{z_q}(z^p),\xi _{m+q}^p(s)=\frac{1}{\sqrt{N}}\frac{^hs}{\overline{z}_q}(z^p)(1qm),$$
(133)
for $`sSH_J^0(M,L^N)`$.
We now write
$$x=(x^1,\mathrm{},x^p),\xi =(\xi _q^p)_{1pn,1q2m},z=(z^1,\mathrm{},z^n).$$
Using (130) and the variables $`x^p,\xi _q^p`$ to make the identification
$$J^1(M,L^N)_{P^1}\mathrm{}J^1(M,L^N)_{P^n}^{n(2m+1)},$$
(134)
we can write
$$𝐃_z^N=D^N(x,\xi ;z)dxd\xi ,$$
where $`dxd\xi `$ denotes Lebesgue measure on $`^{n(2m+1)}`$.
Before proving Theorem 0.2 on the scaling limit of $`𝐃_z^N`$, we state and prove a corresponding result replacing $`(SH_J^0(M,L^N),\nu _N)`$ with the essentially equivalent Gaussian space $`H_J^0(M,L^N)`$ with the normalized standard Gaussian measure
$$\mu _N:=\stackrel{~}{\gamma }_{2d_N}=k_{2d_N}e^{d_N|c|^2}d(c),s=\underset{j=1}{\overset{d_N}{}}c_jS_j^N,$$
(135)
where $`\{S_j^N\}`$ is an orthonormal basis for $`H_J^0(M,L^N)`$. Recall that this Gaussian is characterized by the property that the $`2d_N`$ real variables $`\mathrm{}c_j,\mathrm{}c_j`$ ($`j=1,\mathrm{},d_N`$) are independent random variables with mean 0 and variance $`1/2d_N`$; i.e.,
$$c_j_{\mu _N}=0,c_jc_k_{\mu _N}=0,c_j\overline{c}_k_{\mu _N}=\frac{1}{d_N}\delta _{jk}.$$
(136)
Our normalization guarantees that the variance of $`s_2`$ is 1:
$$s_2^2_{\mu _N}=1.$$
We then consider the Gaussian joint probability distribution
$$\stackrel{~}{𝐃}_{(P^1,\mathrm{},P^n)}^N=\stackrel{~}{D}^N(x,\xi ;z)dxd\xi =(J_{P^1}^1\mathrm{}J_{P^n}^1)_{}\mu _N.$$
(137)
Since $`\mu _N`$ is Gaussian and the map $`J_{P^1}^1\mathrm{}J_{P^n}^1`$ is linear, it follows that the joint probability distribution is a generalized Gaussian measure of the form
$$D^N(x,\xi ;z)dxd\xi =\gamma _{\mathrm{\Delta }^N(z)}.$$
(138)
We shall see below that the covariance matrix $`\mathrm{\Delta }^N(z)`$ is given in terms of the Szegö kernel.
We have the following alternate form of Theorem 0.2:
###### Theorem 5.4.
Let $`L,M,\omega `$ be as above and let $`\{z_j\}`$ be preferred coordinates centered at a point $`P_0M`$. Then
$$\stackrel{~}{𝐃}_{(z^1/\sqrt{N},\mathrm{},z^n/\sqrt{N})}^N𝐃_{(z^1,\mathrm{},n^n)}^{\mathrm{}}$$
where $`𝐃_{(z^1,\mathrm{},z^n)}^{\mathrm{}}`$ is the universal Gaussian measure (supported on the holomorphic 1-jets) of Theorem 0.2.
###### Proof.
The covariance matrix $`\mathrm{\Delta }^N(z)`$ in (138) is a positive semi-definite $`n(2m+1)\times n(2m+1)`$ hermitian matrix. If the map $`J_{z^1}^1\mathrm{}J_{z^n}^1`$ is surjective, then $`\mathrm{\Delta }^N(z)`$ is strictly positive definite and $`\stackrel{~}{D}^N(x,\xi ;z)`$ is a smooth function. On the other hand, if the map is not surjective, then $`\stackrel{~}{D}^N(x,\xi ;z)`$ is a distribution supported on a linear subspace. For example, in the integrable holomorphic case, $`\stackrel{~}{D}^N(x,\xi ;z)`$ is supported on the holomorphic jets, as follows from the discussion below.
By (123), we have
$`\mathrm{\Delta }^N(z)=\left(\begin{array}{cc}A& B\\ B^{}& C\end{array}\right),`$ (141)
$`A=\left(A_p^{}^p\right)=x^p\overline{x}^p^{}_{\mu _N},B=\left(B_{p^{}q^{}}^p\right)=x^p\overline{\xi }_q^{}^p^{}_{\mu _N},C=\left(C_{p^{}q^{}}^{pq}\right)=\xi _q^p\overline{\xi }_q^{}^p^{}_{\mu _N},`$ (142)
$`p,p^{}=1,\mathrm{},n,q,q^{}=1,\mathrm{},2m.`$
(We note that $`A,B,C`$ are $`n\times n,n\times 2mn,2mn\times 2mn`$ matrices, respectively; $`p,q`$ index the rows, and $`p^{},q^{}`$ index the columns.)
We now describe the the entries of the matrix $`\mathrm{\Delta }^N`$ in terms of the Szegö kernel. We have by (136) and (142), writing $`s=_{j=1}^{d_N}c_jS_j^N`$,
$$A_p^{}^p=x^p\overline{x}^p^{}_{\mu _N}=\underset{j,k=1}{\overset{d_N}{}}c_j\overline{c}_k_{\mu _N}S_j^N(z^p,0)\overline{S_k^N(z^p^{},0)}=\frac{1}{d_N}\mathrm{\Pi }_N^{P_0}(z^p,0;z^p^{},0).$$
(143)
We need some more notation to describe the matrices $`B`$ and $`C`$: Write
$$_q=\frac{1}{\sqrt{N}}\frac{^h}{z_q},_{m+q}=\frac{1}{\sqrt{N}}\frac{^h}{\overline{z}_q},1qm.$$
As in §4.1, we let $`_q^1`$, resp. $`_q^2`$, denote the differential operator on $`X\times X`$ given by applying $`_q`$ to the first, resp. second, factor ($`1q2m`$). By differentiating (143), we obtain
$`B_{p^{}q^{}}^p`$ $`=`$ $`{\displaystyle \frac{1}{d_N}}\overline{}_q^{}^2\mathrm{\Pi }_N^{P_0}(z^p,0;z^p^{},0),`$ (144)
$`C_{p^{}q^{}}^{pq}`$ $`=`$ $`{\displaystyle \frac{1}{d_N}}_q^1\overline{}_q^{}^2\mathrm{\Pi }_N^{P_0}(z^p,0;z^p^{},0).`$ (145)
It follows from (143)–(145) and Theorem 2.3, recalling (16) and (95), that
$$\mathrm{\Delta }^N(\frac{z}{\sqrt{N}})\mathrm{\Delta }^{\mathrm{}}(z)=\frac{m!}{c_1(L)^m}\left(\begin{array}{cc}A^{\mathrm{}}(z)& B^{\mathrm{}}(z)\\ B^{\mathrm{}}(z)^{}& C^{\mathrm{}}(z)\end{array}\right)$$
(146)
uniformly, where
$`A^{\mathrm{}}(z)_p^{}^p`$ $`=`$ $`\mathrm{\Pi }_1^𝐇(z^p,0;z^p^{},0)={\displaystyle \frac{1}{\pi ^m}}e^{\psi _2(z^p,z^p^{})},`$
$`B^{\mathrm{}}(z)_{p^{}q^{}}^p`$ $`=`$ $`\{\begin{array}{cc}\frac{1}{\pi ^m}(z_q^{}w_q^{})e^{\psi _2(z^p,z^p^{})}\hfill & \text{for}1qm\hfill \\ 0\hfill & \text{for}m+1q2m\hfill \end{array},`$
$`C^{\mathrm{}}(z)_{p^{}q^{}}^{pq}`$ $`=`$ $`\{\begin{array}{cc}\frac{1}{\pi ^m}(\overline{w}_q\overline{z}_q)(z_q^{}w_q^{})e^{\psi _2(z^p,z^p^{})}\hfill & \text{for}1q,q^{}m\hfill \\ 0\hfill & \text{for}\mathrm{max}(q,q^{})m+1\hfill \end{array}.`$
(Recall that $`\psi _2(u,v)=u\overline{v}\frac{1}{2}(|u|^2+|v|^2)`$.) In other words, the coefficients of $`\mathrm{\Delta }^{\mathrm{}}(z)`$ corresponding to the anti-holomorphic directions vanish, while the coefficients corresponding to the holomorphic directions are given by the Szegö kernel $`\mathrm{\Pi }_1^𝐇`$ for the reduced Heisenberg group and its covariant derivatives.
Finally, we apply Lemma 5.1 to (138) and conclude that
$$\stackrel{~}{𝐃}_{z/\sqrt{N}}^N=\gamma _{\mathrm{\Delta }^N(z/\sqrt{N})}\gamma _{\mathrm{\Delta }^{\mathrm{}}(z)}.$$
Thus Theorem 5.4 holds with $`𝐃_z^{\mathrm{}}=\gamma _{\mathrm{\Delta }^{\mathrm{}}(z)}`$. ∎
Proof of Theorem 0.2: The proof is similar to that of Theorem 5.4. This time we define
$$\mathrm{\Delta }^N=\frac{1}{d_N}𝒥_N𝒥_N^{}:H^0(M,L^N)^{n(2m+1)},$$
where $`𝒥_N=J_{P^1}^1\mathrm{}J_{P^n}^1`$ under the identification (134). We see immediately that $`\mathrm{\Delta }^N`$ is given by (143)–(145) and the conclusion follows from Lemma 5.2 and (146). ∎
Remark: There are other similar ways to define the joint probability distribution that have the same universal scaling limits. One of these is to use the (un-normalized) standard Gaussian measure $`\gamma _{2d_N}`$ on $`H_J^0(M,L^N)`$ in place of the normalized Gaussian $`\mu _N`$ in Theorem 5.4 to obtain joint densities $`D_\mathrm{\#}^N(x,\xi ;z)=D^N(\frac{x}{N^{m/2}},\frac{\xi }{N^{m/2}};z)`$. Then we would have instead
$$D_\mathrm{\#}^N(N^{m/2}x,N^{m/2}\xi ;N^{1/2}z)dxd\xi \gamma _{\mathrm{\Delta }^{\mathrm{}}(z)}.$$
Another similar result is to let $`\lambda _N`$ denote normalized Lebesgue measure on the unit ball $`\{s1\}`$ in $`H_J^0(M,L^N)`$ and to let $`\widehat{𝐃}_z^N=𝒥_N\lambda _N`$. By a similar argument as above, we also have $`\widehat{𝐃}_{z/\sqrt{N}}^N\gamma _{\mathrm{\Delta }^{\mathrm{}}(z)}`$. |
warning/0001/cond-mat0001420.html | ar5iv | text | # Optimizing Traffic in Virtual and Real Space
## Virtual Traffic in the World Wide Web
Anyone who has browsed the World Wide Web has probably discovered the following strategy: whenever a web page takes too long to appear, it is useful to press the reload button. Very often, the web page then appears instantly. This motivates the implementation of a similar but automated strategy for the frequent “web crawls” that many Internet search engines depend on. In order to ensure up-to-date indexes, it is important to perform these crawls quickly. More generally, from an electronic commerce perspective, it is also very valuable to optimize the speed and the variance in the speed of transactions, automated or not, especially when the cost of performing those transactions is taken into account. Again, restart strategies may provide measurable benefits for the user.
The histogram in Figure 1 shows the variance associated with the download times for the text on the main page of over 40,000 web sites. Based on such observations, an economics-based strategy has recently been proposed for quantitatively managing the time of executing electronic transactions . It exploits an analogy with the modern theory of financial portfolio management by associating cost with the time it takes to complete the transaction and taking into account the “risk” given by the standard deviation of that time. Before, such a strategy has already been successfully applied to the numerical solution of hard computational problems .
In modern portfolio theory, risk averse investors may prefer to hold assets from which they expect a lower return if they are compensated for the lower return with a lower level of risk exposure. Furthermore, it is a non-trivial result of portfolio theory that simple diversification can yield portfolios of assets which have higher expected return as well as lower risk. In the case of latencies (waiting times) on the Internet, thinking of different restart strategies is analogous to asset diversification: there is an efficient trade-off between the average time a request will take and the standard deviation of that time (“risk”).
Consider a situation in which data have been requested but not received (downloaded) for some time. This time can be very long in cases where the latency distribution has a long tail. One is then faced with the choice to either continue to wait for the data, to send out another request or, if the network protocols allow, to cancel the original request before sending out another. For simplicity, we consider the case in which it is possible to cancel the original request before sending out another one after waiting for a time period of duration $`\tau `$. If $`p(t)`$ denotes the probability distribution for the download time without restart, the probability $`P(t)`$ that a page has been successfully downloaded in time less than $`t`$ is given by
$$P(t)=\{\begin{array}{cc}p(t)\hfill & \text{ if }t\tau ,\hfill \\ & \\ \left[1_0^\tau 𝑑tp(t)\right]P(t\tau )\hfill & \text{ if }t>\tau .\hfill \end{array}$$
(1)
As a consequence, the resulting average latency $`t`$ and the risk $`\sigma `$ in downloading a page are given by
$$t=\underset{0}{\overset{\mathrm{}}{}}𝑑ttP(t)$$
(2)
and
$$\sigma ^2=(tt)^2=t^2t^2.$$
(3)
If we allow an infinite number of restarts, the recurrence relation (1) can be solved in terms of the partial moments $`M_n(\tau )=_0^\tau 𝑑tt^nP(t)`$:
$`t`$ $`=`$ $`{\displaystyle \frac{1}{M_0}}\left[M_1+\tau (1M_0)\right],`$
$`t^2`$ $`=`$ $`{\displaystyle \frac{1}{M_0}}\left\{M_2+\tau (1M_0)\left[2{\displaystyle \frac{M_1}{M_0}}+\tau \left({\displaystyle \frac{2}{M_0}}1\right)\right]\right\}.`$ (4)
In the case of a log-normal distribution
$$p(t)=\frac{1}{\sqrt{2\pi }x\sigma }\mathrm{exp}\left(\frac{(\mathrm{log}x\mu )^2}{2\sigma ^2}\right),$$
(5)
$`t`$ and $`t^2`$ can be expressed in terms of the error function:
$$M_n(\tau )=\frac{1}{2}\mathrm{exp}\left(\frac{\sigma ^2n^2}{2}+\mu n\right)\left[1+\mathrm{erf}\left(\frac{\mathrm{log}\tau \mu }{\sigma \sqrt{2}}\frac{\sigma n}{\sqrt{2}}\right)\right].$$
(6)
The resulting $`t`$-versus-$`\sigma `$ curve is shown in Fig. 2(a). As can be seen, the curve has a cusp point that represents the restart time $`\tau `$ that is preferable to all others. No strategy exists in this case with a lower expected waiting time or with a lower variance. The location of the cusp can be translated into the optimum value of the restart time to be used to reload the page.
There are many variations to the restart strategy described above. In particular, in Fig. 2(b), we show the family of curves obtained from the same distribution used in Fig. 2(a), but with a restriction on the maximum number of restarts allowed in each transaction. Even a few restarts yield an improvement.
Clearly, in a network without any kind of usage-based pricing, sending many identical data requests, to begin with, would be the best strategy as long as we do not overwhelm the target computer. On the other hand, everyone can reason in exactly the same way, resulting in congestion levels that would render the network useless. This paradoxical situation, sometimes called a social dilemma, arises often in the consideration of ”public goods” such as natural resources and the provision of services which require voluntary cooperation . This explains much of the current interest in determining the details of an appropriate pricing scheme for the Internet, since users do consume Internet bandwidth greedily when downloading large multimedia files for example, without consideration of the congestion caused by such activity.
Note that the histogram in Figure 1 represents the variance in the download time between different sites, whereas a successful restart strategy depends on a variance in the download times for the same document on the same site. For this reason, we cannot use the histogram in Figure 1 to predict the effectiveness of the restart strategy, but need to apply the similarly looking distribution of the respective Internet site. While a spread in the average download times of pages from different sites reduces the gains that can be made using a common restart strategy, it is possible to take advantage of geography and the time of day to fine tune and improve the strategy’s performance. As a last resort, it is possible to fine tune the restart strategy on a site per site basis.
As a final caution, we point out that with current client-server implementations, multiple restarts are detrimental and very inefficient since every duplicated request will enter the server’s queue and are processed separately until the server realizes that the client is not listening to the reply. This is an important issue for a practical implementation, and we neglect it here: our main assumption is that the restart strategy only affects the congestion by modifying the perceived latencies. This is only true if the restart strategy is implemented in an efficient and coordinated way on both the client and server side.
## Real Traffic on Freeways with Ramps
The recent study of the properties of “synchronized” congested highway traffic has generated a strong interest in the rich spectrum of phenomena occuring close to on-ramps . In this connection, a particularly relevant problem is that of choosing an optimal injection strategy of vehicles into the highway. While there exist a number of heuristic approaches to optimizing vehicle injection into freeways by on-ramp controls, the results are still not satisfactory. What is needed is a strategy that is flexible enough to adapt in real time to the transient flow characteristics of road traffic while leading to minimal travel times for all vehicles on the highway.
Our study presents a solution to this problem that explicitly exploits the naturally occuring fluctuations of traffic flow in order to enter the freeway at optimal times. This method leads to a more homogeneous traffic flow and a reduction of inefficient stop-and-go motions. In contrast to conventional methods, the basic performance criterion behind this technique is not the traffic volume, the optimization of which usually drives the system closer to the instability point of traffic flow and, hence, reduces the reliability of travel time predictions . Instead, we will focus on the optimization of the travel time distribution itself, which is a global measure of the overall dynamics on the whole freeway stretch. It allows the evaluation of both the expected (average) travel time of vehicles and their variance, where a high value of the variance indicates a small reliability of the expected travel time when it comes to the prediction of individual arrival times.
Both the average and the variance of travel times are influenced by the inflow of vehicles entering the freeway over an on-ramp. From these two quantities one can again construct a relation between the average payoff (the negative mean value of travel times) and the risk (their standard deviation). The optimal strategy will then correpond to the point in the curve that yields the lowest risk at a high average payoff. In the following, we will show that the variance of travel times has a minimum for on-ramp flows that are different from zero, but only in the congested traffic regime (which shows that the effect is not trivial at all). This finding implies that traffic flow can be optimized by choosing the appropriate vehicle injection rate into the freeway. Hence, in order to reach well predictable and small average travel times at high flows in the overall system, it makes sense to temporarily hold back vehicles by a suitable on-ramp control based on a traffic-dependent stop light . At intersections of freeways, this may require additional buffer lanes .
In order to obtain the travel time distribution of vehicles on a highway, we simulated traffic flow via a discretized and noisy version of the optimal velocity model by Bando et al. , which describes the empirical known features of traffic flows quite well . Moreover, we extended the simulation to several lanes with lane-changing maneuvers and different vehicle types (fast cars and slow trucks) . For lane changes, we assumed symmetrical (“American”) rules, i.e. vehicles could equally overtake on the left-hand or on the right-hand lane. Lane changing maneuvers were performed, when an incentive criterion and a safety criterion were satisfied . The incentive criterion was fulfilled, when a vehicle could go faster on the neighboring lane, while the safety criterion required that lane changing would not produce any accident (i.e., there had to be a sufficiently large gap on the neighboring lane) .
In addition to a two-lane stretch of length $`L=10`$ km, we simulated an on-ramp section of length 1 km with a third lane that could not be used by vehicles from the main road. However, vehicles entered the beginning of the on-ramp lane at a specified injection rate. Injected vehicles tried to change from the on-ramp to the main road as fast as possible, i.e. they cared only about the safety criterion, but not about the incentive criterion. The end of the on-ramp was treated like a resting vehicle, so that any vehicle that approached it had to stop, but it changed to the destination lane as soon as it found a sufficiently large gap. If the on-ramp was completely occupied by vehicles waiting to enter the main road, the injected vehicles formed a queue and entered the on-ramp as soon as possible. After injected vehicles had completed the 10 kilometer long two-lane measurement stretch, they were automatically removed from the freeway .
Our simulations were carried out for a circular road. After the overall density was selected, vehicles were homogeneously distributed over the road at the beginning, with the same densities on both lanes of the main road. The experiments started with uniform distances among the vehicles and their associated optimal velocities. The vehicle type was determined randomly after specifying the percentages $`r`$ of cars ($``$ 90%) and $`(1r)`$ of trucks ($``$ 10%). Notice that the effects discussed in the following should be more pronounced for increased $`r(1r)`$, since lane-changing rates seem to be larger and traffic flow more unstable, then.
We determined the travel times of all vehicles by calculating the difference in the times at which they passed the beginning and the end of the 10 kilometer long two-lane section. For mixed traffic composed of a high percentage of cars and a small percentage of slower trucks, we found narrow travel time distributions at small vehicle densities, where traffic flow was stable, while for unstable traffic flow at medium densities, the travel time distributions were broad (see Fig. 3).
If we plot the average of travel times as a function of their standard deviation (Fig. 4), we obtain curves parametrized by the injection rate of vehicles into the road and find the following: 1. With growing injection rate
$$Q_{\mathrm{rmp}}=\frac{1}{n\mathrm{s}},$$
(7)
the travel time increases monotonically. This is because of the increased density caused by injection of vehicles into the freeway. 2. The average travel time of injected vehicles is higher, but their standard deviation lower than for the vehicles circling on the main road. This is due to the fact that vehicle injection produces a higher density on the lane adjacent to the on-ramp, which leads to smaller velocities. The difference between the travel time distributions of injected vehicles and those on the main road decreases with the length $`L`$ of the simulated road, since lane-changes tend to equilibrate densities between lanes.
In addition, the standard deviation of the travel times has a minimum for finite injection rates, as entering vehicles tend to fill existing gaps and thus homogenize traffic flow. This minimum is optimal in the sense that there is no other value of the injection rate that can produce travel with smaller variance. In particular, gap-filling behavior mitigates inefficient stop-and-go traffic at medium densities. Above a density of 45 vehicles per kilometer and lane on the main road (measured without injection), the minimum of the travel times’ standard deviation occurs for $`n60`$. The reduction of the average travel time by smaller injection rates is less than the increase of their standard deviation. This result suggests that, in order to generate predictable and reliable arrival times, one should operate traffic at medium injection rates. For the case of 40 vehicles per kilometer and lane, the minimum of the standard deviation of travel times is located at $`n30`$, while for 35 vehicles per kilometer and lane, it is at $`n15`$. Below about 25 vehicles per kilometer and lane, vehicle injection does not reduce the standard deviation of travel times, since the travel time distribution is narrow anyway. At these densities, traffic flow is stable and homogeneous, so that no inefficient stop-and-go traffic exists and therefore no large gaps can be filled .
The curves displayed in Fig. 4 correspond to a given density $`\rho _{\mathrm{main}}`$ on the main road without injection of vehicles. The effective density $`\rho _{\mathrm{eff}}`$ on the freeway resulting from the injection of vehicles can be approximated by
$$\rho _{\mathrm{eff}}=\rho _{\mathrm{main}}+\frac{N_{\mathrm{inj}}}{IL},$$
(8)
where $`I=2`$ lanes, $`L=10`$ km. $`N_{\mathrm{inj}}`$ is the average number of injected vehicles present on the main road and can be written as
$$N_{\mathrm{inj}}=N_{\mathrm{tot}}\frac{𝒯_{\mathrm{inj}}}{𝒯_{\mathrm{tot}}𝒯_{\mathrm{inj}}},$$
(9)
where $`N_{\mathrm{tot}}=1000`$ is the total number of injected vehicles during the simulation runs, $`𝒯_{\mathrm{inj}}`$ is their average travel time, and $`𝒯_{\mathrm{tot}}`$ the time interval needed by all $`N_{\mathrm{tot}}=1000`$ vehicles to complete their trip. We point out that, in addition to these measurements, we used two other methods of density measurement which yielded similar results.
In contrast to Fig. 4, we also computed the dependence of the travel time characteristics on the resulting effective densities of vehicles. Fig. 5 shows the average of the travel times for vehicles in the main road as a function of their standard deviation. Once again, we observe a minimum of the standard deviation of travel times at high vehicle densities and medium injection rates. However, this time, an increase of the injection rate reduces the average travel times!
Figure 6 investigates the surprising reduction of the average travel times in more detail. While the injected cars experienced travel times that agreed with the case of no injection, the vehicles on the main road clearly profited from vehicle injection, if the effective density was the same. This means that, for given $`\rho _{\mathrm{eff}}`$, one can actually increase the average velocity $`V_{\mathrm{main}}=L/𝒯_{\mathrm{main}}`$ of vehicles by injecting vehicles at a considerable rate without affecting their travel times. This result is due to the increased degree of homogeneity caused by entering vehicles that fill gaps on the main road, which mitigates the less efficient stop-and-go traffic.
We point out that the injection-based reduction of travel times on the main road at a given effective density $`\rho _{\mathrm{eff}}`$ is related with a higher proportion
$$P=1\frac{\rho _{\mathrm{main}}}{\rho _{\mathrm{eff}}}$$
(10)
of injected vehicles, which implies a reduction in the number of vehicles circling on the main road. The relation between the injection rate and the percentage of injected vehicles is roughly linear (see Fig. 7).
The dependencies of the average travel times and their standard deviation on the proportion of injected vehicles are depicted in Figure 8. We find that the decrease in the average travel times is minor, while a significant reduction of the variance of travel times can be achieved by less than 5% of injected vehicles.
## Summary and Conclusions
Portfolio strategies can be successfully applied to systems in which the distribution of the quantity to be optimized is broad. This is the case for the download times in the World Wide Web as well as the travel time distributions in congested traffic flow. In this article, we showed how to reduce the average waiting times as well as their standard deviation (“risk”) by suitable injection strategies. In the case of the World Wide Web, it is possible to enforce smaller average waiting times by restarting a data request when the data were not received after a certain time interval $`\tau `$. This deforms the waiting time distribution towards smaller values, which automatically reduces the variance as well. Since the long tail of the waiting time distribution comes from the intermittent (“bursty”) behavior of Internet traffic, restart strategies partially manage to “calm down” these “Internet storms” by withdrawing requests in busy periods and restarting them later on. Similarly, the injection of vehicles to a freeway over an on-ramp can homogenize inefficient stop-and-go traffic by filling large gaps. In other words: The strategy exploits the naturally occuring fluctuations of traffic flow in order to allow the entry of new vehicles to the freeway at optimal times. In this way, the variation of travel times can be considerably reduced, which is favourable for more reliable travel time predictions. Moreover, at a given effective density, the average travel time decreases with increasing injection rate, i.e., with an increasing percentage of injected vehicles on the main road.
Acknowledgments: D.H. wants to thank the German Research Foundation (DFG) for financial support (Heisenberg scholarship He 2789/1-1). S.M. thanks for financial support by the Hertz Foundation. |
warning/0001/hep-th0001118.html | ar5iv | text | # A Note on Open Strings in the Presence of Constant 𝐵-Field
## Abstract
We consider the open string $`\sigma `$-model in the presence of a constant Neveu-Schwarz $`B`$-field on the world-sheet that is topologically equivalent to a disk with $`n`$ holes. First, we compute the $`\sigma `$-model partition function. Second, we make a consistency check of ideas about the appearance of noncommutative geometry within open strings.
PACS : 11.25.-w
Keywords: $`\sigma `$-models, strings
hep-th/0001118
HU Berlin-EP-00/08
The string $`\sigma `$-models provide a basic tool to derive stringy low energy effective actions that contain all orders of the fields. One example is the dilaton dependence of the effective action : the dilaton field is coupled to the Euler characteristic $`\chi `$ of the corresponding string world-sheet so that each world-sheet contributes a factor $`\text{e}^{\chi \phi }`$ into the effective action. The second one is the Born-Infeld (BI) action which was derived using the open string $`\sigma `$-model on the disk and adopting standard renormalization schemes.<sup>1</sup><sup>1</sup>1For a review see and references therein. Later it was realized that this action plays a crucial role in D-brane physics . It is also worth mentioning that there are known examples of string partition functions (one loop corrections to effective actions) for toroidal compactifications in the presence of constant metric and antisymmetric tensor which contain all orders of the fields .
Recently, it was pointed out that a special renormalization scheme within the open string $`\sigma `$-model, a point splitting regularization, results in a rather peculiar situation where the space-time (brane) coordinates do not commute (see and refs. therein). As a result, the low energy affective action becomes a noncommutative BI action. Since different renormalization schemes should be equivalent, there exists a change of variables ($`\sigma `$-model couplings) that relates the two BI actions. The aim of this note is to check that the Seiberg-Witten relations are consistent at higher genus topologies.
We study the open string $`\sigma `$-model with a constant background metric and a constant Neveu-Schwarz 2-form field on the world-sheet that is topologically equivalent to a disk with $`n`$ holes. Such topologies appear in the perturbation theory of open orientable strings. Thus the world-sheet action is given by
$$S=\frac{1}{4\pi \alpha ^{}}_{\mathrm{\Sigma }_n}d^2z\left(g_{ij}_aX^i^aX^j2i\pi \alpha ^{}B_{ij}\epsilon ^{ab}_aX^i_bX^j\right)+\chi \phi ,$$
(1)
where $`\mathrm{\Sigma }_n`$ means the disk with $`n`$ holes. $`g_{ij},B_{ij},\phi `$ are the constant metric, antisymmetric tensor and dilaton fields, respectively. $`X^i`$ map the world-sheet to the brane and $`i,j=1,\mathrm{},p+1`$. The world-sheet indices are denoted by $`a,b`$.
A natural object to compute is the $`\sigma `$-model partition function
$$Z_n[\phi ,g,B]=[dq]_n𝒟X\text{e}^S,$$
(2)
where $`[dq]_n`$ means the modular measure for the disk with $`n`$ holes extended to arbitrary dimension. It proves irrelevant for what follows, so we do not explicitly write it down <sup>2</sup><sup>2</sup>2See, e.g., where it is written down for the annulus topology..
To compute the partition function (2) one can follow the approach of namely, first integrate over the internal points of the world-sheet to reduce the integral to the boundaries and then split the integration variable $`X^i`$ on the constant $`x^i`$ and non-constant parts $`\xi ^i`$. Since the action is quadratic in $`\xi ^i`$ the problem is simply reduced to a computation of the corresponding functional determinant. The simplest case to consider is $`n=0`$ i.e., the path integral on the disk. In this case the problem is equivalent to the one solved in . This is clear by rewriting the term $`B_{ij}\epsilon ^{ab}_aX^i_bX^j`$ as a boundary interaction and replacing $`B`$ by $`F`$ <sup>3</sup><sup>3</sup>3For the sake of simplicity, we use the matrix notations here and below.. A subtle point we should mention here is due to a non-diagonal metric $`g`$. So, to get a GL(p+1) invariant answer, one must be careful with the measure of the integration (see ). Thus the partition function on the disk computed using the $`\zeta `$-function regularization is given by
$$Z_0[\phi ,g,B]=\text{e}^\phi [dx]\sqrt{detg}\left[det(1+2\pi \alpha ^{}g^1B)\right]^{\frac{1}{2}},$$
(3)
where $`[dx]=\frac{d^{p+1}x}{(2\pi \alpha ^{})^{p+1}}`$. The last factor is due to the integration over $`\xi ^i`$. This is clear within the perturbation theory where the $`B`$-term serves as an interaction.
Our aim now is to generalize the above result for arbitrary $`n`$. In fact, what we actually need is only a generalization for the last factor in the integrand. Let us give simple, but a little bit heuristic, arguments that lead to a desired answer. It turns out that the problem has a simple solution in the framework of the so-called sewing operation for the world-sheets. The latter is based on the idea of building surfaces by sewing together other ones. So let us begin with two disks and take a cylinder as a propagator between them. It is clear that the sewing operation produces a sphere. A crucial point here is that the partition function on the sphere does not depend on $`B`$. So, restricting ourselves to the $`B`$-dependence of the partition functions, namely $`Z_{\text{sphere}}[B]1,Z_0[B][det(1+2\pi \alpha ^{}g^1B)]^{\frac{1}{2}}`$, etc., we have
$$1Z_{\text{sphere}}[B]Z_0[B]\mathrm{\Pi }[B]Z_0[B].$$
(4)
As a result, we find the normalization of the propagator
$$\mathrm{\Pi }[B][det(1+2\pi \alpha ^{}g^1B)]^1.$$
(5)
Now let us make a consistency check and compute the partition function on the annulus. This can be done at least in two ways. The first one is to sew its boundaries to get the torus topology. The second way is to sew it with two disks. The both ways lead to the same result
$$Z_1[B]det(1+2\pi \alpha ^{}g^1B).$$
(6)
Note that such a result was also found by direct calculation in . This is clear by replacing $`BF`$ and rewriting the corresponding term in (1) as boundary interactions. To be more precise, what we found corresponds to orientable non-planar diagrams for vector fields (see also where this case corresponds to the two field strengths at the two boundaries set to be opposite, $`F_1=F_2`$).
It is now straightforward to get $`Z_n[B]`$. It is simply
$$Z_n[B][det(1+2\pi \alpha ^{}g^1B)]^{\frac{1+n}{2}}.$$
(7)
A crucial point here is that this factor does not depend on moduli. Hence the partition function is given by
$$Z_n[\phi ,g,B,\alpha ^{}]=\text{e}^{\chi \phi }[dq]_n[dx]\sqrt{detg}\left[det(1+2\pi \alpha ^{}g^1B)\right]^{1\frac{1}{2}\chi }.$$
(8)
In above we have used the fact that the Euler characteristic $`\chi `$ of a planar disc surface with $`n`$ holes is equal to $`1n`$.
Let us now give another way to derive the above result. The use of the point splitting regularization assumes that the metric $`g`$ becomes a new metric $`G`$ while all dependence on $`B`$ can be absorbed into the so-called star product that provides a multiplication law for other background fields. In fact, in this case the action for the kinetic term is given by
$$S=\frac{1}{4\pi \alpha ^{}}_{\mathrm{\Sigma }_n}d^2zG_{ij}_aX^i^aX^j+\chi \widehat{\phi },$$
(9)
while interaction terms include the build in star products. Since there are no interaction terms in the problem at hand, the partition function should have a simple structure due to standard dependence on the dilatonic field and our definition of the measure. Thus we have
$$\widehat{Z}_n[\widehat{\phi },G,\alpha ^{}]=\text{e}^{\chi \widehat{\phi }}[dq]_n[dx]\sqrt{detG}.$$
(10)
The new variables found by Seiberg and Witten in case of the disk topology are
$$\widehat{\phi }=\phi +\frac{1}{2}\mathrm{ln}det\left(G(g+2\pi \alpha ^{}B)^1\right),G=(g2\pi \alpha ^{}B)g^1(g+2\pi \alpha ^{}B).$$
(11)
A simple algebra shows that the partition functions (8) and (10) coincide. So, the Seiberg-Witten relations (11) hold on higher topologies too.
Finally, let us make some remarks.
(i) First, let us remark that what we found can be reinterpreted in terms of vector fields. Indeed, we can consider a set of Abelian vector fields with constant field strengths $`F^{(i)}`$ such that different $`A^{(i)}`$’s coupled to different boundaries (in other words, take $`n+1`$ Wilson factors as interactions). A configuration of the $`F`$’s that allows to rewrite the boundary interactions as the bulk term exactly corresponds what we considered. From the physical point of view, such a configuration represents $`n+1`$ free Wilson factors with each factor contributing the Born-Infeld determinant.
(ii) Second, it was conjectured in that there exists a suitable regularization that interpolates between the Pauli-Villars ($`\zeta `$-function) regularization and the point splitting one. In the framework of the open string $`\sigma `$-model such a regularization was further developed in where it was proposed to use the world sheet action
$$S=\frac{1}{4\pi \alpha ^{}}_{\mathrm{\Sigma }_n}d^2z\left(\stackrel{~}{G}_{ij}_aX^i^aX^j2i\pi \alpha ^{}\mathrm{\Phi }_{ij}\epsilon ^{ab}_aX^i_bX^j\right)+\chi \stackrel{~}{\phi },$$
(12)
where
$$\stackrel{~}{\phi }=\phi +\frac{1}{2}\mathrm{ln}det\left(\frac{\stackrel{~}{G}+2\pi \alpha ^{}\mathrm{\Phi }}{g+2\pi \alpha ^{}B}\right),\stackrel{~}{G}=(G^1\frac{1}{(2\pi \alpha ^{})^2}\theta _0G\theta _0)^1,\mathrm{\Phi }=\frac{1}{(2\pi \alpha ^{})^2}\stackrel{~}{G}\theta _0G.$$
(13)
Here $`\theta _0`$ is a free matrix parameter. $`G`$ is given by Eq. (11).
Repeating the analysis that led us to Eq. (8), we can easily write down the partition function in the case of interest
$$\stackrel{~}{Z}_n[\stackrel{~}{\phi },\stackrel{~}{G},\mathrm{\Phi },\alpha ^{}]=\text{e}^{\chi \stackrel{~}{\phi }}[dq]_n[dx]\sqrt{det\stackrel{~}{G}}\left[det(1+2\pi \alpha ^{}\stackrel{~}{G}^1\mathrm{\Phi })\right]^{1\frac{1}{2}\chi }.$$
(14)
A simple algebra shows that $`\stackrel{~}{Z}_n`$ coincides with $`Z_n`$, so it also passes a consistency check on higher genus topologies.
(iii) It is not difficult to formally repeat the previous analysis for superstring. To do so, it is more convenient to consider the NSR formalism within the point splitting regularization. In other words, we add a set of the fermionic fields $`\psi ^i`$ whose metric also is $`G_{ij}`$. It is simply to suggest what the superstring partition function should be.
$$\widehat{𝐙}_n[\phi ,G,\alpha ^{}]=\text{e}^{\chi \widehat{\phi }}[d𝐪]_n[dx]\sqrt{detG},$$
(15)
where $`[d𝐪]_n`$ means a proper modular measure for superstring. Clearly, there is no problem with rewriting this expression in terms of $`g,B`$ and $`\phi `$.
Acknowledgements. We would like to thank H. Dorn, R. Metsaev, and especially A.A. Tseytlin for useful discussions and critical comments. This work is supported in part by the Alexander von Humboldt Foundation and the European Community grant INTAS-OPEN-97-1312. |
warning/0001/astro-ph0001094.html | ar5iv | text | # Photoionization of Galactic Halo Gas by Old Supernova Remnants
## 1 Introduction
The nature of diffuse warm ionized gas (or warm ionized medium, WIM) in the interstellar medium has been puzzling since its discovery. In a series of articles, Reynolds (1984, 1985a, 1985b, 1988, 1989a, 1989b, and Reynolds & Tufte 1995) has pointed out several unusual properties of this gas including: high \[S II\] $`\lambda `$6716/H$`\alpha `$ ratios; low \[O I\] $`\lambda `$6300/H$`\alpha `$, \[O III\] $`\lambda `$5007/H$`\alpha `$ and He I $`\lambda `$5876/H$`\alpha `$ ratios; high scale height; and large power requirement for its ionization.
Among the characteristics of the WIM that are the most difficult to make sense of are its power requirement, its large scale height and its relative smoothness. The power required to maintain the ionization has been estimated by Reynolds (1984) to be $`9\times 10^5`$ ergs s<sup>-1</sup> cm<sup>-2</sup>, putting severe constraints on the source of ionizing photons. Only OB stars and supernovae appear to fulfill this power requirement. Citing the high efficiency required for conversion of explosion energy into photoionization, Reynolds (1984) favored early type stars as the primary source for maintaining the WIM. There are two principal difficulties with this source, however. First, several lines of sight in the plane apparently do not pass sufficiently close to any O or B stars (Reynolds, 1990). Equivalently, the relative smoothness of the H$`\alpha `$ background appears inconsistent with the very clumpy distribution of early type stars. Second, the large scale height observed for the emission, $`1`$kpc, is much higher than the stellar scale height. This problem has been addressed by Miller & Cox (1993) and by Dove & Shull (1994) by modeling the escape of Lyman continuum photons from the disk, particularly from OB associations. For some cases, however, even the existence of large “HII chimneys” cannot explain the ionization at high latitudes. A prime example of this is the line of sight towards HD 93521 that we discuss in detail below.
Cooling hot gas avoids some of these difficulties. First, although supernovae should have a scale height essentially identical to OB stars, hot gas is buoyant in the galactic gravitational field and rises to much larger scale heights. This is directly observed in the large scale heights of the highly ionized species C<sup>+3</sup>, N<sup>+4</sup> and Si<sup>+3</sup> (Savage, Sembach & Lu, 1997, and references therein). Second, hot gas by its very nature is a more smoothly distributed source of ionizing radiation than stars. While the effect of a single star or star cluster drops off with the square of the distance, the flux from hot gas, even distributed non-uniformly in regions of high emissivity, depends on the line of sight emission measure and thus tends to be much smoother. Nevertheless, the high efficiency required for conversion of SN energy into photoionizations remains. We do not argue in this paper that old supernova remnants are the sole or even primary source of the ionization responsible for the entire diffuse galactic H$`\alpha `$ background. Rather we show that cooling hot gas sets a floor on the ionization level for diffuse gas in the ISM in general and the galactic halo in particular. Moreover, for at least the line of sight towards HD 93521 and most probably for many high galactic latitude sight lines, cooling hot gas *is* the dominant source of the ionization. Thus we find that in many cases there is no need to devise a means of transporting photons from early type stars in the disk to the galactic halo.
The ionizing flux produced by cooling hot gas cannot be measured directly because the photons responsible for most of the ionization lie in the 13.6 – 50 eV range, a part of the EUV for which the diffuse background has not yet been definitively detected (see Vallerga & Slavin, 1998). We are therefore left to theoretical modeling to estimate the ionizing flux. Our procedure, detailed below, is to calculate a space and time averaged emissivity due to supernova remnants expanding in the ISM. The mean emissivity (along with an estimate of the fraction of ionizing photons that escape from the disk) gives us a prediction for the emission measure generated by photoionization due to hot gas emission. Using the emissivity and a simple model for the mean opacity in a cloudy medium interspersed with hot gas, we derive the mean intensity incident on a cloud face. \[It should be noted that we use the term “cloud” to refer to any cold or warm ($`T10^4`$ K) region of gas that is embedded in much lower density ($`n5\times 10^3`$ cm<sup>-3</sup>), hot ionized medium (HIM).\] We are then able to calculate such observables as the C<sup>+∗</sup> (C<sup>+</sup> in the excited fine structure level), S<sup>+</sup>, and S<sup>++</sup> column densities per cloud as well as a typical value for the galactic soft X-ray background and the X-ray surface brightness of the Galaxy as viewed from the outside. As we show below, our model is remarkably successful in explaining a variety of observations. Thus, while our model is clearly too simple to explain all the detailed observations of the soft X-ray background and ionization in the halo, its success provides evidence that cooling hot gas is an important source of ionization for the WIM.
In §2 below we describe current evidence on emission from old SNRs and give details about our model calculations of that emission. In §3 we use our model to predict the emission measure and H$`\alpha `$ intensity generated by ionization due to hot gas emission and compare with observations. In §4 we look at other ionization predictions from the model and compare with observations towards the line of sight towards HD 93521. We discuss how observations of S II, S III, and C II absorption lines constrain the morphology of the clouds and their thermal pressure, whereas \[N II\], \[S II\] and \[O I\] emission lines constrain the temperature. In §5 we describe how the ionization of helium in our model affects the X-ray opacity of the halo. In §6 we discuss the uncertainties in the model spectrum, the need for additional heating in the slow clouds, and the decaying neutrino model for the ionization proposed by Sciama (1990). Finally, we summarize our conclusions in §7.
## 2 The Soft X-ray Background from Old Supernova Remnants
### 2.1 Observations of the Local Background
The direct evidence for soft X-ray emission from hot gas comes primarily from proportional counter observations of the diffuse background by the Wisconsin group (e.g., McCammon et al., 1983) and by ROSAT (e.g., Snowden et al., 1997). The most important of these for the issue of ionization are the lower energy bands: the B and C bands ($`150`$ and 250 eV) of the Wisconsin instruments and the R1 and R2 bands ($`200`$ and 250 eV) defined for ROSAT. The all-sky maps produced by both sets of observations are essentially consistent with each other when allowances are made for the difference in filters, sensitivity and angular resolution (Snowden et al., 1995).
In the broadest terms, the observations show the whole sky to be glowing with soft X-ray emission with a general trend of greater brightness towards the poles. While there has been much discussion of the proper interpretation of the results (e.g., Cox, 1998; McKee, 1998), there is now a general consensus that most of the observed emission comes from a cavity in the neutral hydrogen that surrounds the Sun and is filled with hot ($`T10^6`$K), very low density gas. At high latitude the 1/4 keV (C band) background also has substantial contributions from more distant emission, in some directions dominating the local component. In the galactic plane this seems not to be the case, though the presence of low column density clouds, seen by optical and EUV absorption, has yet to be fully taken into account in the modeling of the observed flux (see Slavin, 1998).
### 2.2 Models for the Soft X-ray Background
Several studies have been carried out to explain the origins of the soft X-ray emission. McKee & Ostriker (1977) put forward a comprehensive model for the diffuse ISM in which supernova remnants sweep the warm and cold gas into shells and isolated clouds, leaving most of the medium filled with hot gas. The observed local emission is seen as being somewhat greater than in the typical ISM due to our being in an old supernova remnant (their estimate corresponds to an age of about $`3\times 10^5`$ yr). Cox & collaborators (Cox & Anderson, 1982; Edgar, 1986; Edgar & Cox, 1993; Smith & Cox, 1998) have explored a number of models aimed specifically at explaining the SXRB as the result of the Sun’s position inside a supernova remnant (possibly due to several explosions). Our model that we present below is not aimed specifically at explaining the observed SXRB and is, in this sense, more like the McKee & Ostriker model. We do not attempt, however, to model the medium as a whole and our model of the emissivity does not rest on any particular model for the morphology of the ISM. On the other hand, our model for the flux incident on clouds does require assumptions about the distribution of opacity in the medium and we use this fact below to draw conclusions on typical cloud (i.e., WNM and WIM) sizes. In addition, the data for the line of sight towards HD 93521 gives us information on the typical pressures in clouds that we can use to set limits to the filling factor of warm clouds in this direction. We find, in accordance with Reynolds (1991) and Spitzer & Fitzpatrick (1993), that the warm clouds have a fairly small filling factor, $`10`$% (see §4.3 below).
A model for the mean emission from a population of evolving SNRs has previously been calculated by Cox & McCammon (1986) for the purpose of constraining the properties of SNRs in the galaxy M101. They carried out analytical calculations of the distribution over temperature of the surface averaged emission measure. A set of numerical calculations similar to the ones we present here was performed by Miller (1994) again with observations of external galaxies in mind. The aim of our calculations is somewhat different, being focussed on the ionizing properties of the x-ray/EUV emission. At the time of publication of Cox & McCammon (1986), only upper limits to the soft X-ray emission were available. With the advent of *ROSAT*, the X-ray emission has been detected for M101 and a number of other galaxies. As we discuss below, the fraction of soft X-ray emission that escapes the galaxy provides one of the tests of our model. In addition, we examine in more detail the spectrum produced by the cooling, supernova-shocked gas and its effect on the ionization structure of the WIM and the WNM.
### 2.3 The Hot Gas Emission Model
The primary purpose of this paper is to explore the effects of photoionization by cooling hot gas in the galaxy. As a consequence, we have focussed on the calculation of the emission and have made very simple assumptions regarding the ISM. To calculate the mean emissivity, we assume independent evolution of SNRs expanding in a uniform density medium. The mean emissivity of the hot gas is then
$$ϵ_\nu =S𝑑tϵ_\nu 𝑑V$$
(1)
where $`S`$ is the supernova rate per unit volume, $`ϵ_\nu `$ is the emissivity per unit frequency in the remnant as a function of position and time within the remnant. The integration is carried out over the volume of the remnant and time evolution of the remnant. We shall generally work with quantities measured per unit area of the disk, so we define $`S_AS𝑑z`$ as the supernova rate per unit area and $`ϵ_{\nu A}ϵ_\nu 𝑑z`$ as the mean emissivity per unit area. Normalizing to the explosion energy, $`E_{\mathrm{SN}}`$ we define:
$$\varphi _\nu \frac{ϵ_\nu }{SE_{\mathrm{SN}}}=\frac{ϵ_{\nu A}}{S_AE_{\mathrm{SN}}}.$$
(2)
Thus $`\varphi _\nu `$ is the fraction of the supernova power radiated per unit frequency interval. Equation (2) shows that the total supernova power, or equivalently, the supernova power per unit disk area, is separable from the calculation of the mean spectrum that determines the emission characteristics of the SNRs. This formulation brings out the fact that it is this distribution in frequency that is critical in determining the photoionizing properties of the cooling hot gas in SNRs.
To carry out the calculations of $`\varphi _\nu `$ we perform high resolution numerical calculations using a 1-D (spherically symmetric) hydrodynamics code. This code, written by us, borrows from the VH-1 code (see http://wonka.physics.ncsu.edu/pub/VH-1/index.html) and uses the same piecewise parabolic method (PPM) with a Lagrangian step followed by a remap onto the fixed grid. We include the non-equilibrium ionization of the gas and radiative cooling appropriate to the ionization. For these calculations the ionization, recombination and radiative cooling rates have been generated using the Raymond & Smith (1977, plus updates) codes. Care has been taken to assure that the remapping step conserves the mass in each ion stage. We have carried out resolution studies and have found that the results converged at high resolution. The spectra used in the results that follow are for our highest resolution (0.07 pc for the $`n_a=0.1`$cm<sup>-3</sup> case, where $`n_a`$ is the average ambient preshock hydrogen nucleus density of the medium into which supernovae expand). The evolution of a remnant has been followed until nearly all the energy has been radiated and a small fraction of the explosion energy remains in thermal energy. The spectrum is then generated using the temperature, density and ionization profiles as a function of remnant age. We have done some runs in which thermal conduction was included (using an operator splitting technique) with no inhibition of conduction other than the limitation imposed by saturation (Cowie & McKee, 1977). We have found the differences in the resulting spectra to be relatively small and therefore do not present results for these cases in this paper. We have done runs for a variety of ambient density cases, $`n_a=0.04`$, 0.1, 0.3 and 1.0 cm<sup>-3</sup>. We assume an explosion energy of 10$`^{51}`$ergs in all cases. We do not include any magnetic field effects; this assumption is discussed below.
Having generated $`\varphi _\nu `$ (and thus the mean emissivity, $`ϵ_\nu `$, assuming $`SE_{\mathrm{SN}}`$ is known) we then need a model for the opacity of the medium in order to calculate the mean intensity, $`J_\nu `$. The true opacity of the ISM to ionizing photons is clearly extremely complicated. It depends on the distribution of neutral hydrogen column density on scales ranging from $`10^{18}`$cm<sup>-2</sup> to $`10^{20}`$cm<sup>-2</sup>, including the correlations among clouds and local large scale structures. To create a realistic model for the opacity of the ISM is thus well beyond the scope of this paper. Nevertheless, any given cloud that is to be subject to the ionizing field generated by hot gas will receive radiation from sight lines that will pass through a range of cloud columns present in the medium. For this reason we adopt a model for the opacity that is as simple as possible while still allowing for an interspersion of emitting and absorbing material, as must be the case in the ISM.
Our approach is to assume a uniform medium in the sense that the clouds, assumed to be of some typical optical depth, $`\tau _{c\nu }`$, are distributed randomly in the hot gas. The optical depth due to these clouds at frequency $`\nu `$ is
$$\tau _\nu =𝒩_{\mathrm{los}}\left(1e^{\tau _{c\nu }}\right)$$
(3)
(Bowyer & Field, 1969, see the Appendix), where $`𝒩_{\mathrm{los}}`$ is the expected number of clouds along the line of sight. The average value of the mean intensity in the disk and halo is
$`\overline{J}_\nu `$ $`=`$ $`{\displaystyle \frac{1}{\tau _{0\nu }}}{\displaystyle _0^{\tau _{0\nu }}}𝑑\tau _\nu {\displaystyle \frac{1}{4\pi }}{\displaystyle I_\nu 𝑑\mathrm{\Omega }},`$ (4)
$`=`$ $`{\displaystyle \frac{ϵ_{\nu A}}{4\pi \tau _{0\nu }}}(1\eta _\nu )`$ (5)
where $`\tau _{0\nu }`$ is the optical depth through the full disk and halo,
$$\eta _\nu =\frac{1}{\tau _{0\nu }}\left[\frac{1}{2}+E_3(\tau _{0\nu })\right],$$
(6)
is the mean escape probability, and $`E_3(\tau _{0\nu })`$ is an exponential integral. As shown in the Appendix, this value of the mean intensity is approximately equal to the expected value of the mean intensity in the intercloud medium, provided it has a substantial filling factor.
Thus our model for the mean intensity incident on a cloud depends on five parameters: the cloud optical depth $`\tau _{c\nu }`$, the disk optical depth, $`\tau _{0\nu }`$, the SN explosion energy, $`E_{\mathrm{SN}}`$, the supernova rate per unit area $`S_A`$, and $`\varphi _\nu `$. The optical depths are related to the column densities by $`\tau _{c\nu }=N_{\mathrm{H}^0c}\sigma _\nu `$ and $`\tau _{0\nu }=N_{\mathrm{H}^0}/N_{\mathrm{H}^0c}[1\mathrm{exp}(N_{\mathrm{H}^0c}\sigma _\nu )]`$ (eq A22), where $`N_{\mathrm{H}^0c}`$ is the HI column through the standard cloud and $`N_{\mathrm{H}^0}`$ is the total HI column of the disk and the halo perpendicular to the disk plane. The values for the column densities are discussed in §3 below. For $`S_A`$ we use the results from McKee & Williams (1997), $`S_A=3.8\times 10^5`$kpc<sup>-2</sup> yr<sup>-1</sup> at the Solar Circle. However, for the line of sight toward HD 93521 we find it necessary to effectively reduce this number slightly in order to match the observed H$`\alpha `$ intensity. In general, $`E_{\mathrm{SN}}`$ and $`\varphi _\nu `$ are not independent, since the explosion energy will affect the evolution (temperature, ionization, etc.) of the remnant which in turn determines $`\varphi _\nu `$. Nevertheless, the separation of the two quantities emphasizes the fact that $`\varphi _\nu `$ is the efficiency with which the available energy is radiated over frequency. Thus our formulation could be extended to other potential sources of ionizing radiation such as the conduction fronts surrounding evaporating clouds (McKee & Cowie, 1977) or the interfaces in turbulent mixing layers (Slavin, Shull & Begelman, 1993).
### 2.4 The Model X-ray/EUV Spectrum
Figure 1 shows our model X-ray/EUV spectrum compared with a collisional ionization equilibrium (CIE), $`T=10^6`$K, unabsorbed spectrum. The latter is scaled so as to match the all-sky average B band count rate observed by the Wisconsin group (McCammon et al., 1983). Of particular note are the greater fluxes at low energies ($`h\nu 13.630`$ eV) in the cooling model as opposed to the CIE spectrum. As can be seen from the spectrum, these photons dominate the higher energy photons in producing ionization in interstellar clouds.
A great deal of effort has gone into modeling the observed diffuse soft X-ray background (SXRB) (e.g. Cox & Anderson, 1982; Jakobsen & Kahn, 1986; Snowden et al., 1998; Smith & Cox, 1998). Our model has not been created with the aim of matching the SXRB observations. In particular, we do not attempt to match the plane to pole variation in intensity or the variation of the observed X-ray band count rate ratios with $`N`$(H I). These characteristics of the data apparently demand the existence of an irregularly shaped Local Bubble and as such are inconsistent with any simple global model. On the contrary, we have calculated a time and space averaged emission spectrum that should approximate the observed spectrum at high galactic latitudes, where the lower H I density allows us to see farther along a line of sight and where there is clear evidence that the observed flux has substantial contributions from distant emission (Garmire et al., 1992).
For comparison of our model results to the data, we need to specify all the model parameters as described above. Of particular importance is $`N_{\mathrm{H}^0}`$, as it effectively determines the mean free path for the soft X-ray photons. The average disk column density at the solar circle is $`6.2\times 10^{20}`$ cm<sup>-2</sup> (Dickey & Lockman, 1990). However, the Local Bubble apparently has carved out a hole in the H I distribution, resulting in a substantially lower value of $`N_{\mathrm{H}^0}`$ locally. We use a value of $`3.2\times 10^{20}`$ cm<sup>-2</sup> (Kulkarni & Fich, 1985) for our comparisons at high latitude. For the typical cloud column density, we adopt $`N_{\mathrm{H}^0c}=1.46\times 10^{19}`$ cm<sup>-2</sup>, the mean value observed for clouds along the line of sight towards HD 93521. We use our standard spectrum ($`n_a=0.1`$cm<sup>-3</sup>, no conduction, and $`E_{\mathrm{SN}}=10^{51}`$ergs) to determine $`\varphi _\nu `$. In order to match the observed H$`\alpha `$ intensity toward HD 93521, we find it necessary to reduce $`S_A`$ by a factor 0.68 (§3). In other words, the average supernova rate $`S_A`$ at the Solar Circle creates enough cooling hot gas to produce somewhat more hydrogen ionization than is observed toward the particular line of sight to HD 93521; this clearly demonstrates that the ionizing power of supernova remnants is significant at high latitude. A slight reduction in $`S_A`$ along this line of sight, or of the radiation field along this line of sight, is quite reasonable in view of the observed variations in the soft X-ray background (e.g. Snowden et al., 1997).
The most important comparison we can make with the data is a spatially averaged measurement at low energy. For this purpose, the Wisconsin B and C band data are the best choice. From the publicly available maps we find that for $`|b|>45`$° the average count rates for the B and C bands are 63 and 184 counts s<sup>-1</sup>, respectively. For the line of sight towards HD 93521, the count rates are 84 and 210 counts s<sup>-1</sup>. There is a fair amount of scatter, with the B band rate ranging from 18 to 134 and the C band rate from 62 to 302. The C/B ratio shows somewhat less scatter, ranging from 1.7 to 8.2. For the parameter choices detailed above we find a B band rate of 70 and a C band rate of 202 counts s<sup>-1</sup>. Given that we have made no attempt to vary the parameters of our model to fit the SXRB data, we find this agreement remarkable.
### 2.5 Soft X-ray Emission from External Galaxies
Another test of our model is to compare it with the soft X-ray emission from external galaxies. This comparison has the advantage that local variations are averaged out, but the disadvantage that various conditions affecting the strength and spectrum of the emission (e.g. $`n_a`$, $`S_A`$) could be significantly different in other galaxies. We estimate the X-ray luminosity of the Galaxy using the average conditions at the solar circle; thus, we use $`N_{\mathrm{H}^0}=6.2\times 10^{20}`$cm<sup>-2</sup> for the disk thickness, $`S_A=3.8\times 10^5`$ kpc<sup>-2</sup> yr<sup>-1</sup> for the SN rate per unit area, and 530 kpc<sup>2</sup> for the effective area of the disk (McKee & Williams, 1997). Using equations (5) and (6), we find that the total luminosity in ionizing photons that escape from the disk in this case is $`1.1\times 10^{40}`$ ergs s<sup>-1</sup>, while the luminosity in X-ray photons ($`E>100`$ eV) is $`2.2\times 10^{39}`$ ergs s<sup>-1</sup>. Note that our estimate for the ionizing luminosity is very sensitive to our assumption that the emission and absorption are uniformly mixed; our estimate for the X-ray luminosity also depends on this assumption, although to a lesser extent (see the Appendix).
Observations of diffuse soft X-ray emission in external galaxies are quite difficult. Proper accounting for point sources, backgrounds and galactic absorption are required (Cui, Sanders & McCammon, 1996, hereafter CSM). The intensity or emission measure estimates that result are not tightly constrained. CSM observed five galaxies (NGC 3184, NGC 4736, M101, NGC 4395 and NGC 5055) in the R12 or R12L bands of ROSAT ($`E100`$-284 eV). In all cases the emission is sharply peaked near the center of the galaxy and drops off quickly with galactocentric radius. For our standard case ($`n_a=0.1`$cm<sup>-3</sup>, no conduction) we find in the R12L band a rate of $`550\times 10^6`$ counts s<sup>-1</sup> arcmin<sup>-2</sup> (this is the standard unit for such ROSAT observations). Comparing this result to the surface brightness measured by CSM at the boundary between rings 2 and 3 (which corresponds approximately to the Solar Circle), we find that our result lies below the 95% upper limits for all the galaxies observed, and above the 95% confidence lower limits for all but two of the galaxies listed, namely NGC 4736 and NGC 5055. The fact that we get agreement with the observations, despite the large luminosity is a result of the softness of our spectrum. As has been pointed out previously (Cox & McCammon, 1986), the relative faintness of the observed soft X-rays from external galaxies requires that much of the supernova energy is absorbed in the disk. Our model provides a natural explanation of why that occurs. Most of the radiated supernova energy is generated at low energies (see Figure 1) and few of the photons escape.
## 3 Photoionization Due to the Hot Gas Emission
Soft X-ray/EUV photons that are radiated by hot gas from SNRs will either be absorbed or escape the galaxy. Since the mean half-thickness of the disk is $`N_{\mathrm{H}^0}3\times 10^{20}`$cm<sup>-2</sup>, nearly all of the ionizing photons less energetic than 0.5 keV will be absorbed. The surprisingly low soft X-ray brightness observed for external galaxies (CSM) is an indication that a large fraction of the X-rays generated by SNRs are being retained by the galaxies and therefore should contribute to their ionization. The photoionization rate of species $`i`$ per unit area of disk is then
$`\zeta _{iA}`$ $`=`$ $`{\displaystyle _{\nu _i}^{\mathrm{}}}𝑑\nu (1\eta _\nu )\left({\displaystyle \frac{\sigma _{i\nu }}{\sigma _\nu }}\right){\displaystyle \frac{ϵ_{\nu A}}{h\nu }},`$ (7)
$``$ $`y_i{\displaystyle _{\nu _i}^{\mathrm{}}}𝑑\nu {\displaystyle \frac{ϵ_{\nu A}}{h\nu }},`$ (8)
since $`\eta _\nu 0`$; here $`\nu _i`$ is the frequency threshold for ionization of species $`i`$ and $`y_i`$ is the fraction of the photons absorbed by species $`i`$. For hydrogen, $`y_\mathrm{H}`$ is very close to unity since almost every helium ionization results in the emission of a hydrogen ionizing photon (Osterbrock, 1989); for helium, $`y_{\mathrm{He}^0}0.6`$ in weakly ionized regions since hydrogen competes for the ionizing photons. Defining $`ϵ_{iA}^{}_{\nu _i}𝑑\nu ϵ_{\nu A}/h\nu `$, which is the rate of generation of photons with $`\nu \nu _i`$ per unit area, we can express the ionization equation as simply
$$\zeta _{iA}=y_iϵ_{iA}^{}.$$
(9)
Note that here and below we use the superscript (\*) to denote units of photons instead of units of energy; this should not be confused with the conventional notation for absorption lines from excited fine structure states, such as $`\mathrm{C}^+`$. The ionization rate is thus almost independent of the morphology of the $`N`$(H I) distribution—it does not matter whether the H<sup>0</sup> is in small, optically thin clouds or large optically thick slabs (for helium, there can be a weak dependence of $`y_{\mathrm{He}^0}`$ on the cloud morphology). The ionization of hydrogen determines the emission measure
$$EM_{\mathrm{pc}}=\frac{ϵ_{\mathrm{H}A}^{}}{\alpha _\mathrm{H}^{(2)}}=1.25\times 10^6T_4^{0.8}ϵ_{\mathrm{H}A}^{}\mathrm{cm}^6\mathrm{pc}$$
(10)
from equation (A12). The subscript “pc” on $`EM_{\mathrm{pc}}`$ indicates that it is measured in units of cm<sup>-6</sup> pc; $`ϵ_{\mathrm{H}A}^{}`$ has units photons cm<sup>-2</sup> s<sup>-1</sup>; and $`EM_{\mathrm{pc}}`$ is the emission measure normal to the plane of the galaxy and through the entire disk. Photoionization by emission from hot gas thus sets a floor on the emission measure that should be observable on nearly all lines of sight.
To relate the ionization rate to the properties of individual SNRs, we define
$$\varphi (>\nu _i)_{\nu _i}^{\mathrm{}}𝑑\nu \varphi _\nu $$
(11)
as the fraction of the SN energy emitted above $`\nu _i`$ and
$$\frac{1}{h\overline{\nu }_i}\frac{1}{\varphi (>\nu _i)}_{\nu _i}^{\mathrm{}}𝑑\nu \frac{\varphi _\nu }{h\nu }$$
(12)
is the inverse of the mean energy of ionizing photons for species $`i`$. In terms of these quantities, the number of photons emitted above $`\nu _i`$ by an SNR is
$$𝒩_\gamma (>\nu _i)=\frac{\varphi (>\nu _i)E_{\mathrm{SN}}}{h\overline{\nu }_i}=4.6\times 10^{61}\varphi (>\nu _i)\left[\frac{E_{\mathrm{SN}}}{10^{51}\mathrm{erg}}\right]\left[\frac{13.6\mathrm{eV}}{h\overline{\nu }_i}\right],$$
(13)
and the rate of generation of ionizing photons per unit area is
$`ϵ_{iA}^{}`$ $`=`$ $`S_A𝒩_\gamma (>\nu _i)`$ (14)
$`=`$ $`1.26\times 10^6\left({\displaystyle \frac{S_A}{3.8\times 10^5\mathrm{kpc}^2\mathrm{yr}^1}}\right)\left[{\displaystyle \frac{𝒩_\gamma (>\nu _i)}{10^{61}\mathrm{photons}}}\right]\mathrm{photons}\mathrm{cm}^2\mathrm{s}^1.`$ (15)
In Table 1 we illustrate the nature of the spectrum emitted by SNR for the different model assumptions by listing the fraction of the SN energy emitted in several different energy bands and the commonly used ionizing photon ratio, $`Q`$(He<sup>0</sup>)/$`Q`$(H<sup>0</sup>). This latter is the ratio of the total number of He<sup>0</sup> ionizing photons to H<sup>0</sup> ionizing photons emitted. The energy bands are defined as follow: $`\varphi _\mathrm{H}`$: 13.6–24.6 eV, $`\varphi _{\mathrm{He}^0}`$: 24.6–54.4 eV, $`\varphi _{\mathrm{He}^+}`$: 54.4–284 eV, $`\varphi _X`$: $`>284`$ eV. The first two are defined by the hydrogen and helium ionization edges. The edge at 284 eV is the carbon edge above which C and other elements are responsible for a significant amount of the opacity. The C edge also corresponds (not coincidentally) to the cutoff in response in the Wisconsin C band and the *ROSAT* R1 and R2 bands. It is clear that while there is a general trend towards harder spectra for SNR expanding in higher ambient density, $`n_a`$, the trend is not monotonic in all the energy bands. This is due to fact that the emission spectrum is line-dominated so that a few strong lines from ions of an abundant element can substantially affect the total emission within an energy band.
Table 2 lists the emission measures generated by our modeled SNRs, together with characteristics of the ionizing spectrum, for several values of the ambient density and thermal conductivity. We have assumed $`T=8000`$ K in evaluating the emission measure, and we have divided the value given in equation (10) by 2 in order to give the value that would be measured from the midplane of the disk. Typically somewhat more that half the energy of a SNR is emitted in ionizing photons, and the mean energy of these photons is only about 20 eV.
The effects on the spectrum of varying conductivity are quite small, despite the fact that the temperature at the center of the remnant is radically different for the two cases. When thermal conduction is included, the central temperature flattens quickly and the thermal energy is shared with matter further out in the remnant. The effects on the emission are small, however, because the emissivity is sharply peaked towards the edge of the remnant where the temperature and density profiles are little effected by conduction.
The effects of varying the ambient density are similarly small. To first order, the luminosity, spectrum, and $`EM`$ are independent of $`n_a`$. However, as can be seen from Tables 1 and 2, there is a weak trend in that higher ambient density leads to slightly harder emission spectrum. This is due to the fact that a supernova remnant evolving in a higher density medium becomes radiative earlier, when the gas inside is at a higher temperature. Since the radiative phase is the period of greatest luminosity for the remnant, this has an effect on the mean emissivity even though the radiative phase lasts for only a small fraction of the remnant lifetime. The differences in the hardness of the spectrum for the different cases can be seen both from the mean energy of ionizing photons and the fraction of the emitted energy in ionizing photons.
It is noteworthy that while the ambient density in our models varies by a factor of 25, the emission measure generated for the extreme cases differ by only 34%. The intensity of an H$`\alpha `$ line in Rayleighs is related to the emission measure by $`I^{}(\mathrm{H}\alpha )=0.445EM_{\mathrm{pc}}`$ R for $`T=8000`$ K from equation (A14). (Note that although the relation between $`I^{}`$ and $`EM`$ depends on temperature, the relation of $`I^{}`$ to the underlying emissivity $`ϵ_{\mathrm{H}A}^{}`$ is almost independent of temperature—see eq. A15.) Our standard model gives $`(EM_{\mathrm{pc}}/2)=1.2`$ cm<sup>-6</sup> pc for the half disk, which corresponds to a predicted intensity at $`b=90\mathrm{°}`$ of $`I^{}(\mathrm{H}\alpha )=0.53`$ R. By comparison, Reynolds (1991) has suggested 1 R as a “typical” minimum value for the diffuse H$`\alpha `$ background at $`b=90\mathrm{°}`$, with the emission roughly following a $`\mathrm{csc}b`$ law. Those conclusions were based on a sampling of regions observed at moderately high latitudes using a less sensitive instrument than the current instrument, WHAM (Wisconsin H$`\alpha `$ Mapper). Even those data showed substantial variation in the value of $`I\mathrm{sin}b`$, and more recent data reveal the existence of many regions of low emission, well below the 1 R zenith value of the earlier estimate. One of these is the direction towards HD 93521 discussed in detail below. Thus, SNRs are capable of accounting for the observed H$`\alpha `$ intensity in at least some directions in the halo. It is clear that the cooling hot gas in our model makes a substantial contribution to the ionization of hydrogen in the warm ionized medium of our galaxy and sets a lower limit on the ionization.
Another way of testing our model is to compare it to the ionization of other elements. Absorption line data provide particularly good tests since high spectral resolution observations allow us to separate different velocity components and thus model individual clouds. A particularly useful example is the line of sight towards HD 93521.
## 4 The HD 93521 Line of Sight
### 4.1 The Data
The line of sight towards the halo star HD 93521 ($`\mathrm{}=183\mathrm{°},b=62\mathrm{°}`$), while in some ways quite unusual, reveals important information about the background ionizing radiation field in the galaxy. The line of sight is unusual in two ways. First, there appear to be no stars close enough to it to be significant ionizing sources (with the exception, of course, of HD 93521 itself, which could ionize no more than one of the clouds, see below). Note that it is not unusual for any given point in the galaxy to be far from O stars, rather it is the fact that the entire 1.7 kpc line of sight is far from any O stars that makes it somewhat exceptional. This can be seen by examining figure 16 (particularly 16a and 16b) of Miller & Cox (1993). In that figure, they present the outlines of the extended H II regions resulting from the observed distribution of nearby O stars. The line of sight does not pass through any of the H II regions, reflecting the fact that no O stars are close to the line of sight. \[Note: in figure 16c the Sun and the line of sight towards HD 93521 appear to be on the edge of an H II region. It is unclear which star might be responsible for this ionization (perhaps $`\zeta `$ Oph) and it is unlikely to be capable of ionizing clouds towards HD 93521.\] Secondly, observations by Spitzer & Fitzpatrick (1993, hereafter SF) show that there is extremely little cold neutral gas along the line of sight. Thus we have in this line of sight an excellent opportunity to observe ionization of warm diffuse gas by non-stellar sources without confusion by interspersed cold gas.
The absorption line data for HD 93521, including lines of S II, S III, C II, and several other lines, reveal a number of velocity resolved features, or “clouds”, all indicating partial ionization (SF). Complementing the absorption line data are new observations using the Wisconsin H$`\alpha `$ Mapper (WHAM; Reynolds et al., 1998) of diffuse emission in H$`\alpha `$, \[N II\]$`\lambda `$6584, and \[S II\]$`\lambda `$6716 (Pifer et al., 1999). The velocity components in the absorption lines split roughly into two groups of 4 clouds each, the “fast” clouds with $`70\mathrm{km}\mathrm{s}^1<v_{\mathrm{LSR}}<30\mathrm{km}\mathrm{s}^1`$ and the “slow” clouds with $`20\mathrm{km}\mathrm{s}^1<v_{\mathrm{LSR}}<10\mathrm{km}\mathrm{s}^1`$. The emission lines cannot be resolved into individual cloud components at this time, but they do split up into broad “fast” and “slow” components that correspond to the absorption line groups. Hausen et al. (1999) find that for the slow clouds, $`I^{}(\mathrm{H}\alpha )=0.22\pm 0.06`$R, and for fast clouds, $`I^{}(\mathrm{H}\alpha )=0.16\pm 0.03`$R. Pifer et al. (1999) find $`I^{}`$(\[N II\]$`\lambda 6584)=0.22`$ R for the slow clouds and 0.29 R for the fast ones; they also find $`I^{}`$(\[S II\]$`\lambda 6716)=0.24`$ R for the slow clouds and 0.19 R for the fast ones (see Table 3).
### 4.2 Modeling the Ionization
In order to calculate the ionization of clouds subjected to the SNR-generated flux of our model (Figure 1) we make use of the photoionization/thermal equilibrium code CLOUDY (Ferland, 1996). The flux has been scaled by reducing $`S_A`$ by 0.68 in order to give the observed value for the H$`\alpha `$ intensity. In order to approximate the effect of isotropic, diffuse radiation using CLOUDY, which is designed for point source calculations, we use the “illuminate” command with the parameter value 60, which simulates flux coming into the cloud face at an angle of 60 degrees. We also use the “case b” command to avoid the artificial excitation of H$`\alpha `$ emission due to absorption of FUV background photons at Lyman $`\beta `$. We use a fixed temperature and have explored a range from 6000 K to 9000 K. Fixed $`T`$ avoids the complications of heating/cooling balance that involve the amount of dust heating (particularly due to PAHs) (Bakes & Tielens, 1998), which is only crudely included in the code, and the uncertain details of the elemental abundances. We discuss heating further below.
Although the usual approach is to assume that individual velocity features are physical “clouds”, we examined the possibility that the features were actually formed from collections of smaller clouds or “cloudlets”. This affects the opacity of the medium and thus the hardness of the radiation field incident on a typical cloud, but it does not affect the emission measure along the line of sight since nearly all ionizing photons are absorbed somewhere in the disk. The difference between many small cloudlets and a few big clouds is greatest near the Lyman limit, where for small cloudlets the incident flux is much less. We have constructed two models to try to differentiate between these two cases. The “cloudlet” model has cloudlet column densities of $`N_{\mathrm{H}^0}=7\times 10^{17}`$cm<sup>-2</sup>. The simple cloud model has $`N_{\mathrm{H}^0}=1.46\times 10^{19}`$cm<sup>-2</sup>, the mean for the observed clouds towards HD 93521. The only data we have found that differentiates these two fairly extreme cases from each other is the ratio of $`N_{\mathrm{S}^{++}}/N_{\mathrm{S}^+}`$. This ratio has unfortunately only been observed for the fast clouds, but if we assume that the fast and slow clouds are at least geometrically similar (which might be doubtful, see below), then our results favor larger clouds over cloudlets. This is best seen by forming a ratio that is insensitive to the pressure of the clouds but still sensitive to the cloud size, $`(N_{\mathrm{S}^{++}}/N_{\mathrm{S}^+})\times (N_{\mathrm{C}^+}/N_\mathrm{S})`$, where $`N_\mathrm{S}N_{\mathrm{S}^+}+N_{\mathrm{S}^{++}}`$ to good accuracy. The lack of sensitivity to pressure derives from the fact that $`N_{\mathrm{S}^{++}}/N_{\mathrm{S}^+}`$ goes as 1/$`n_e`$ (for a fixed input spectrum), while $`N_{\mathrm{C}^+}/N_\mathrm{S}`$ goes as $`n_e`$. Comparing this ratio for both cloud size assumptions, we find that the cloudlet model results in too low a ratio, $`0.0007`$, whereas the larger clouds yield a ratio of $`0.002`$, close to the observed values. For this reason, and because of the various conceptual difficulties inherent in the cloudlet picture, we have focussed our attention on the more standard cloud picture (i.e. velocity features corresponding to spatially coherent clouds) in our comparisons with the data. It should be noted, however, that lacking the SIII line observations for the slow clouds, we are unable to differentiate between the cloud size models for those clouds.
Up to this point we have not discussed the ionization due to HD 93521 itself. According to Vacca, Garmany & Shull (1996), HD 93521, an O9.5V star, should have a luminosity in ionizing photons of $`2.40\times 10^{48}`$s<sup>-1</sup>. This then implies that the star can create enough ionization to produce $`I^{}(\mathrm{H}\alpha )=9320/r_{\mathrm{pc}}^2`$ Rayleigh from equation (A15), where $`r_{\mathrm{pc}}`$ is the distance of the cloud from the star in pc. We do not have any direct information on the distance between HD 93521 and the cloud nearest to it. HD 93521 is determined to be at a height above the plane of $`z1500`$pc, well above scale height of neutral and even ionized gas. From this we estimate that the contribution of HD 93521 to the H$`\alpha `$ observed towards it is $`10`$%. In any case, since all the clouds observed by SF are optically thick in the EUV, the first cloud along the line of sight will absorb all the ionizing radiation and the other clouds are not affected by HD 93521. We do not include its contribution to the ionization in our calculations.
### 4.3 Model Results
Our model results for both the absorption lines and emission lines, combined into “fast” and “slow” cloud groups are compared with the observations in Table 3. We use results from a series of CLOUDY runs with the same input spectrum (SNR model with $`n_a=0.1`$, $`E_0=10^{51}`$ ergs, no conduction) and pressures ranging from $`P/k=1000`$$`2\times 10^4`$cm$`^3`$K. The model value for the observed column densities and intensities for each cloud is determined by finding which pressure value led to the best match with the observed C II/S II line ratio at the observed value of $`N_{\mathrm{S}^+}`$. We then use the results of the best match pressure run for determining the values for the other observables. We choose to peg our results to the C II/S II line ratio because it is one of the best determined observationally for the line of sight.
Because we fix the other quantities calculated in our models by their values at the observed $`N_{\mathrm{S}^+}`$ values for each cloud, the model and observed values for $`N_{\mathrm{S}^+}`$ are identical. We also match the total line of sight (fast+slow) value for H$`\alpha `$ intensity by scaling the input spectrum as described in §3. The model values for $`N`$(C II) are also very close to the observed values, as a result of our fixing the pressure value by finding the best match for the C II/S II line ratio as described above. Thus the first three lines of the table consititute inputs for the model. For the rest of the table, the model results are truly outputs, not having been fixed to match the data. For these results (as well as for $`N_{\mathrm{C}^+}`$ for completeness) we have presented ranges for the model results corresponding to a (fixed) cloud temperature of 6000 K and 8000 K respectively. This range of values might correspond to the temperatures in the neutral and ionized portions of the clouds as we discuss below.
One of the most confusing aspects of the absorption line data is the apparently large neutral column densities of the clouds, as seen in either H I or S II, coupled with large ionized columns, as inferred from the relatively high columns of $`\mathrm{C}^+`$, which is mainly excited by collisions with electrons in diffuse, $`n0.2`$ cm<sup>-3</sup> gas. This is especially true for the slow clouds. If one interprets each velocity component as being due to a partially ionized cloud, these data require either a very large ionizing flux and/or a high thermal pressure. As we discuss in more detail below, this is because $`N_{\mathrm{C}^+}/N_{\mathrm{H}^0}xn/(1x)`$, where $`xn_{\mathrm{H}^+}/n`$ is the fractional ionization of hydrogen. The observed values of $`N_{\mathrm{C}^+}/N_{\mathrm{H}^0}`$ imply high values for $`x`$ and/or $`n`$ for the clouds at the same time that $`N_{\mathrm{H}^0}`$ is large ($`3\times 10^{18}2\times 10^{19}`$cm<sup>-2</sup>).
The H$`\alpha `$ data taken with WHAM (Hausen et al., 1999), however, provide tight limits on the ionizing flux. For an assumed temperature of 8000 K, their data imply an emission measure of 0.85 cm$`^6`$pc for this line of sight. With the added information from the emission measure ruling out a very large photoionizing flux, we are faced with the prospect of large thermal pressures in the slow clouds. This is demonstrated in figure 2, where we plot the $`N_{\mathrm{C}^+}/N_{\mathrm{S}^+}`$ ratio based on calculations using our model flux. The flux, which is assumed to be the same for all the clouds, has been adjusted to match the total H$`\alpha `$ intensity observed for the line of sight. For our standard model this required reducing the flux by a factor of 0.68. The $`x`$-axis is the total thickness of a cloud, as opposed to depth into a cloud. The points are labelled as in SF but with cloud no. 5 (the smallest and most poorly detected component) excluded. The fast clouds are numbered 1–4 while the slow clouds are 6–9. We are able to match the observed ratios fairly well, although rather high thermal pressures are demanded for the slow clouds, $`P/k75002\times 10^4`$cm$`^3`$K. There is evidence that in regions as diverse as the fairly quiescent hot gas of the Local Bubble and the cold neutral gas at greater distances (Jenkins, Jura & Loewenstein, 1983), high thermal pressures, $`P/k10^42\times 10^4`$cm$`^3`$K, do exist in the ISM. A difficulty with the high value of the inferred pressure, however, is that it could lead to more cooling than can be balanced by the known available heating sources at $`T7000`$ K, which would transform the cloud from the warm phase to the cold phase. We discuss the problem with the heating/cooling balance further below in §6.2.
The reason for the high pressures demanded for the slow clouds can be seen analytically (see also McKee & Slavin, 1999, for a more extensive discussion). Under the assumption of uniform ionization and temperature (as used by SF), the ratio of H$`\alpha `$ intensity to $`\mathrm{C}^+`$ column density gives the hydrogen ionization fraction, $`x`$,
$$\frac{I^{}(\mathrm{H}\alpha )}{N_{\mathrm{C}^+}}=\frac{\alpha _{\mathrm{H}\alpha }(T)A_{21}}{10^6𝒜_{\mathrm{Cg}}\gamma _{12}}x,$$
(16)
where $`\alpha _{\mathrm{H}\alpha }`$ is the effective recombination coefficient for emission of H$`\alpha `$ (see the Appendix), $`𝒜_{\mathrm{Cg}}`$ is the abundance of gaseous carbon, $`\gamma _{12}`$ is rate coefficient for excitation of C<sup>+</sup> to the $`J=3/2`$ level and $`A_{21}`$ is the downward transition probability. $`I^{}(\mathrm{H}\alpha )`$ is assumed to be in Rayleighs and gaseous carbon is assumed to be entirely singly ionized. Applying equation (16) to the data towards HD 93521 (and assuming a temperature of 8000 K) yields $`x=0.37\pm 0.075`$ and $`x=0.14\pm 0.038`$ for the fast and slow clouds, respectively. (Note that the temperature of 8000 K was used here rather than the 6000 K of SF because of results from recent emission line measurements that we discuss below.) This estimate of the ionization can then be used with the observations of the H I column densities to get constraints on the pressure. Assuming that H and He are equally ionized we derive
$$\frac{N_{\mathrm{C}^+}}{N_{\mathrm{H}^0}}\left(\frac{P}{kT}\right)\frac{\gamma _{12}}{A_{21}}𝒜_{\mathrm{Cg}}\left(\frac{x}{1x^2}\right).$$
(17)
Using our results from above this implies (again assuming $`T=8000`$K) $`P/k=1120\pm 340`$ and $`9220\pm 2770`$ for the fast and slow clouds, respectively, where most of the uncertainty derives from the $`20`$% uncertainty in the H$`\alpha `$ intensities. For the more realistic case in which ionization varies with depth into the cloud, as in our model, the value of $`x`$ derived in equation (16) is weighted by $`n_e`$, thereby emphasizing the most highly ionized part of the cloud. As a result, the true mean ionization will be less than that estimated assuming uniform ionization. Equation (17) carries with it a different weighting of $`x`$, introducing more uncertainty in the determination of $`P`$. Again the effect is to emphasize regions of higher ionization resulting in an *underestimate* of the pressure. Despite these limitations we infer that the fast and slow clouds are markedly different in their mean ionization level and their pressures. The slow clouds have high thermal pressures, $`10^4`$ cm$`^3`$K. The fast clouds are more highly ionized and have lower thermal pressures, $`P/k2000`$ cm$`^3`$K, more typical of diffuse interstellar clouds (Jenkins, Jura & Loewenstein, 1983).
As detailed above, we have no difficulty matching the observed H$`\alpha `$ intensities towards HD 93521. The other emission line data, observations of \[S II\] $`\lambda `$6716 and \[N II\] $`\lambda `$6584 (Pifer et al., 1999), are more problematic, however. These data are consistent with the trend recently discovered in the new WHAM data (e.g. Hausen et al., 1999) of both \[S II\]/H$`\alpha `$ and \[N II\]/H$`\alpha `$ to increase to high values as H$`\alpha `$ decreases and as galactic latitude increases. The values for these ratios are a factor of 3 or more larger than is typical for diffuse emission in the galactic plane. The \[N II\]/H$`\alpha `$ ratio is a good temperature diagnostic since the ionization potentials of hydrogen and nitrogen are within 1 eV of each other. The reported values for $`I(6584)/I`$(H$`\alpha `$) of 1 and 1.8 (for the slow and fast clouds respectively) demand high temperatures, $`T8000`$ K and $`T9000`$K. Another temperature diagnostic is the $`N_{\mathrm{C}^+}/I`$(\[S II\]) ratio, because $`N_{\mathrm{C}^+}`$ is relatively insensitive to $`T`$ whereas \[S II\] increases rapidly with rising $`T`$. Since both C<sup>+</sup> and S<sup>+</sup> are expected to be the dominant stages of ionization in WIM/WNM gas, both quantities in the ratio go as $`n_en_\mathrm{H}`$ and the density dependence cancels in the ratio. In contrast to \[N II\]/H$`\alpha `$, however, the observed values for $`N_{\mathrm{C}^+}/I`$(\[S II\]) indicate relatively low temperatures, $`T<6000`$ K for the slow clouds and $`T7000`$ K for the fast clouds. A potential solution to this seeming contradiction is to have substantially lower temperatures in the neutral regions of the clouds relative to the ionized regions. $`I(6584)`$ and $`I`$(H$`\alpha `$) are weighted towards the ionized part of the cloud (both the H<sup>+</sup> and N<sup>+</sup> densities are highest there), while $`N_{\mathrm{C}^+}`$ and $`I`$(\[S II\]) have greater contributions from the neutral part of the clouds. As we noted above, we are unable to model the thermal structure of the WIM/WNM accurately using CLOUDY because it does not yet include PAH heating accurately. Calculations by Wolfire et al (1995a) indicate that temperatures in the WNM range from $`5000`$–9000 K depending on various model parameters including the assumed pressure. Higher pressures lead to lower temperatures which is consistent with the observation that the high pressure, slow clouds have a high $`N_{\mathrm{C}^+}/I`$(\[S II\]) ratio while the lower pressure, fast clouds have a lower $`N_{\mathrm{C}^+}/I`$(\[S II\]) ratio. In summary, it appears that the slow clouds have $`T5000`$ K in the more neutral regions and $`T8000`$ K in the more highly ionized regions, while the fast clouds are characterized by $`T7000`$ K in the more neutral regions to $`T9000`$ K in the more ionized regions.
One additional output from our modeling of the clouds is the filling factor of cloud material. This is more clearly stated in terms of the total line of sight distance occupied by the clouds as determined by their densities and column densities. For our “best fit” models as described above we find that the clouds account for only 100–150 pc of the $`1700`$ pc line of sight towards HD 93521. Thus only 6–9% of the line of sight is occupied by the warm ionized and neutral gas. It is likely that this is not evenly distributed but rather that there are more clouds near the plane. Nevertheless, even assuming the clouds are confined within 1 kpc of the galactic plane their filling factor would only be 9–13%. Thus for this low column density, high latitude line of sight, warm diffuse gas appears to have a low filling factor.
## 5 Helium Ionization and the X-ray Opacity of the Halo
A possible constraint on the nature of the ionizing flux in the halo comes from the recent study of the halo’s X-ray opacity by Arabadjis & Bregman (1999, hereafter A&B). A&B inferred the opacity $`\tau `$ in Galactic gas needed to account for the X-ray spectra of a number of galaxy clusters (including 13 at high latitude) observed with *ROSAT*. In each case, they determined an effective $`\mathrm{H}^0`$ absorption column density $`N_{\mathrm{H},X}`$ based on the assumptions that the absorbing gas is neutral and contains 10% helium by number,
$$\tau N_{\mathrm{H},X}(\sigma _{\mathrm{H}^0}+0.1\sigma _{\mathrm{He}^0}),$$
(18)
where the cross sections, $`\sigma _{\mathrm{H}^0}`$ and $`\sigma _{\mathrm{He}^0}`$ are evaluated at the energy appropriate to the X-ray observations, 0.25 keV. They then compared these column densities with those determined directly from 21 cm observations, $`N_{\mathrm{H}^0}`$ (which they termed $`N_{\mathrm{H},21\mathrm{c}\mathrm{m}}`$). Because the values of $`N_{\mathrm{H},X}`$ make no allowance for ionized gas, which presumably contains some neutral and once ionized helium, one would expect $`N_{\mathrm{H},X}`$ to be substantially larger than the actual $`N_{\mathrm{H}^0}`$. What they found instead was that for the low column, high latitude directions, $`N_{\mathrm{H},X}/N_{\mathrm{H}^0}=0.972\pm 0.022`$.
As A&B pointed out, this result is difficult to understand. The two column densities are related by
$$N_{\mathrm{H},X}(\sigma _{\mathrm{H}^0}+0.1\sigma _{\mathrm{He}^0})=N_{\mathrm{H}^0}\sigma _{\mathrm{H}^0}+N_{\mathrm{He}^0}\sigma _{\mathrm{He}^0}+N_{\mathrm{He}^+}\sigma _{\mathrm{He}^+},$$
(19)
where at the typical energy of the observations (0.25 keV), the cross sections are $`\sigma _{\mathrm{H}^0}=1.08\times 10^{21}`$ cm<sup>-2</sup>, $`\sigma _{\mathrm{He}^0}=2.81\times 10^{20}`$ cm<sup>-2</sup>, and $`\sigma _{\mathrm{He}^+}=2.04\times 10^{20}`$ cm<sup>-2</sup>. Following A&B, we define $`CN_{\mathrm{H},X}/N_{\mathrm{H}^0}`$. Let $`x_{\mathrm{He}}`$ be the helium abundance relative to hydrogen, and let $`\chi _{\mathrm{H}^+}N_{\mathrm{H}^+}/N_H`$, $`\chi _{\mathrm{He}^+}N_{\mathrm{He}^+}/N_{\mathrm{He}}`$, etc, be the fractional ionizations averaged along the line of sight. Inserting the values of the cross sections, we obtain
$$C=0.278+\frac{0.722}{1\chi _{\mathrm{H}^+}}\left(\frac{x_{\mathrm{He}}}{0.1}\right)\left(10.274\chi _{\mathrm{He}^+}\chi _{\mathrm{He}^{++}}\right).$$
(20)
A&B adopted $`\chi _{\mathrm{H}^+}=0.27`$ and $`\chi _{\mathrm{He}^+}=0.5\chi _{\mathrm{H}^+}`$ for the ISM at the solar circle, which gives $`C=1.23`$ for a helium abundance $`x_{\mathrm{He}}=0.1`$. If the helium abundance were $`x_{\mathrm{He}}=0.09`$, as suggested by Baldwin et al (1991) in their analysis of the Orion Nebula, then $`C=1.14`$. Neither value comes close to matching the observed value, $`C=0.972`$.
Our analysis of the line of sight toward HD 93521 gives a significantly lower hydrogen ionization than adopted by A&B, $`\chi _{\mathrm{H}^+}=0.18`$, but a comparable helium ionization, $`\chi _{\mathrm{He}^+}=0.15`$. The column density of doubly ionized helium is negligible, $`\chi _{\mathrm{He}^{++}}=0.0017`$. The difference in the hydrogen ionization fractions is due in part to the fact that the fully ionized HIM does not enter into our analysis, whereas it does contribute to pulsar dispersion measures (Wolfire et al, 1995b) and thus to the value of $`N_{\mathrm{H}^+}`$ used by A&B, as they recognized. Our lower hydrogen ionization gives lower values for $`C`$, 1.12 and 1.04 for $`x_{\mathrm{He}}=0.1`$ and 0.09, respectively. In view of the many possible systematic effects that could enter into determining $`C`$ observationally, we do not regard the discrepancy between our model and A&B’s observation as significant.
## 6 Discussion
### 6.1 The Calculated Spectrum
The main results of this paper involve the ionization caused by cooling hot gas from supernova remnants. Thus questions of the accuracy of the spectra we have produced are of central importance. There are two concerns in this regard: the supernova remnant evolution model and the plasma emission model.
It is clear that our SNR models are highly simplified. In reality, SNRs evolve in a very inhomogenous medium that contains material with a wide range of densities, ionization states, magnetic field strengths and temperatures. In the case of type II and type Ibc supernovae, the stellar progenitor will have shaped the medium into which the remnant evolves. The magnetic field in the ISM can have important influences on the radiative properties of remnants by limiting the compression in the cold shell that is formed during the radiative phase (Slavin & Cox, 1992, 1993). Supernovae in OB associations combine to form superbubbles that can grow to kpc size, sometimes venting their hot gas into the halo. We justify the simplifications we have made by appealing to our results. Despite the drastically different temperature and density profiles that we see for the range of parameter values we explore ($`n_a=0.04`$–1.0 cm<sup>-3</sup>; conduction turned on or off), the time-averaged spectra are not radically different. Over its lifetime, the hot gas radiates its energy primarily in the same lines in all cases, though differences in ionization and temperature shift the balance from one set of lines to another for the different cases. Because we are concerned with the time and volume averaged spectrum, the substantial differences between remnants at any given time in their evolution is smoothed out. Thus, while a strong magnetic field can cause a delay in the radiation from a SNR, we do not expect the time and space averaged spectrum to be substantially different from the spectrum for the no field case.
The reliability of the plasma emission code that we use, the “Raymond & Smith” code, is currently an area of active study. Soon more up-to-date and detailed plasma emission codes will become available (see Brickhouse, 1999). The computed emission spectra and even the total cooling rates of hot plasmas using these codes may differ significantly from those calculated using the R&S code. However, for the reasons discussed in the preceding paragraph, the photoionization rates of the time and space averaged spectra will likely not be substantially changed.
A potentially more important limitation of the calculations we have presented is the lack of inclusion of interface radiation generated in the boundaries between cold and hot gas. Both thermal evaporation (McKee & Cowie, 1977) and turbulent mixing (Slavin, Shull & Begelman, 1993) at those boundaries could lead significant enhancement of cooling in those regions as well as an emission spectrum that is spectrally different from those that we have presented here. It is unclear at this time what fraction of the energy radiated by SNRs could be emitted in such interface regions. This is a worthy subject for future investigations.
### 6.2 Ionization and Thermal Balance in the Clouds
An important difficulty for our modeling of the clouds towards HD 93521 concerns the heating/cooling balance. As we mentioned in our description of our model, we chose to use a constant temperature in the clouds to avoid the sensitivity to gas phase elemental abundances and dust heating rates in our cloud modeling. Nevertheless, we have used CLOUDY, to examine the thermal balance in the gas for some cases of interest. We have found that, while at low pressure ($`P/k=1000`$–2000 cm$`^3`$K) CLOUDY can produce temperatures in the 5000–8000 K range in the ionized portion of the clouds, it has insufficient heating to maintain temperatures above 1000 K in the neutral portions of the cloud. For the low pressure (fast) clouds, we attribute this to the grain photoelectric heating model in CLOUDY. Using the best grain heating rates available, Wolfire et al (1995a) have shown that WNM (i.e, warm, $`7000`$ K, neutral medium) can be maintained by dust and PAH heating at these pressures. However, the slow clouds have large neutral hydrogen column densities, and, from our modeling of the C II, large pressures. Wolfire et al. also showed that the WNM cannot be maintained at pressures above 10<sup>4</sup> cm$`^3`$K, as seems to be demanded for some of the clouds. This difficulty can be ameliorated, but not entirely solved, if the slow clouds are broken up into small cloudlets. The cloudlet model has the advantage that, because of the low column density in an individual cloudlet, the ionization is more uniform. As discussed in §4.3, more uniform ionization decreases the pressure required to produce the observed $`N_{\mathrm{C}^+}/N_{\mathrm{H}^0}`$. However, Eq. (17) deomstrates that even under the extreme assumption of uniform ionization, fairly high pressures are demanded for the slow clouds, $`P/k10^4`$cm$`^3`$K. Some additional heating source seems to be needed to maintain these clouds at their observed temperatures of $`6000`$–8000 K. Reynolds, Haffner & Tufte (1999) have also presented evidence from the WHAM survey for an additional heating mechanism in the halo. Some possibilities for such a heating source include photoelectric heating by small dust grains (Reynolds & Cox, 1992), magnetic reconnection heating (Raymond, 1992; Birk, Lesch & Neukirch, 1998) and dissipation of turbulence (Minter & Balser, 1998).
Another proposed source of ionization for the WIM is radiation from decaying neutrinos (Sciama, 1990). This source was seen to have the advantage that, because neutrinos are essentially uniformly distributed throughout the Galaxy, the ionization rate per unit volume is nearly constant. Thus, clouds with large H I columns could be evenly and partially ionized throughout. There are several difficulties with this, however, that make neutrinos a very unlikely candidate for the ionization of the clouds towards HD 93521. First, we have identified a source of ionization associated with the observed X-ray background that can account for essentially all the ionization observed along this line of sight. While the details in our model can be improved upon by using more realistic models for the evolution of supernova remnants and for the structure of the ISM, it is difficult to escape the conclusion that the Galactic X-ray background has a low energy tail that can produce significant photoionization. Second, the photon energy from the decaying neutrinos is estimated to be only about 13.7 eV (Sciama, 1995), which is inadequate to account for the observed N<sup>+</sup> (Pifer et al., 1999) or $`\mathrm{S}^{++}`$ (Spitzer & Fitzpatrick, 1993). Thus another source capable of creating these ions is needed in any case, and this source will ionize hydrogen as well. Finally, we note that we find, in agreement with Reynolds (1991) and Spitzer & Fitzpatrick (1993), that most of the volume along the line of sight to HD 93521 is apparently empty. In Sciama’s (1997) model, however, much of this volume is filled with gas that is maintained in a fully ionized state by the decaying neutrinos. This gas should have detectable absorption lines, but none were found by Spitzer & Fitzpatrick.
## 7 Conclusions
Ionization in the warm diffuse interstellar medium of our galaxy is substantially influenced by the soft X-ray/EUV emission from cooling hot gas in supernova remnants. In the galactic halo, hot gas emission is especially important, dominating the ionization in some regions. In this paper we have presented calculations of the spectrum from the cooling hot gas and the ionization that results. We have found that:
* the hot gas is capable of producing roughly 50% of the EM observed for a typical line of sight through the galactic disk, and more than enough to account for that observed toward HD 93521;
* our calculated spectrum is consistent with the observed soft X-ray diffuse background at high latitudes;
* the flux we calculate is also consistent with the soft X-ray emission observed from other spiral galaxies;
* we predict a low strength of the He I $`\lambda `$5876 recombination line along the HD 93521 sight line, consistent with low values seen in large beam observations;
* the ionization in our models allows us to match the $`\mathrm{C}^+/\mathrm{S}^+`$ and the $`\mathrm{S}^{++}/\mathrm{S}^+`$ ratios seen with absorption measurements along the HD 93521 line of sight, though for the “slow” clouds large thermal pressures ($`10^42\times 10^4`$cm$`^3`$K) are required;
* by varying the assumed temperature in our clouds from $`60009000`$ K, we are able to match the observed emission line strengths of \[S II\] and \[N II\], and predict the intensity of \[O I\] $`\lambda `$6300;
* with allowance for the uncertainty in the helium abundance and possible systematic effects in the analysis of the data, we regard our value for the X-ray opacity toward HD 93521 as being consistent with the value inferred by Arabadjis & Bregman (1999) for the halo;
* by accounting for the observed ionization toward HD 93521 with a theoretical extrapolation of a known source of ionization, the Galactic X-ray background, we obviate the need to invoke exotic ionization mechanisms such as decaying neutrinos (Sciama, 1990).
We are thankful to Ron Reynolds, Steve Tufte, and Nancy Hausen for sharing their observational results with us prior to publication. We also thank Gary Ferland for discussions regarding CLOUDY and sulphur dielectronic recombination rates. We wish to recognize the contribution of Xander Tielens who participated in our early discussions of this project. CFM gratefully acknowledges support by an NSF grant (AST95-30480), a grant from the Guggenheim Foundation, and a grant from the Sloan Foundation to the Institute for Advanced Study. The authors acknowledge support from the NASA Astrophsical Theory Program.
## Appendix A TRANSFER OF IONIZING RADIATION IN THE HALO
In the absence of scattering, the equation of transfer is
$$\mu \frac{dI_\nu }{d\tau _\nu }=S_\nu I_\nu ,$$
(A1)
in standard notation; recall that $`S_\nu =j_\nu /\kappa _\nu `$ is the source function. We model the disk and halo of a galaxy as having an emissivity $`j_\nu `$ and opacity $`\kappa _\nu `$ that depend only on the distance from the midplane, $`z`$. Solving this equation for the flux, $`F_\nu \mu I_\nu 𝑑\mathrm{\Omega }`$, we find that the flux measured above the halo is
$$F_{\nu +}=F_\nu +4\pi _0^{\tau _{0\nu }}S_\nu 𝑑\tau _\nu 4\pi _0^{\tau _{0\nu }}J_\nu 𝑑\tau _\nu ,$$
(A2)
where $`F_\nu `$ is the flux measured below the halo, $`J_\nu =(1/4\pi )I_\nu 𝑑\mathrm{\Omega }`$ is the mean intensity, and $`\tau _{0\nu }`$ is the total optical depth of the galaxy (measured normal to the plane). Note that $`4\pi S_\nu 𝑑\tau _\nu =4\pi j_\nu 𝑑z`$ is the total emission per unit area. The expected value of $`j_\nu `$ is $`ϵ_\nu /4\pi `$, where the expected volume emissivity $`ϵ_\nu `$ is defined in equation (1). We define the emissivity per unit area as
$$ϵ_{\nu A}_{\mathrm{}}^{\mathrm{}}𝑑zϵ_\nu =4\pi _0^{\tau _{0\nu }}S_\nu 𝑑\tau _\nu .$$
(A3)
The mean escape probability—i.e., the fraction of the radiation that escapes the galaxy—is
$$\eta _\nu =\frac{F_{\nu +}F_\nu }{4\pi {\displaystyle _0^{\tau _{0\nu }}}S_\nu 𝑑\tau _\nu }=\frac{F_{\nu +}F_\nu }{ϵ_{\nu A}}.$$
(A4)
In terms of the mean escape probability, the mean value of the average intensity is
$$\overline{J}_\nu \frac{1}{\tau _{0\nu }}_0^{\tau _{0\nu }}J_\nu 𝑑\tau _\nu =\frac{(1\eta _\nu )}{\tau _{0\nu }}_0^{\tau _{0\nu }}S_\nu 𝑑\tau _\nu $$
(A5)
from equation (A2), so that
$$\overline{J}_\nu =\frac{(1\eta _\nu )}{\tau _{0\nu }}\left(\frac{ϵ_{\nu A}}{4\pi }\right).$$
(A6)
In order to calculate the mean escape probabability, we must make an assumption about the spatial distribution of the emissivity and absorption. We make the simplest assumption possible: The source function is constant ($`j_\nu \kappa _\nu `$). In that case, the flux escaping from the top of the halo is
$$F_{\nu +}=𝑑\mathrm{\Omega }_0^{\tau _{0\nu }}S_\nu e^{\tau _\nu /\mu }𝑑\tau _\nu =2\pi S_\nu _0^1𝑑\mu _0^{\tau _{0\nu }}e^{\tau _\nu /\mu }𝑑\tau _\nu .$$
(A7)
The magnitude of the flux escaping from the bottom of the halo is the same. Evaluating the integrals and using equation (6), we find
$$\eta _\nu =\frac{1}{\tau _{0\nu }}\left[\frac{1}{2}E_3(\tau _{0\nu })\right],$$
(A8)
where $`E_3`$ is an exponential integral. Although we have numerically evaluated $`E_3(\tau _{0\nu })`$ in applying this equation in the text, it is sometimes useful to have an approximate analytic expression. The expression
$$\eta _\nu \frac{1}{1+2\tau _{0\nu }}$$
(A9)
is accurate to within 20%, whereas the corresponding approximation for the absorbed fraction,
$$1\eta _\nu \frac{2\tau _{0\nu }}{1+2\tau _{0\nu }},$$
(A10)
has the same accuracy provided $`\tau _{0\nu }>0.03`$. For ionizing photons with energies somewhat above 1 Ryd, the escape probability is small; as a result, the absorbed fraction is close to unity even if our assumption of a constant source function is badly violated, whereas the escaping fraction is quite sensitive to this assumption. For X-rays, the optical depth of the disk and halo is of order unity, so that the predicted X-ray luminosity of a galaxy is somewhat sensitive to the assumption of a constant source function.
### A.1 Intensity of Recombination Lines
Essentially all the ionizing photons that are absorbed in the disk and halo ionize hydrogen; even those that initially ionize helium lead to the production of a hydrogen ionizing photon (Osterbrock, 1989). If we ignore ionization by secondary electrons, the ionization equilibrium of hydrogen is governed by
$$_{\nu _\mathrm{H}}^{\mathrm{}}𝑑\nu (1\eta _\nu )\frac{ϵ_{\nu A}}{h\nu }=\alpha ^{(2)}_{\mathrm{}}^{\mathrm{}}𝑑zn_en_p=\alpha ^{(2)}EM_{},$$
(A11)
where $`h\nu _\mathrm{H}`$ is the ionization potential of hydrogen, $`\alpha ^{(2)}2.59\times 10^{13}T_4^{0.8}`$ cm<sup>3</sup> s<sup>-1</sup> is the rcombination coefficient to the excited states of hydrogen, and $`EM_{}`$ is the emission measure normal to the plane of the galaxy. Ionization by secondary electrons becomes dominant at large column densities (Maloney, Hollenbach & Tielens, 1996), but for the relatively small column densities encountered in our model ($`N1.5\times 10^{19}`$—see §4), they contribute less than 2% to the emission measure based on calculations with CLOUDY. In fact, the escape probability is small for almost all the ionizing photons ($`\eta _\nu 1`$), so the ionization balance can be approximately described by
$$ϵ_{\mathrm{H}A}^{}=\alpha ^{(2)}EM_{}=7.99\times 10^5T_4^{0.8}EM_{\mathrm{pc}}\mathrm{photons}\mathrm{cm}^2\mathrm{s}^1,$$
(A12)
where
$$ϵ_{\mathrm{H}A}^{}_{\nu _\mathrm{H}}^{\mathrm{}}𝑑\nu \frac{ϵ_{\nu A}}{h\nu }$$
(A13)
is the ionizing photon luminosity per unit area of the galactic disk and the subscript “pc” on the emission measure indicates that it is measured in units of cm<sup>-6</sup> pc.
In the absence of extinction, the H$`\alpha `$ photon intensity in Rayleighs is
$$I^{}(\mathrm{H}\alpha )=\frac{\alpha _{\mathrm{H}\alpha }}{4\pi }EM\left(\frac{1\mathrm{R}}{10^6/4\pi \mathrm{photons}}\right)=0.364T_4^{0.9}EM_{\mathrm{pc}}\mathrm{R},$$
(A14)
where $`\alpha _{\mathrm{H}\alpha }1.18\times 10^{13}T_4^{0.9}`$ cm<sup>3</sup> s<sup>-1</sup> is the effective recombination coefficient for H$`\alpha `$ (from Osterbrock, 1989). The relationship between the H$`\alpha `$ intensity normal to the disk and the surface emissivity (eq. A12) is almost independent of temperature,
$$I_{}^{}(\mathrm{H}\alpha )=4.55\times 10^7T_4^{0.1}ϵ_{\mathrm{H}A}^{}\mathrm{R}.$$
(A15)
These results can be readily extended to the recombination lines of helium, such as He I $`\lambda \mathrm{\hspace{0.17em}5876}`$. Ionization equilibrium in a column normal to the disk gives
$$\zeta _{\mathrm{He}^0A}=y_{\mathrm{He}^0}ϵ_{\mathrm{He}^0A}^{}\alpha _{\mathrm{He}^0}^{(2)}𝑑zn_en(\mathrm{He}^+),$$
(A16)
where $`y_{\mathrm{He}^0}`$ is the fraction of the helium ionizing photons that are absorbed by neutral helium and $`\alpha _{\mathrm{He}^0}^{(2)}2.73\times 10^{13}T_4^{0.67}`$ cm<sup>3</sup> s<sup>-1</sup> is the recombination coefficient to excited states of He<sup>0</sup> (from Osterbrock, 1989). As written, the equation is approximate because, in contrast to hydrogen, not all the recombinations to the ground state result in the ionization of another helium atom. The intensity of the $`\lambda `$ 5876 line is
$$I(5876)=\frac{\alpha _{5876}}{4\pi }𝑑zn_en(\mathrm{He}^+),$$
(A17)
where $`\alpha _{5876}4.9\times 10^{14}T_4^{1.1}`$ cm<sup>3</sup> s<sup>-1</sup> is the effective recombination coefficient for this line (from Osterbrock, 1989). As a result, we find
$$I^{}(5876)=1.8\times 10^7y_{\mathrm{He}^0}T_4^{0.4}ϵ_{\mathrm{He}^0A}^{}\mathrm{R}.$$
(A18)
### A.2 Opacity in a Cloudy Medium
The interstellar medium is highly inhomogeneous. In the simplest idealization, this inhomogeneity can be represented by clouds of opacity $`\tau _c`$ embedded in a transparent intercloud medium. The opacity in a cloudy medium was derived by Bowyer & Field (1969). If the expected number of clouds along the line of sight is $`𝒩_{\mathrm{los}}`$, then the expected value of the ratio of the observed intensity $`I_\nu `$ to the emitted intensity $`I_{\nu 0}`$ is
$`I_\nu /I_{\nu 0}`$ $`=`$ $`\left(1+𝒩_{\mathrm{los}}e^{\tau _{c\nu }}+{\displaystyle \frac{1}{2!}}𝒩_{\mathrm{los}}^2e^{2\tau _{c\nu }}\mathrm{}\right)e^{𝒩_{\mathrm{los}}},`$ (A19)
$`=`$ $`\mathrm{exp}[𝒩_{\mathrm{los}}(1e^{\tau _{c\nu }})],`$ (A20)
where the different terms in the series represent the contribution to the expected intensity if there are no clouds along the line of sight, one cloud, two clouds, etc. Thus, the expected value of the optical depth in a medium composed of identical clouds is
$$\tau _\nu =𝒩_{\mathrm{los}}\left(1e^{\tau _{c\nu }}\right).$$
(A21)
In the limit of small $`\tau _{c\nu }`$, the expected optical depth reduces to $`𝒩_{\mathrm{los}}\tau _{c\nu }`$, the value for a uniform medium. On the other hand, in the limit of large $`\tau _{c\nu }`$, the expected optical depth is $`𝒩_{\mathrm{los}}`$, since the probability that a ray will not encounter any clouds is $`\mathrm{exp}(𝒩_{\mathrm{los}})`$.
If the total HI column density of the disk and halo is $`N_{\mathrm{H}^0}`$ and that of an individual cloud is $`N_{\mathrm{H}^0c}`$, then
$$\tau _{0\nu }=\frac{N_{\mathrm{H}^0}}{N_{\mathrm{H}^0c}}\left(1e^{N_{\mathrm{H}^0c}\sigma _\nu }\right),$$
(A22)
where $`\sigma _\nu `$ is the absorption cross section of the cloud material.
If the clouds are not all identical, the above argument is readily generalized (Bowyer & Field, 1969). The expected opacity is
$$\kappa _\nu =\frac{d\tau _\nu }{ds}=\frac{d𝒩_{\mathrm{los}}}{ds}\left(1e^{\tau _{c\nu }}\right)\frac{1}{\lambda _c}\left(1e^{\tau _{c\nu }}\right),$$
(A23)
where $`ds`$ is an element of distance along the line of sight and $`\lambda _c`$ is the cloud mean free path. In terms of the area of a cloud, $`A_c`$, and the number of clouds per unit volume, $`d𝒩/dV`$, the mean free path is
$$\frac{1}{\lambda _c}=A_c\frac{d𝒩}{dV}.$$
(A24)
If the clouds have a spectrum of masses, then the expected opacity becomes
$$\kappa _\nu =𝑑MA_c(M)\frac{d^2𝒩}{dVdM}\left[1e^{\tau _{c\nu }(M)}\right].$$
(A25)
Finally, we note that the specific intensity calculated with this value of $`\kappa _\nu `$ is approximately equal to that in the intercloud medium, provided the filling factor of the intercloud medium $`f_{ic}`$ is close to unity. The solution of the equation of radiative transfer for a constant source function is
$$I_\nu =S_\nu (1e^{\tau _\nu /\mu }).$$
(A26)
If $`\tau _\nu `$ or $`\tau _{c\nu }`$ is small, there is no distinction between the intensity in the clouds and that in the intercloud medium, so we focus on the case in which both optical depths are large. In that case,
$$I_\nu S_\nu =\frac{j_\nu }{\kappa _\nu }j_\nu \lambda _c.$$
(A27)
In a cloudy medium, $`j_\nu =j_{\nu ,ic}f_{ic}`$, where $`j_{\nu ,ic}`$ is the emissivity in the intercloud medium. If the intercloud filling factor is close to unity, then $`I_\nu j_{\nu ,ic}\lambda _c`$, which is the expected value of the intensity in the intercloud medium when $`\tau _\nu `$ and $`\tau _{c\nu }`$ are large. |
warning/0001/hep-th0001185.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Integrable quantum field theories which interpolate between different conformal field theories have been recognised as an intriguing feature of two-dimensional models ever since the initial paper of A.B. Zamolodchikov . The first examples arose as $`\varphi _{13}`$ perturbations of unitary minimal models $`_{p,p+1}`$. The existence of higher-spin integrals of motion indicates that such perturbations should be integrable; in (see also ) perturbative arguments were used to show that for one sign of the coupling constant the resulting flow is to the neighbouring unitary minimal model, namely $`_{p1,p}`$, at least for $`p1`$. Further support for this picture, this time valid at all values of $`p`$, was provided when ideas from the thermodynamic Bethe ansatz (TBA) enabled exact integral equations describing the flows to be proposed .
Unitarity of a conformal field theory is not a necessary requirement for the existence of integrable perturbations, and in the arguments of were extended to the $`\varphi _{13}`$ perturbations of more general minimal models $`_{p,q}`$. The results indicated that there should exist integrable flows between the models $`_{p,q}`$ and $`_{2pq,p}`$, at least in the region $`pqp`$ where calculations perturbative in $`(qp)/p`$ can be trusted. (It remains an open problem to give a TBA treatment of the nonunitary flows, which would not suffer from this caveat.) The unitary and nonunitary models can be put on a common footing by defining a parameter $`\zeta =p/(qp)`$<sup>*</sup><sup>*</sup>*$`\zeta `$ was denoted $`m`$ in ; the symbol $`p`$ is also often used, but we are reserving this for the integer-valued first index of the minimal models $`_{p,q}`$. The value of $`\zeta `$ suffices to identify $`p`$ and $`q`$ uniquely, since they must be coprime. In all cases the predicted flow is from the model specified by $`\zeta `$ to the model specified by $`\zeta 1`$.
Many further massless flows have since been discovered and studied, often by means of the TBA technique: the papers are a sample of this work. Our main interest here will be the flows found by Martins and further studied by Ravanini et al.. These are perturbations of the minimal models $`_{p,2p1}`$ and $`_{p,2p+1}`$ by the operators $`\varphi _{21}`$ and $`\varphi _{15}`$ respectively, and the predicted trajectories are:
$`\varphi _{21}:_{p,2p1}`$ $``$ $`_{p1,2p1};`$ (1.1)
$`\varphi _{15}:_{p,2p+1}`$ $``$ $`_{p,2p1}.`$ (1.2)
These flows can be chained together to form a single sequence, along which the perturbing operator alternates between $`\varphi _{21}`$ and $`\varphi _{15}`$.
In this paper we will extend this picture by proposing further sequences of flows for which the perturbing operator alternates in the same way. As in studies based on the TBA technique, our main tool will be a conjecture for an exact equation expressing the finite-volume ground state energy of each model as a function of the system size $`R`$. In fact, it will be most convenient to work with an associated scaling function known as the ‘effective central charge’, $`c_{\mathrm{eff}}(r)`$. This is related to the ground state energy as
$$E_0(M,R)=E_{\mathrm{bulk}}(M,R)\frac{\pi }{6R}c_{\mathrm{eff}}(r),r=MR.$$
(1.3)
In this equation, $`E_0`$ is the ground state energy, $`E_{\mathrm{bulk}}`$ the irregular bulk part, and $`M`$ some mass scale, set for example by an infinite-volume one-particle state. The effective central charge depends on the system size only through the dimensionless combination $`MR`$; for a model which happens to be scale invariant (conformal), it is constant and equal to $`c24d`$, where $`c`$ is the model’s central charge and $`d`$ its lowest conformal dimension. The previously-studied sequence (1.1), (1.2) is picked out by the monotonicity of the function $`c_{\mathrm{eff}}`$ as a function of $`r`$. In all other cases, somewhat to our surprise, we found that the effective central charge undergoes a number of oscillations as it interpolates between its short and long distance limits.
Our approach differs from the TBA method, and is much closer to that of the papers , in that a single nonlinear integral equation (NLIE) is proposed to describe infinitely-many different perturbed conformal field theories, each being picked out by an appropriate choice of certain parameters. For massless flows, the one previous example of such an equation was found by Al. Zamolodchikov, and his paper formed a large part of the motivation for our work.
The particular nonlinear integral equation that we will be using is obtained in §2; then in §3 we discuss the nature of the flows that it predicts, and in §4 we take a more detailed look at various asymptotics. This allows us to back up our conjectures with a comparison with UV conformal perturbation theory. In §5 we comment on some features of the $`\varphi _{13}`$ flows, and in the concluding §6 we indicate some open problems that remain for future work.
## 2 The nonlinear integral equation
The work of built in part on the nonlinear integral equation of the scale-invariant form of this equation had arisen previously in the context of integrable lattice models , which encodes the finite-volume ground state energy of the sine-Gordon model for a continuous range of the coupling $`\beta `$, and at general ‘twist’ $`\alpha `$. Setting $`\zeta =\beta ^2/(8\pi \beta ^2)`$, it can be shown that for $`\zeta =p/(qp)`$ and $`\alpha =1/p`$, the ground-state energy of the $`\varphi _{13}`$ perturbation of a general minimal model $`_{p,q}`$ is also matched . (Note that this definition of $`\zeta `$ is thus in line with the one given in the introduction.) Since the sine-Gordon model is massive, the minimal model perturbations reproduced in this way are also massive. In contrast, the modification to the equation found in describes a massless perturbation of the sine-Gordon model, flowing from the model at $`\zeta `$ to the model at $`\zeta 1`$, and thus matching the pattern of massless $`\varphi _{13}`$ flows discussed in the introduction. While only the $`\alpha =0`$ case was treated explicitly in , we will confirm below that with a suitable nonzero value of $`\alpha `$ the massless $`\varphi _{13}`$ perturbations of minimal models are also obtained, with results that agree with the perturbative picture. In the nonunitary cases this yields some genuinely new information, since, as already mentioned, TBA systems giving the exact evolution of the ground-state energy for these models are not known.
However this is not the main purpose of our paper: rather, we would like to perform a similar trick for the $`\varphi _{21}`$ and $`\varphi _{15}`$ perturbations of minimal models. Recently, a nonlinear integral equation describing finite size effects in the $`a_2^{(2)}`$ model was conjectured as for the sine-Gordon model, a scale invariant version of this equation had previously been found in the context of integrable lattice models , and checked to describe the massive $`\varphi _{12}`$ , $`\varphi _{21}`$ and $`\varphi _{15}`$ perturbations of minimal models on the imposition of a suitably-chosen twist. (The connection between the $`a_2^{(2)}`$ model and these perturbations can be understood through quantum group reduction .) Our aim is to modify this equation so as to find out about the massless, interpolating, flows that may also be induced by these perturbations.
The nonlinear integral equation found in is
$$f(\theta )=i\pi \alpha ir\mathrm{sinh}\theta +_{𝒞_1}\phi (\theta \theta ^{})\mathrm{ln}(1+e^{f(\theta ^{})})𝑑\theta ^{}_{𝒞_2}\phi (\theta \theta ^{})\mathrm{ln}(1+e^{f(\theta ^{})})𝑑\theta ^{}$$
(2.1)
with the effective central charge given in terms of its solution as
$$c_{\mathrm{eff}}(r)=\frac{3ir}{\pi ^2}\left(_{𝒞_1}\mathrm{sinh}\theta \mathrm{ln}(1+e^{f(\theta )})d\theta _{𝒞_2}\mathrm{sinh}\theta \mathrm{ln}(1+e^{f(\theta )})d\theta \right).$$
(2.2)
The contours $`𝒞_1`$ and $`𝒞_2`$ run from $`\mathrm{}`$ to $`+\mathrm{}`$, just below and just above the real $`\theta `$-axis, and $`r`$ is equal to $`MR`$, $`R`$ being the size of the system and $`M`$ the mass scale, set by the fundamental kink. The kernel
$$\phi (\theta )=_{\mathrm{}}^{\mathrm{}}\frac{e^{ik\theta }\mathrm{sinh}(\frac{\pi }{3}k)\mathrm{cosh}(\frac{\pi }{6}k(12\xi ))}{\mathrm{cosh}(\frac{\pi }{2}k)\mathrm{sinh}(k\frac{\pi }{3}\xi )}\frac{dk}{2\pi }$$
(2.3)
is equal to $`i/2\pi `$ times the logarithmic derivative of the scalar factor in the S-matrix of the $`a_2^{(2)}`$ (Tzitzéica-Izergin-Korepin-Bullough-Dodd-Zhiber-Mikhailov-Shabat …) model. The parameter $`\xi `$, related to the $`a_2^{(2)}`$ coupling $`\gamma `$ as $`\xi =\gamma /(2\pi \gamma )`$, is the analogue of the sine-Gordon parameter $`\zeta `$ mentioned earlier, while $`\alpha `$ corresponds to the twist. To obtain massive $`\varphi _{12}`$, $`\varphi _{21}`$ and $`\varphi _{15}`$ perturbations of a minimal model $`_{p,q}`$, the values of $`\xi `$ and $`\alpha `$ must be chosen as follows :
$`\varphi _{12}`$ $`:`$ $`\xi ={\displaystyle \frac{1}{\frac{2q}{p}1}},\alpha =2/p,p<2q;`$ (2.4)
$`\varphi _{21}`$ $`:`$ $`\xi ={\displaystyle \frac{1}{\frac{2p}{q}1}},\alpha =2/q,p>q/2;`$ (2.5)
$`\varphi _{15}`$ $`:`$ $`\xi ={\displaystyle \frac{1}{\frac{q}{2p}1}},\alpha =1/p,p<q/2;`$ (2.6)
In each case, the inequality delimits the region of the $`(p,q)`$ plane within which $`\xi `$ is positive, and the corresponding perturbation is relevant. (The first inequality is not strictly necessary, since we will be adopting the convention that $`p<q`$ in the specification of $`_{p,q}`$.) The UV effective central charge, $`c_{\mathrm{eff}}(0)=1\frac{3\xi }{\xi +1}\alpha ^2`$, is always equal to $`16/pq`$, but the subsequent terms in the expansions differ, as expected given the different perturbations being described.
The modification to the massive equation of found in made use of elements of the corresponding massless scattering theory, proposed in . In our case we do not have a description of the massless scattering in terms of an S-matrix but, proceeding by analogy with the results of , we will substitute the single equation (2.1) with two equations describing hypothetical left and right movers. From the formulae (2.5) and (2.6), we notice that the massive versions of the perturbations (1.1) and (1.2) have $`\xi =2p1`$ and $`\xi =2p`$ respectively. We therefore seek a massless modification of the NLIE (2.1) which will interpolate in general between an ultraviolet theory with parameter $`\xi `$ and an infrared theory with parameter $`\xi 1`$. To this end, we introduce two analytic functions $`f_R(\theta )`$ and $`f_L(\theta )`$, couple them together via
$`f_R(\theta )`$ $`=`$ $`i{\displaystyle \frac{r}{2}}e^\theta +i\pi \alpha ^{}`$ (2.7)
$`+{\displaystyle _{𝒞_1}}\varphi (\theta \theta ^{})\mathrm{ln}(1+e^{f_R(\theta ^{})})𝑑\theta ^{}{\displaystyle _{𝒞_2}}\varphi (\theta \theta ^{})\mathrm{ln}(1+e^{f_R(\theta ^{})})𝑑\theta ^{}`$
$`+{\displaystyle _{𝒞_1}}\chi (\theta \theta ^{})\mathrm{ln}(1+e^{f_L(\theta ^{})})𝑑\theta ^{}{\displaystyle _{𝒞_2}}\chi (\theta \theta ^{})\mathrm{ln}(1+e^{f_L(\theta ^{})})𝑑\theta ^{}`$
$`f_L(\theta )`$ $`=`$ $`i{\displaystyle \frac{r}{2}}e^\theta i\pi \alpha ^{}`$ (2.8)
$`+{\displaystyle _{𝒞_2}}\varphi (\theta \theta ^{})\mathrm{ln}(1+e^{f_L(\theta ^{})})𝑑\theta ^{}{\displaystyle _{𝒞_1}}\varphi (\theta \theta ^{})\mathrm{ln}(1+e^{f_L(\theta ^{})})𝑑\theta ^{}`$
$`+{\displaystyle _{𝒞_2}}\chi (\theta \theta ^{})\mathrm{ln}(1+e^{f_R(\theta ^{})})𝑑\theta ^{}{\displaystyle _{𝒞_1}}\chi (\theta \theta ^{})\mathrm{ln}(1+e^{f_R(\theta ^{})})𝑑\theta ^{}`$
and replace the expression (2.2) for the effective central charge with
$`c_{\mathrm{eff}}(r)`$ $`=`$ $`{\displaystyle \frac{3ir}{2\pi ^2}}[{\displaystyle _{𝒞_1}}e^\theta \mathrm{ln}(1+e^{f_R(\theta )})d\theta {\displaystyle _{𝒞_2}}e^\theta \mathrm{ln}(1+e^{f_R(\theta )})d\theta `$ (2.9)
$`+{\displaystyle _{𝒞_2}}e^\theta \mathrm{ln}(1+e^{f_L(\theta )})d\theta {\displaystyle _{𝒞_1}}e^\theta \mathrm{ln}(1+e^{f_L(\theta )})d\theta ].`$
As before, the contours $`𝒞_1`$ and $`𝒞_2`$ run from $`\mathrm{}`$ to $`+\mathrm{}`$ just above and just below the real axis, and $`r=MR`$. However, since the flows are massless $`M`$ no longer has a direct interpretation as the mass of an asymptotic state, but rather sets the crossover scale.
In the far infrared, $`r\mathrm{}`$ and the two equations (2.7) and (2.8) decouple. The result is two copies of the UV limit of (2.1), with the kernel $`\phi (\theta )`$ substituted by $`\varphi (\theta )`$. Since the IR destination is to be the model at $`\xi 1`$, we take $`\varphi (\theta )`$ to be the massive kernel (2.3) with the substitution $`\xi \xi 1`$:
$$\varphi (\theta )=_{\mathrm{}}^{\mathrm{}}\frac{e^{ik\theta }\mathrm{sinh}(k\frac{\pi }{3})\mathrm{cosh}(\frac{\pi }{6}k(32\xi ))}{\mathrm{cosh}(\frac{\pi }{2}k)\mathrm{sinh}(k\frac{\pi }{3}(\xi 1))}\frac{dk}{2\pi }.$$
(2.10)
To find the other kernel, $`\chi (\theta )`$, we consider the UV behaviour of the new system. This should coincide with the UV limit of the massive system (2.1). As $`r0`$, the solution to (2.7), (2.8) splits in the standard way into a pair of kink systems, centered at $`\theta =\pm \mathrm{ln}(1/r)`$. We focus on one of them by replacing $`\theta `$ by $`\theta \mathrm{ln}(1/r)`$, and keeping this variable finite as the limit is taken. The exponential term in (2.7) is replaced by $`\frac{i}{2}e^\theta `$, while to leading order the exponential in (2.8) can be neglected. Furthermore, the singularities in $`\mathrm{ln}(1+e^{\pm f_L(\theta )})`$ previously found on the real axis are pushed to $`\mathrm{}`$, allowing integration contours in any integral involving $`f_L`$ to be moved across the real $`\theta `$ axis at will. It will be convenient to shift $`𝒞_2`$ down to coincide with $`𝒞_1`$ in all such integrals, and to take $`𝒞_1`$ to be the line $`\mathrm{}m\theta =\eta `$, with $`\eta `$ a suitably-small real number.
For $`f_R(\theta )`$, the singularities do not disappear but using the reality properties of $`f_R`$ we can at least rewrite the integrations above and below the axis in terms of the imaginary part of a single integration along $`𝒞_1`$. The final equation is
$`f_R(\theta )`$ $`=`$ $`{\displaystyle \frac{i}{2}}e^\theta +i\pi \alpha ^{}+2i{\displaystyle _{\mathrm{}}^{\mathrm{}}}\varphi (\theta (\theta ^{}i\eta ))\mathrm{}m(\mathrm{ln}(1+e^{f_R(\theta ^{}i\eta )}))𝑑\theta ^{}`$ (2.11)
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\chi (\theta (\theta ^{}i\eta ))f_L(\theta ^{}i\eta )𝑑\theta ^{}`$
$`f_L(\theta )`$ $`=`$ $`i\pi \alpha ^{}+{\displaystyle _{\mathrm{}}^{\mathrm{}}}\varphi (\theta (\theta ^{}+i\eta ))f_L(\theta ^{}+i\eta )𝑑\theta ^{}`$ (2.12)
$`2i{\displaystyle _{\mathrm{}}^{\mathrm{}}}\chi (\theta (\theta ^{}i\eta ))\mathrm{}m(\mathrm{ln}(1+e^{f_R(\theta ^{}i\eta )}))𝑑\theta ^{}.`$
Taking the Fourier transform of (2.12) yields
$$(1\stackrel{~}{\varphi }(k))\stackrel{~}{f}_L(k)=2i\pi ^2\alpha ^{}\delta (k)2i\stackrel{~}{\chi }(k)\stackrel{~}{L}_R(k)$$
(2.13)
where $`\stackrel{~}{L}_R(k)`$ denotes the transform of $`\mathrm{}m(\mathrm{ln}(1+e^{f_R(\theta ^{}i\eta )}))`$ and our convention for the Fourier transform of $`f(\theta )`$ is
$$[f(\theta )]=_{\mathrm{}}^{\mathrm{}}f(\theta )e^{i\theta k}𝑑\theta =\stackrel{~}{f}(k)$$
(2.14)
with corresponding inverse
$$^1[\stackrel{~}{f}(k)]=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}\stackrel{~}{f}(k)e^{i\theta k}𝑑k=f(\theta ).$$
(2.15)
Inserting (2.13) into the F.T. of (2.11) (first taking $`\frac{i}{2}e^\theta `$ to the left hand side to ensure the existence of the transform), we obtain
$$[f_R(\theta )+\frac{i}{2}e^\theta ]=2i\pi ^2\alpha ^{}\delta (k)\left(1+\frac{\stackrel{~}{\chi }(k)}{1\stackrel{~}{\varphi }(k)}\right)2i\stackrel{~}{L}_R(k)\left(\stackrel{~}{\varphi }(k)+\frac{\stackrel{~}{\chi }(k)^2}{1\stackrel{~}{\varphi }(k)}\right).$$
(2.16)
This can be compared with the Fourier transform of the kink limit of (2.1), which is
$$[f(\theta )\frac{i}{2}e^\theta ]=2i\pi ^2\alpha \delta (k)2i\stackrel{~}{L}_R(k)\stackrel{~}{\phi }(k).$$
(2.17)
The two will match if
$$\stackrel{~}{\phi }(k)=\stackrel{~}{\varphi }(k)+\frac{\stackrel{~}{\chi }(k)^2}{1\stackrel{~}{\varphi }(k)}$$
(2.18)
and
$$\alpha =\alpha ^{}\left(1+\frac{\stackrel{~}{\chi }(0)}{1\stackrel{~}{\varphi }(0)}\right).$$
(2.19)
After some elementary manipulations we can invert the Fourier transform of $`\chi `$, resulting in
$$\chi (\theta )=_{\mathrm{}}^{\mathrm{}}\frac{e^{ik\theta }\mathrm{sinh}(\frac{\pi k}{3})\mathrm{cosh}(\frac{\pi k}{6})}{\mathrm{cosh}(\frac{\pi k}{2})\mathrm{sinh}(\frac{\pi k}{3}(\xi 1))}\frac{dk}{2\pi }$$
(2.20)
and, from (2.19),
$$\alpha ^{}=\alpha \frac{\xi }{(\xi 1)}.$$
(2.21)
Since (2.18) involves only $`\chi (k)^2`$ we have made a choice of sign when taking the square root. We took the negative sign since only then do the values taken by $`c_{\mathrm{eff}}(\mathrm{})`$ when flowing from a minimal model $`_{p,q}`$ always assume the form $`16/p^{}q^{}`$, as has to be the case if the infrared destination of the flow is to be another minimal model. It is possible that the flows found for the other sign choice also have an interpretation in some other context, but we will not explore this question here.
## 3 The flows
By construction, the new massless system exhibits the same ultraviolet central charge as the corresponding massive system:
$$c_{\mathrm{eff}}(0)=1\frac{3\xi }{(\xi +1)}\alpha ^2.$$
(3.1)
Furthermore, the fact that the kink limit of the new system can be mapped exactly onto that of the massive system makes it natural to suppose that the recipe for choosing $`\xi `$ and $`\alpha `$ given by (2.4), (2.5) and (2.6) also holds for the massless perturbations. We take this as a working hypothesis for now; in the next section it will be subject to some more detailed checks. However, one immediate difficulty should be mentioned: with our convention that the coprime pair $`(p,q)`$ specifying $`_{p,q}`$ has $`p<q`$, the value assigned to $`\xi `$ by (2.4) for $`\varphi _{12}`$ perturbations is less than $`1`$. The same is true for $`\varphi _{15}`$ perturbations of $`_{p,q}`$ whenever $`q>4p`$. This is not a problem for the massive perturbations, but the massless kernel $`\varphi (\theta )`$ has a pole at $`\theta =\frac{2}{3}(\xi 1)\pi i`$, which crosses the real $`\theta `$ axis as $`\xi `$ dips below $`1`$. In addition, the formula (2.21) clearly has a pole at $`\xi =1`$. All of this leads to technical complications, and, while they could perhaps be overcome by analytic continuation, we will leave the issue to one side for now. In the rest of this paper we will focus solely on the $`\varphi _{21}`$ and $`\varphi _{15}`$ perturbations, found via equations (2.5) and (2.6) respectively, and restrict the $`\varphi _{15}`$ perturbations to models $`_{p,q}`$ with $`q<4p`$.
In the infrared, the massless system mimics a massive system with $`\xi `$ replaced by $`\xi 1`$ and $`\alpha `$ replaced by $`\alpha ^{}=\alpha \xi /(\xi 1)`$, and so
$$c_{\mathrm{eff}}(\mathrm{})=1\frac{3(\xi 1)}{\xi }(\alpha ^{})^2=1\frac{3\xi }{(\xi 1)}\alpha ^2.$$
(3.2)
Consider first the perturbation of $`_{p,q}`$ by $`\varphi _{21}`$. For it to be relevant, we must impose that $`p>q/2`$. Substituting the requisite values of $`\xi `$ and $`\alpha `$ (taken from (2.5)) into (3.1) and (3.2) we find:
$$_{p,q}+\varphi _{21}:p>q/2,c_{\mathrm{eff}}(0)=1\frac{6}{pq},c_{\mathrm{eff}}(\mathrm{})=1\frac{6}{(qp)q}.$$
(3.3)
Similarly, for the $`\varphi _{15}`$ perturbations we have
$$_{p,q}+\varphi _{15}:p<q/2,c_{\mathrm{eff}}(0)=1\frac{6}{pq},c_{\mathrm{eff}}(\mathrm{})=1\frac{6}{p(4pq)}.$$
(3.4)
As promised at the end of the last section, the infrared limiting values of the effective central charges are all consistent with the destinations of the flows being minimal models<sup>§</sup><sup>§</sup>§had the opposite sign choice for $`\chi (\theta )`$ been taken in (2.20), the values of $`c_{\mathrm{eff}}(\mathrm{})`$ for the $`\varphi _{21}`$ and $`\varphi _{15}`$ perturbations would have been $`16(qp)/(qp^2)`$ and $`16(4pq)/(pq^2)`$ respectively. Unfortunately, a knowledge of $`c_{\mathrm{eff}}`$ alone is generally not enough to identify a nonunitary minimal model uniquely, so the destinations cannot be completely pinned down by this information. Moreover, there are cases in which even the knowledge of the full $`c_{\mathrm{eff}}(r)`$ is not sufficient to identify the model or models involved. The possibility of such an ambiguity was first noted, for massive perturbations, in . It goes by the name of the ‘type II conjecture’, and explicitly it equates the effective central charges of the following models:
$$_{p,q}+\varphi _{15}_{q/2,2p}+\varphi _{21}_{2p,q/2}+\varphi _{12}(p=2n+1,q=2m).$$
(3.5)
The second equivalence is just a question of labelling, but the first is much less trivial. However, all follow from the massive NLIE (2.1) and the recipe (2.4)–(2.6), assuming the correctness of these equations (see ). Clearly the massless theories, whose NLIEs are parametrised according to essentially the same recipe, will suffer similar ambiguities; some examples will be encountered later.
In spite of these provisos, the forms taken by the results (3.3) and (3.4) make for some obvious conjectures, and these will receive further support, both analytical and numerical, later in the paper. Such checks will never resolve type II ambiguities, so for these cases the conjectures are better motivated by the belief that the pattern observed in other cases should hold in complete generality. Taking into account the ‘$`p<q`$’ labeling convention for a minimal model $`_{p,q}`$, the $`\varphi _{15}`$ case splits into two and we have the following predictions for massless perturbations of general minimal models:
$`_{p,q}+\varphi _{21}`$ $``$ $`_{qp,q}(p<q<2p),`$ (3.6)
$`_{p,q}+\varphi _{15}`$ $``$ $`_{p,4pq}(2p<q<3p),`$ (3.7)
$`_{p,q}+\varphi _{15}`$ $``$ $`_{4pq,p}(3p<q<4p).`$ (3.8)
There is only a single flow on this list from any given minimal model: when $`\varphi _{21}`$ is relevant, $`\varphi _{15}`$ is irrelevant, and vice versa. Note also that the case (3.8) can only occur as the last member of a sequence: for $`3p<q<4p`$, $`\xi `$ lies between $`1`$ and $`2`$ and the next step would therefore involve a value of $`\xi `$ less than $`1`$, and we have already decided to exclude such cases.
The previously-known flows (1.1) and (1.2) are reproduced on setting $`q=2p1`$ in (3.6) and $`q=2p+1`$ in (3.7). In order to understand the more general pattern, it is convenient to define an an ‘index’ $`I=2pq`$, and to rephrase (3.6) and (3.7) as
$`_{p,2pI}+\varphi _{21}`$ $``$ $`_{pI,2pI},(\xi ,\alpha ^{})=(\frac{2p}{I}1,\frac{1}{pI}),`$ (3.9)
$`_{p,2p+I}+\varphi _{15}`$ $``$ $`_{p,2pI},(\xi ,\alpha ^{})=(\frac{2p}{I},\frac{2}{2pI}).`$ (3.10)
For reference we have included the values of $`\xi `$ and $`\alpha ^{}`$ that should be used in the NLIE in order to dial up the corresponding flow. In all cases $`I(\mathrm{IR})=I(\mathrm{UV})`$; the flows (1.1) and (1.2) make up the unique sequence with $`|I|=1`$. For $`|I|>1`$, there may be more than one sequence, with flows sharing the same value of $`|I|`$ interlacing each other. This is reminiscent of the generalised staircase models studied in , but is a little more complicated owing to the constraint that the pair of integers labelling a minimal model must always be coprime. The number of different sequences with index $`\pm I`$ is therefore given by the Euler $`\phi `$-function $`\phi (|I|)`$, equal to the number of integers less than $`|I|`$ which are coprime to $`|I|`$ (so for $`n=1\mathrm{}6`$, $`\phi (n)=1,1,2,2,4,2`$ ). The index measures the distance from the line $`q=2p`$ across which the relevance and irrelevance of the fields $`\varphi _{21}`$ and $`\varphi _{15}`$ swap over. All of this is perhaps best seen pictorially, and in figure 1 some of the predicted flows are plotted, superimposed on a grid of the minimal models $`_{p,q}`$. The horizontal arrows of length $`|I|`$ correspond to $`\varphi _{21}`$ perturbations, while the vertical arrows, of length $`2|I|`$, are $`\varphi _{15}`$ perturbations.
As a first check on our results, we evaluated the effective central charge numerically for a number of members of the $`|I|=1`$ series, and made a comparison with the results from the massless TBA equations discussed in . In such equations were written in a ‘universal’ form of the kind first described in , but for numerical work it is more convenient to write the equations in the following way:
$$\epsilon _a(\theta )=\nu _a(\theta )\underset{b=1}{\overset{n}{}}(l_{ab}^{(A_n)}\delta _{ab})_{\mathrm{}}^{\mathrm{}}𝒦(\theta \theta ^{})\mathrm{ln}(1+e^{\epsilon _b(\theta ^{})})𝑑\theta ^{},$$
(3.11)
$$c_{\mathrm{eff}}(r)=\frac{3}{\pi ^2}\underset{a=1}{\overset{n}{}}_{\mathrm{}}^{\mathrm{}}\nu _a(\theta )\mathrm{ln}(1+e^{\epsilon _a(\theta )})𝑑\theta $$
(3.12)
where $`n\xi 3`$ is integer, $`l_{ab}^{(A_n)}`$ is the incidence matrix of the $`A_n`$ Dynkin diagram,
$$\nu _a(\theta )=\frac{r}{2}(e^\theta \delta _{a,1}+e^\theta \delta _{a,n}),$$
(3.13)
and
$$𝒦(\theta )=\frac{\sqrt{3}}{\pi }\frac{\mathrm{sinh}2\theta }{\mathrm{sinh}3\theta }.$$
(3.14)
The agreement between the NLIE and the TBA is extremely good, and is illustrated in table 1 and figure 2.
Satisfied that our NLIE correctly matches the previously known flows, we can turn to the new families with index $`|I|>1`$. Figure 3 shows the (unique) series with $`|I|=2`$, and figure 4 one of the two possible series with $`|I|=3`$. For all of these flows the effective central charge initially increases from its UV value, oscillates as the system size is increased, before finally flowing to the predicted IR fixed point. This perhaps-surprising behaviour is in contrast to the $`|I|=1`$ series where all of the flows are monotonic. We will return to this point at the end of section 4.
Note that no TBA equations are known for the massless flows with $`|I|>1`$. This puzzle can be highlighted by looking at some cases where the corresponding massive TBA system can be conjectured.
We consider three families of massive systems found in the ‘ADET’ class of models classified in . The first set is obtained from (3.11) simply by replacing the driving term (3.13) by
$$\nu _a(\theta )=r\delta _{a,1}\mathrm{cosh}\theta .$$
(3.15)
This gives the massive flows corresponding to the $`|I|=1`$ massless flows. For the second set, we continue to use the driving term (3.15), but replace $`l_{ab}^{(A_n)}`$ with $`l_{ab}^{(T_n)}`$ (the incidence matrix of the ‘tadpole’ graph $`T_n=A_{2n}/_2`$), letting $`1`$ label the node furthest from the ‘tadpole’ node (see figure 5). In these systems were identified with the models $`_{n+2,2n+2}+\varphi _{21}`$ for $`n`$ odd and $`_{n+1,2n+4}+\varphi _{15}`$ for $`n`$ even. These are precisely the models in the $`|I|=2`$ series. The massless NLIE predicts, in addition to these previously-known massive flows, the existence of interpolating flows ($`nn1`$) within this family.
The final family of massive TBA equations is obtained by replacing the $`A_n`$ incidence matrix with the $`D_n`$ one. For these cases the type II conjecture mentioned above plays a rôle, and for each TBA system there are two possible identifications:
$$A_n^+_{2n1,2+4n}+\varphi _{15}\mathrm{and}A_n^{}_{1+2n,4n2}+\varphi _{21}.$$
(3.16)
The sets $`A_n^\pm `$ correspond, in the ultraviolet, to the two series of models with index $`|I|=4`$. Turning to the massless flows implied by the NLIE, there is a further ambiguity in the infrared destinations. However, the general pattern of (3.6), (3.7) and (4.6)–(4.8) suggests
$$A_n^\pm A_{n1}^{}.$$
(3.17)
We have checked the low-lying members of each family of massive TBA equations against the massive NLIE (2.1), finding agreement to our numerical accuracy (about 14 digits) in each case. However, returning to the question of finding massless versions of these equations, we observe a key difference between the first family and the other two: as is clear from figure 5, the graphs for the $`|I|=1`$ systems have a $`_2`$ symmetry which is in general absent from those for $`|I|=2`$ or $`4`$ (the single exception occurs when $`|I|=4`$, $`n=4`$). The trick that allowed us to move between massive and massless systems for $`|I|=1`$ by swapping the $`\nu `$’s defined in (3.15) for those defined in (3.13) relied completely on this symmetry. This ‘$`_2`$ trick’ was first employed in ; more generally, the nodes ‘$`1`$’ and ‘$`n`$’ could be replaced by any pair of nodes related by a $`_2`$ graph symmetry. To the best of our knowledge, all massless TBA systems that have been discovered to date are related to massive systems in essentially this way. Thus if the $`T_n`$ and $`D_n`$ TBA systems do have associated massless versions, they are likely to be of a somewhat different nature to all previously-encountered examples.
The $`|I|=4`$, $`n=4`$ system might appear to offer a counterexample, since the exceptional symmetry of its graph does allow a massless version of the massive TBA to be constructed. However this equation turns out not to reproduce the massless NLIE. For example, for the massless TBA $`c_{\mathrm{eff}}(\mathrm{})=5/7`$, which identifies the IR destination of the flow as $`_{3,7}`$ rather than the $`_{7,10}`$ or $`_{5,14}`$ found by the NLIE. Furthermore, the short-distance expansion of $`c_{\mathrm{eff}}(r)`$ turns out to be inconsistent with the massless TBA flow being related to the massive one by an analytic continuation of the coupling $`\lambda `$ – rather, it is the flow produced by the massless NLIE which has this property. This failure of the $`_2`$ trick to produce the analytically-continued massless flow can be put into a more general context. Recall that TBA equations have associated Y-systems, and that these entail a periodicity $`Y_a(\theta )=Y_a(\theta +iP)`$ for certain functions $`Y_a(\theta )`$, where $`P`$ depends on the particular Y-system (see and also ). Suppose that a diagram symmetry relates nodes $`a`$ and $`\stackrel{~}{a}`$ for this system. Then it can be argued that the associated massive and massless TBA equations will be related by analytic continuation if $`Y_a(\theta )=Y_{\stackrel{~}{a}}(\theta +iP/2)`$, and not otherwise. In particular, for systems related to the ADET diagrams, such a property holds if and only if the nodes $`a`$ and $`\stackrel{~}{a}`$ are ‘conjugate’, where conjugation acts on TBA diagrams in the same way as charge conjugation acts on the particles in an affine Toda field theory . For the $`D_n`$ diagrams, the fork nodes are related by charge conjugation for $`n`$ odd, but not for $`n`$ even. This thus matches the observation above that the massive and massless $`D_4`$-related systems are not related by continuation in $`\lambda `$; it also helps to understand from the TBA point of view why the models $`H_N^{(\pi )}`$ of are related to $`H_N^{(0)}`$ by analytic continuation for $`N`$ odd, but not for $`N`$ even.
We conclude this section with some further observations. Following we note that two models in the $`|I|=2`$ set possess $`N=1`$ supersymmetry in the ultraviolet. In the notation of we have
$$S_{2,8}+\widehat{\varphi }_{13}^{top}_{3,8}+\varphi _{15},$$
(3.18)
and
$$S_{3,7}+\widehat{\varphi }_{15}^{top}_{7,12}+\varphi _{21}.$$
(3.19)
The operators $`\widehat{\varphi }_{13}^{top}`$ and $`\widehat{\varphi }_{15}^{top}`$ are SUSY-preserving, so the massless flows $`_{3,8}_{3,4}`$ and $`_{7,12}_{5,12}`$ should exhibit spontaneous breaking of $`N=1`$ SUSY. The situation is reminiscent of the flow between the tricritical Ising and Ising models discussed in . In particular, the flow $`_{3,8}_{3,4}`$ is to a theory of a single massless free Majorana fermion, just as was the case in . This particle, in turn, was identified in as the massless Goldstone fermion. Finally, we mention that that the massive TBA system for $`_{3,8}+\varphi _{15}`$ was alternatively obtained in as a folding of the $`N=2`$ supersymmetric model with spontaneously-broken $`_k`$ symmetry of , for $`k=3`$. This suggests that there should also be a (probably nonunitary) flow from the ‘parent’ $`N=2`$ supersymmetric theory into a massless free fermionic theory. Given that a similar interpolating phenomenon is also present for $`k=2`$for more on the $`k=2`$ case, and its connection with the theory of dense polymers, see it is natural to suppose that the same behaviour will be exhibited for all of the $`_k`$-related models of . It would be interesting to check this conjecture using the $`a_{k1}^{(1)}`$-generalisations of our NLIE.
## 4 Perturbation theory
The claims of the last section can be put on a more secure footing by taking a closer look at the behaviour of the effective central charges at small and large $`r`$. For this we will need the leading asymptotics of the kernels $`\varphi (\theta )`$ and $`\chi (\theta )`$ as $`\theta \mathrm{}`$. That of $`\varphi `$ can be found from the residues of the poles in the integrand of (2.10) at $`k=i`$ and $`k=3i/(\xi 1)`$, and is:
$$\varphi (\theta )\frac{\sqrt{3}\mathrm{sin}(\frac{\pi }{3}\xi )}{\pi \mathrm{sin}(\frac{\pi }{3}(\xi 1))}e^\theta +\frac{3\mathrm{sin}(\frac{\pi }{(\xi 1)})\mathrm{cos}(\frac{\pi }{2(\xi 1)})}{\pi (\xi 1)\mathrm{cos}(\frac{3\pi }{2(\xi 1)})}e^{3\theta /(\xi 1)}+\mathrm{},\theta \mathrm{}.$$
(4.1)
Similarly,
$$\chi (\theta )\frac{3}{2\pi \mathrm{sin}(\frac{\pi }{3}(\xi 1))}e^\theta \frac{3\mathrm{sin}(\frac{\pi }{(\xi 1)})\mathrm{cos}(\frac{\pi }{2(\xi 1)})}{\pi (\xi 1)\mathrm{cos}(\frac{3\pi }{2(\xi 1)})}e^{3\theta /(\xi 1)}+\mathrm{},\theta \mathrm{}.$$
(4.2)
We first analyse the NLIE as $`r\mathrm{}`$, to see whether its behaviour is compatible with the claimed infrared destinations of the massless flows. In spite of the fact that conformal perturbation theory about the infrared fixed point is not renormalisable, there are a number of unambiguous predictions against which the equation can be checked. The key ideas are set out in , and are further discussed in, for example, .
In general, the infrared model will be described by an action of the form
$$S=S_{\mathrm{IR}}^{}+\mu _1\psi d^2x+\mu _2T\overline{T}d^2x+(\mathrm{further}\mathrm{terms})$$
(4.3)
where $`S_{\mathrm{IR}}^{}`$ is the action, and $`\psi `$ one of the (irrelevant) primary fields, of the infrared conformal field theory. $`\psi `$ might be absent, in which case the only operators attracting the flow to the IR fixed point would be the descendents of the identity, $`T\overline{T}`$ being the least irrelevant example. On dimensional grounds, the couplings $`\mu _1`$ and $`\mu _2`$ are related to the single crossover scale $`M`$ as
$$\mu _1=\kappa _1M^{22h},\mu _2=\kappa _2M^2,$$
(4.4)
with $`\kappa _1`$ and $`\kappa _2`$ dimensionless constants and $`h`$ the conformal dimension of $`\psi `$.
Standard methods of perturbed conformal field theory can now be used to calculate the first corrections to $`c_{\mathrm{eff}}(\mathrm{})`$, with results that are good at least up to the order at which the ‘further terms’ cut in. The first perturbing term in (4.3) yields a series in $`\mu _1R^{22h}=\kappa _1r^{22h}`$ if all powers of $`\mu _1`$ contribute, or $`\mu _1^2R^{44h}=\kappa _1^2r^{44h}`$ if only even powers appear, as happens when $`\psi `$ is odd under some symmetry of $`S_{\mathrm{IR}}^{}`$. The second term always contributes a series in $`\mu _2R^2=\kappa _2r^2`$, the first two terms of which were found explicitly in . Putting everything together gives $`c_{\mathrm{eff}}(r)`$ the following large-$`r`$ expansion:
$`c_{\mathrm{eff}}(r)`$ $``$ $`c_{\mathrm{eff}}(\mathrm{})+(\text{a series in }\kappa _1^{}r^{22h}\text{ or }\kappa _1^2r^{44h})`$ (4.5)
$`{\displaystyle \frac{\pi ^3c_{\mathrm{eff}}(\mathrm{})^2}{6}}\kappa _2^{}r^2+{\displaystyle \frac{\pi ^6c_{\mathrm{eff}}(\mathrm{})^3}{18}}\kappa _2^2r^4+(\mathrm{further}\mathrm{terms}).`$
In contrast to the UV situation to be discussed shortly, this series is only expected to be asymptotic, and furthermore there is very little control over the omitted ‘further terms’. Nevertheless, it allows for some useful comparisons with results from the NLIE.
When $`r`$ is large, the nontrivial behaviours of the functions $`f_L`$ and $`f_R`$ appearing in (2.7) and (2.8) are concentrated in ‘kink’ regions near $`\theta =\mathrm{ln}r`$ and $`\theta =\mathrm{ln}r`$ respectively. So long as the contours $`𝒞_1`$ and $`𝒞_2`$ are kept a finite distance from the real $`\theta `$-axis, the functions $`\mathrm{ln}(1+e^{\pm f_L})`$ and $`\mathrm{ln}(1+e^{\pm f_R})`$ appearing in these equations are doubly-exponentially suppressed in the central zone between the kinks, and so the principal interaction between equations (2.7) and (2.8), and hence the principal correction to $`c_{\mathrm{eff}}(r)`$, comes from the exponential tail of the kernel function $`\chi (\theta )`$, given by (4.2). Since we are only worrying about the first few terms in the expansion, and we are free to vary $`\xi `$ so as to avoid any ‘resonance’ effects, we can discuss effects of the two terms in (4.2) separately. Consider first the second term, decaying as $`e^{3\theta /(\xi 1)}`$ as $`\theta \mathrm{}`$. Inserted into (2.7) and considered iteratively, corrections to $`f_R(\theta )`$ as a series in $`r^{6/(\xi 1)}`$ will be generated. Feeding through into $`c_{\mathrm{eff}}(r)`$, these corrections can be matched against the $`\kappa _1`$ terms in (4.5), allowing $`h`$, the conformal dimension of the field $`\psi `$, to be extracted. Comparing with the Kac formula $`h_{ab}=((bp^{}aq^{})^2(p^{}q^{})^2)/(4p^{}q^{})`$ at the appropriate (IR) values of $`p^{}`$ and $`q^{}`$ leads us to conjecture the following pattern of arriving operators $`\psi `$:
$`_{p,q}+\varphi _{21}`$ $``$ $`_{qp,q}(p<q<2p)\text{arriving via }\varphi _{21},`$ (4.6)
$`_{p,q}+\varphi _{15}`$ $``$ $`_{p,4pq}(2p<q<3p)\text{ arriving via }\varphi _{15},`$ (4.7)
$`_{p,q}+\varphi _{15}`$ $``$ $`_{4pq,p}(3p<q<4p)\text{ arriving via }\varphi _{51}.`$ (4.8)
In making these identifications, we took account of the fact that $`\varphi _{21}`$ is an odd operator, while $`\varphi _{15}`$ is even.
Exactly the same line of argument starting from the $`e^\theta `$ term in (4.2) reveals a further series of corrections to $`c_{\mathrm{eff}}`$ as powers of $`r^2`$, perfectly adapted to match the $`T\overline{T}`$ terms in (4.5). However this time it is possible to say more, by observing that the effect of this part of $`\chi (\theta )`$ on equations (2.7) and (2.8) can be reabsorbed into a shift of $`r`$ by a constant, and furthermore that this constant can be expressed in terms of $`c(r)`$. This allows the first few terms of the iterative expansion to be found exactly. The same idea was used in to extract IR asymptotics from various massless TBA systems; in the current context we find
$$c_{\mathrm{eff}}(r)c_{\mathrm{eff}}(\mathrm{})\frac{\pi c_{\mathrm{eff}}(\mathrm{})^2}{\mathrm{sin}\frac{\pi }{3}(\xi 1)}r^2+\frac{2\pi ^2c_{\mathrm{eff}}(\mathrm{})^3}{\mathrm{sin}^2\frac{\pi }{3}(\xi 1)}r^4+\mathrm{}$$
(4.9)
Matching $`r^2`$ terms in (4.5) and (4.9) gives the exact relation $`\kappa _2=6/(\pi ^2\mathrm{sin}\frac{\pi }{3}(\xi 1))`$ ; the fact that the $`r^4`$ terms then agree provides a nontrivial check on the IR behaviour of the massless NLIE.
We also performed some numerical fits on the IR data. Our results are summarised in table 2. Accuracy was fairly low, but where relevant we have also included the predicted (‘exact’) values of coefficients, obtained from equation (4.9).
The situation in the ultraviolet is in many respects much simpler. Conformal perturbation theory gives direct access to a function $`c_{\mathrm{pert}}(r)`$, related to the ground state energy (1.3) as $`E_0(M,R)=\frac{\pi }{6R}c_{\mathrm{pert}}(r)`$. Note that $`c_{\mathrm{pert}}`$ contains a bulk part which must be subtracted before comparisons are made with the NLIE. For a theory perturbed by a relevant primary operator $`\varphi `$ with scaling dimensions $`(h_{\mathrm{U}V},h_{\mathrm{U}V})`$ , $`c_{\mathrm{pert}}`$ has the expansion
$$c_{\mathrm{pert}}(r)=c(0)+\underset{n=1}{\overset{\mathrm{}}{}}C_n(\lambda R^y)^n$$
(4.10)
and, in contrast to the situation in the IR, the series is expected to have a finite radius of convergence. Here $`y=22h_{\mathrm{U}V}`$, $`\lambda `$ is the coupling, and the coefficients $`C_n`$ are given in terms of the connected correlation functions of the perturbing field on the plane as
$$C_n=\frac{12(1)^n}{n!(2\pi )^{yn1}}\underset{j=2}{\overset{n}{}}\frac{d^2z_j}{|z_j|^y}V(0)\varphi (1,1)\varphi (z_2,\overline{z}_2)\mathrm{}\varphi (z_n,\overline{z}_n)V(\mathrm{})_\mathrm{C}^{},$$
(4.11)
where $`V`$ creates the CFT ground state on the cylinder. This is the state with lowest conformal dimension, and for a general minimal model $`_{p,q}`$ it corresponds to the field $`\varphi _0\varphi _{ab}`$ with $`a`$ and $`b`$ integers satisfying $`bpaq=1`$. (Only for the unitary models $`_{p,p+1}`$ does it coincide with the conformal vacuum $`\varphi _{11}`$.)
For the perturbing operator, we will be interested in $`\varphi =\varphi _{15}`$ and $`\varphi =\varphi _{21}`$. As mentioned above, $`\varphi _{21}`$ is odd, so only the even coefficients $`C_{2n}`$ are nonzero for this case. Since $`\lambda `$ must be related to the crossover scale $`M`$ as
$$\lambda =\kappa M^y$$
(4.12)
with $`\kappa `$ a dimensionless constant, (4.10) is a series in $`r^y`$ for the $`\varphi _{15}`$ perturbations, and $`r^{2y}`$ for the $`\varphi _{21}`$ perturbations.
The effective central charge calculated using the NLIE is expected to expand as
$$c_{\mathrm{eff}}(r)=c_{\mathrm{eff}}(0)+B(r)+\underset{n=1}{\overset{\mathrm{}}{}}c_nr^{yn}$$
(4.13)
for some value of $`y`$. Periodicity arguments suggest that this will be a series in $`r^{6/(1+\xi )}`$, and if $`\xi `$ is chosen according to (2.5) or (2.6) then $`y=22h_{\mathrm{U}V}`$ and the perturbative expansion (4.10) is matched, with $`h_{\mathrm{U}V}`$ the conformal dimension of either $`\varphi _{21}`$ or $`\varphi _{15}`$, and all odd terms zero for the $`\varphi _{21}`$ case . The irregular bulk term $`B(r)`$ must be subtracted before $`c_{\mathrm{pert}}`$ can be compared with $`c_{\mathrm{eff}}`$. Fortunately, it can be obtained exactly from the NLIE, using a small generalisation of arguments used in and . The term we need is given by the behaviour as $`r0`$ of
$$2\frac{3ir}{2\pi ^2}\left[_{\genfrac{}{}{0pt}{}{𝒞_2}{>0}}e^\theta \frac{d}{d\theta }\mathrm{ln}(1+e^{f_L(\theta )})𝑑\theta _{\genfrac{}{}{0pt}{}{𝒞_1}{>0}}e^\theta \frac{d}{d\theta }\mathrm{ln}(1+e^{f_L(\theta )})𝑑\theta \right],$$
(4.14)
where the ‘$`>0`$’ indicates that only those parts of the contours $`𝒞_1`$ and $`𝒞_2`$ with positive real part should be taken, and the symmetry between $`f_L`$ and $`f_R`$ was used to trade the first two integrals in (2.9) for the prefactor $`2`$. Consider the $`r0`$ limit of (2.7) and (2.8) in the region $`0\mathrm{}e\theta \mathrm{ln}(1/r)`$ , where the driving term $`i\frac{r}{2}e^\theta `$ in (2.8) can be dropped. Take the derivative of these equations with respect to $`\theta `$, and extract the contributions to the convolutions proportional to $`e^\theta `$ using (4.1) and (4.2). These should cancel either against the remaining driving term or between themselves to ensure that the functions $`f_L`$ and $`f_R`$ have no such dependency. This leads to the equations
$`i{\displaystyle \frac{r}{2}}`$ $`=`$ $`{\displaystyle \frac{\sqrt{3}\mathrm{sin}(\frac{\pi }{3}\xi )}{\pi \mathrm{sin}(\frac{\pi }{3}(\xi 1))}}\left[{\displaystyle _{\genfrac{}{}{0pt}{}{𝒞_1}{>0}}}e^\theta ^{}{\displaystyle \frac{d}{d\theta ^{}}}\mathrm{ln}(1+e^{f_R(\theta ^{})})𝑑\theta ^{}{\displaystyle _{\genfrac{}{}{0pt}{}{𝒞_2}{>0}}}e^\theta ^{}{\displaystyle \frac{d}{d\theta ^{}}}\mathrm{ln}(1+e^{f_R(\theta ^{})})𝑑\theta ^{}\right]`$
$`{\displaystyle \frac{3}{2\pi \mathrm{sin}(\frac{\pi }{3}(\xi 1))}}\left[{\displaystyle _{\genfrac{}{}{0pt}{}{𝒞_1}{>0}}}e^\theta ^{}{\displaystyle \frac{d}{d\theta ^{}}}\mathrm{ln}(1+e^{f_L(\theta ^{})})𝑑\theta ^{}{\displaystyle _{\genfrac{}{}{0pt}{}{𝒞_2}{>0}}}e^\theta ^{}{\displaystyle \frac{d}{d\theta ^{}}}\mathrm{ln}(1+e^{f_L(\theta ^{})})𝑑\theta ^{}\right];`$
$`0`$ $`=`$ $`{\displaystyle \frac{\sqrt{3}\mathrm{sin}(\frac{\pi }{3}\xi )}{\pi \mathrm{sin}(\frac{\pi }{3}(\xi 1))}}\left[{\displaystyle _{\genfrac{}{}{0pt}{}{𝒞_2}{>0}}}e^\theta ^{}{\displaystyle \frac{d}{d\theta ^{}}}\mathrm{ln}(1+e^{f_L(\theta ^{})})𝑑\theta ^{}{\displaystyle _{\genfrac{}{}{0pt}{}{𝒞_1}{>0}}}e^\theta ^{}{\displaystyle \frac{d}{d\theta ^{}}}\mathrm{ln}(1+e^{f_L(\theta ^{})})𝑑\theta ^{}\right]`$
$`{\displaystyle \frac{3}{2\pi \mathrm{sin}(\frac{\pi }{3}(\xi 1))}}\left[{\displaystyle _{\genfrac{}{}{0pt}{}{𝒞_2}{>0}}}e^\theta ^{}{\displaystyle \frac{d}{d\theta ^{}}}\mathrm{ln}(1+e^{f_R(\theta ^{})})𝑑\theta ^{}{\displaystyle _{\genfrac{}{}{0pt}{}{𝒞_1}{>0}}}e^\theta ^{}{\displaystyle \frac{d}{d\theta ^{}}}\mathrm{ln}(1+e^{f_R(\theta ^{})})𝑑\theta ^{}\right].`$
Solving for the integrals needed for the evaluation of (4.14), we find
$$B_{massless}(r)=\frac{3}{\pi }\frac{\mathrm{sin}(\frac{\pi }{3}\xi )\mathrm{sin}(\frac{\pi }{3}(\xi 1))}{\mathrm{sin}(\pi \xi )}r^2.$$
(4.15)
A similar if slightly simpler analysis of the massive equation (2.1) predicts
$$B_{massive}(r)=\frac{\sqrt{3}\mathrm{sin}(\frac{\pi }{3}\xi )}{2\pi \mathrm{sin}(\frac{\pi }{3}(\xi +1))}r^2.$$
(4.16)
It can be checked that this formula combines the results given separately for $`\varphi _{12}`$, $`\varphi _{21}`$ and $`\varphi _{15}`$ perturbations in .
The formula (4.15) has a pole when $`\xi +1`$ is an integer multiple of 3. Since the overall result remains finite this infinity should cancel against one of the terms in the regular expansion of $`c_{\mathrm{pert}}`$, leaving a logarithmic contribution to the final result (see for example ). At this point the calculation might appear to split into two cases: the term proportional to $`r^2`$ in the perturbative expansion occurs at an order $`n`$ which depends both on the dimension and on the ‘parity’ of the perturbing operator: $`n=2/y`$ for $`\varphi _{15}`$, and $`1/y`$ for $`\varphi _{21}`$ . However $`y=6/(\xi +1)`$ for $`\varphi _{15}`$ and $`3/(\xi +1)`$ for $`\varphi _{21}`$ , so $`n=(\xi +1)/3`$ in both cases and they can be treated simultaneously. The logarithm is found by evaluating
$$\underset{\xi 3n1}{lim}\frac{3\mathrm{sin}(\frac{\pi }{3}\xi )\mathrm{sin}(\frac{\pi }{3}(\xi 1))}{\pi \mathrm{sin}(\pi \xi )}\left(r^2r^{2n\frac{3}{\xi +1}}\right)$$
(4.17)
which yields
$$B_{massless}(r)|_{\xi =3n1}=(1)^n\frac{3}{2\pi ^2n}r^2\mathrm{ln}r.(n)$$
(4.18)
Similarly, and for exactly the same values of $`\xi `$, logarithms arise in the massive perturbations. In these cases we find
$$B_{massive}(r)|_{\xi =3n1}=\frac{3}{2\pi ^2n}r^2\mathrm{ln}r.(n)$$
(4.19)
This seems to leave the coefficient of $`r^2`$ term unknown, but one piece of information can be extracted, which will be used later. Recall that the logarithm is linked to a divergence in the perturbative integral defining the term $`C_n`$. Providing the same regularisation scheme is used for both flow directions, the perturbative contributions should cancel upon taking their sum or difference. The remaining finite contribution can be found directly from (4.15) and (4.16):
$$(B_{massless}(r)\pm B_{massive}(r))|_{\xi =3n1}=\pm \frac{\sqrt{3}}{4\pi }r^2(n)$$
(4.20)
where the plus signs should be chosen if $`n`$ is odd and the minus signs if $`n`$ is even.
Using an iterative method to solve (2.7) and (2.8), $`c_{\mathrm{eff}}(r)`$ can be calculated to high accuracy. After subtracting the ultraviolet central charge and the relevant bulk term, the perturbative coefficients $`c_n`$ can be estimated via a polynomial fit. To make the comparison with CPT, we must also fix the value of $`\kappa `$. For the massive systems this was determined exactly in , and it turns out, numerically at least, that the same relationship between $`\lambda `$ and $`M`$ holds in the massless cases, apart from a factor of either $`1`$ or $`i`$. As explained in , the relationship between the massive and massless perturbations depends on the parity of the perturbing operator. For $`\varphi _{15}`$ , the massive behaviour is related to the massless by flipping the sign of the coupling constant $`\lambda `$. This has no effect for the odd operator $`\varphi _{21}`$; instead the required transformation is $`\lambda i\lambda `$. The CPT coefficients $`C_n`$ are the same for both flow directions so, provided the mass and the crossover scales $`M`$ are equal, we expect
$`\varphi _{15}:c_n`$ $`=`$ $`(1)^n\stackrel{~}{c}_n,`$ (4.21)
$`\varphi _{21}:c_{2n}`$ $`=`$ $`(1)^n\stackrel{~}{c}_{2n},`$ (4.22)
where $`\stackrel{~}{c}_n`$ denotes the expansion coefficients obtained using the massive NLIE. Tables 3 and 4 present the first few perturbative coefficients for the models $`_{4,11}`$ and $`_{5,8}`$ perturbed in both massless and massive directions. The relative signs of the coefficients confirm relations (4.21) and (4.22). Similar results were obtained for the models $`_{3,8}`$ , $`_{5,12}`$ and $`_{3,10}`$ perturbed by $`\varphi _{15}`$ , and the models $`_{3,4}`$ , $`_{3,5}`$ and $`_{7,11}`$ perturbed by $`\varphi _{21}`$. These provide strong support for our claim that, modulo the bulk terms, the results from the massive and massless integral equations are related by analytic continuation.
For the $`\varphi _{15}`$ perturbations, it is also simple to check the value of $`\kappa `$ directly, since the first term in the perturbative expansion (4.10) is just given by
$$C_1=12(2\pi )^{(1y)}C_{\varphi _0\varphi _{15}\varphi _0}$$
(4.23)
where $`C_{\varphi _0\varphi _{15}\varphi _0}`$ , the operator product coefficient between the perturbing operator $`\varphi _{15}`$ and the ground state $`\varphi _0`$ , can be found in . Comparing (4.10) and (4.13) at order $`n=1`$, we have
$$\kappa =\frac{c_1}{C_1}.$$
(4.24)
Using this formula and the massless NLIE, we estimated $`\kappa ^2`$ numerically for a number of models. Table 5 reports the results, and compares them with the exact expression for $`\kappa ^2`$ for (massive) $`\varphi _{15}`$ perturbations given in :
$$\lambda ^2=\frac{4^2(\xi 1)^2\gamma (\frac{2+\xi }{2(1+\xi )})\gamma (\frac{5\xi }{2(1+\xi )})}{\pi ^2(23\xi )^2(2\xi )^2\gamma ^2(\frac{3+\xi }{1+\xi })}\left[\frac{M\mathrm{\Gamma }(\frac{\xi +1}{3})}{\sqrt{3}\mathrm{\Gamma }(\frac{1}{3})\mathrm{\Gamma }(\frac{\xi }{3})}\right]^{\frac{12}{\xi +1}}$$
(4.25)
where $`\gamma (x)=\frac{\mathrm{\Gamma }(x)}{\mathrm{\Gamma }(1x)}`$ and $`M`$ is the mass of the lightest kink.note the relation was given in in terms of $`\xi ^{\mathrm{FLZZ}}=\frac{p}{qp}`$ and $`m`$, the mass of the lightest breather, related to $`M`$ as $`m=2M\mathrm{sin}(\frac{\pi }{3}\xi )`$ The agreement is clearly very good. Note that $`\kappa ^2`$ is sometimes negative – these are cases where, in the normalisations of , $`C_{\varphi _0\varphi _{15}\varphi _0}`$ turns out to be purely imaginary.
Now we return to the behaviour of the flows (1.1) and (1.2), to mention one reason why, exceptionally, $`c_{\mathrm{eff}}(r)`$ should be a monotonic function of $`r`$ for these flows. The behaviour of a flow is determined at small $`r`$ by the first nonzero term in the perturbative expansion. Thus whether $`c_{\mathrm{eff}}`$ initially increases or decreases will typically be determined by the sign of $`C_1\lambda `$ ($`\varphi _{15}`$) or $`C_2\lambda ^2`$ ($`\varphi _{21}`$). In particular, if $`C_1`$ (respectively $`C_2`$) is nonzero then for one sign of $`\lambda `$ (or $`\lambda ^2`$), $`c_{\mathrm{eff}}`$ will initially increase, leading to an immediate violation of the ‘$`c_{\mathrm{eff}}`$-theorem’, and a flow which must be non-monotonic if $`c_{\mathrm{eff}}(\mathrm{})`$ is to be less than $`c_{\mathrm{eff}}(0)`$. But for the model $`_{3,5}+\varphi _{21}`$, the first coefficient $`C_2`$ was calculated to be zero in . In this case the asymptotic UV behaviour is instead controlled by $`C_4\lambda ^4`$ and this permits $`c_{\mathrm{eff}}`$ to decrease initially for both (massless and massive) signs of $`\lambda ^2`$. A similar calculation for the other models in the $`|I|=1`$ series finds $`C_1`$ ($`\varphi _{15}`$) or $`C_2`$ ($`\varphi _{21}`$) to be zero and thus, as in the first model, all of the flows are able to be monotonic. A numerical fit of the data for models higher up in the series confirms these coefficients are zero within our numerical accuracy. For all other sequences of models, such vanishings of $`C_1`$ or $`C_2`$ do not occur, forcing at least one of each pair of massive and massless flows to be non-monotonic. One remaining mystery, out of many, is why it should always be the massless flow which is non-monotonic.
The $`|I|=1`$ sequence of flows was found in via an associated staircase model, and the final step of this staircase interpolates in the infrared to a massive model with $`c_{\mathrm{eff}}=0`$. It turns out that the NLIE also reproduces this behaviour: even though the operator $`\varphi _{15}`$ is not a member of the Kac table for $`_{2,5}`$ , (3.7) formally predicts a further flow, to a theory with $`c_{\mathrm{eff}}=0`$: $`_{2,5}+\varphi _{15}_{2,3}`$. There is nothing to stop us using the massless NLIE (2.7,2.8) and the $`\varphi _{15}`$ recipe (2.6) to compute an effective central charge for this putative flow. Our numerical results show $`c_{\mathrm{eff}}`$ to have an exponential behaviour in the infrared, as would be expected for a massive, rather than a massless, flow. Furthermore, comparing $`c_{\mathrm{eff}}`$ (massless) with that of the massive flow $`_{2,5}+\varphi _{15}`$ calculated using (2.1) and (2.6), we find the effective central charges match exactly. They also turn out to coincide with the results from the more standard massive flow $`_{25}+\varphi _{12}`$, computed using (2.1) and (2.4). This means that the flow is at least physically reasonable, since $`\varphi _{12}`$ is the single relevant primary field in $`_{2,5}`$. However, it remains a curiosity that the same flow can be found from three different nonlinear integral equations. A sample of our numerical results is shown in table 6. Note that this example shows that the straightforward prediction of scaling dimensions based on periodicity arguments is not always correct: for $`_{2,5}+\varphi _{15}`$ one would have expected a series in $`r^{6/5}`$ for $`c_{\mathrm{eff}}(r)`$, whereas in fact the expansion is in powers of $`r^{12/5}`$.
If we compare the infrared destinations of (3.7) and (3.8) we see it is possible for two different ultraviolet models, both perturbed by $`\varphi _{15}`$, to flow to the same infrared fixed point, one attracted via the irrelevant operator $`\varphi _{15}`$ and the other via $`\varphi _{51}`$. Figure 6 illustrates two such flows. As was noted in §3, the sequence attracted by $`\varphi _{51}`$ necessarily stops at this model, but the other sequence may continue to flow down further.
For some models, the predictions (4.6)–(4.8) must be treated with caution, as we were unable to check them explicitly. Examples are the minimal models $`_{p,q}`$ perturbed by $`\varphi _{15}`$ when $`p`$ and $`q`$ are related as $`4pq=1`$. For these cases the twist $`\alpha ^{}`$ is equal to $`2`$, and a trivial shift of $`f_{L/R}(\theta )`$ leaves a system with $`\alpha ^{}=0`$. This prevents us from tuning in a straightforward way onto the required ultraviolet model: numerically, the massless NLIE simply recovers a flow from $`c_{\mathrm{eff}}(0)=1`$ to $`c_{\mathrm{eff}}(\mathrm{})=1`$ .
We also observe that (3.6) formally predicts the following flows from unitary minimal models:
$$_{p,p+1}+\varphi _{21}_{1,p},(\xi ,\alpha ^{})=(\frac{p}{p+1},1).$$
(4.26)
However it is a little hard to decide what is meant by $`_{1,p}`$ , and indeed, while our numerical results indicate that these flows behave as expected for small $`r`$, at some intermediate scale $`c_{\mathrm{eff}}(r)`$ appears to have a discontinuity. This could be linked to $`\alpha ^{}`$ being equal to $`1`$. In any event, we suspect there may be a square root singularity though a detailed investigation will have to be left for future work. Note that by the type II equivalence, a similar phenomenon occurs for the models $`_{p,4p2}`$ perturbed by $`\varphi _{15}`$. In one case a detailed check on these exceptional flows can be made with relative ease: the Ising model $`_{3,4}`$ perturbed by $`\varphi _{21}`$. For real coupling $`\lambda `$, an exact expression for the (massive) effective central charge was given in . Using this we can make an exact prediction for $`c_{\mathrm{eff}}`$ in the massless case: we send $`rir`$ in the perturbative expansion, swap the logarithmic bulk term for that of the massless model (4.15) and find the coefficient of the $`r^2`$ term using (4.20). The result is:
$`c_{\mathrm{eff}}(r)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{3r^2}{2\pi ^2}}\left[\mathrm{ln}r{\displaystyle \frac{1}{2}}\mathrm{ln}\pi +\gamma _E+{\displaystyle \frac{\pi }{2\sqrt{3}}}\right]`$ (4.27)
$`+{\displaystyle \frac{6}{\pi }}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\left(\sqrt{(2k1)^2\pi ^2r^2}(2k1)\pi +{\displaystyle \frac{r^2}{2(2k1)\pi }}\right)`$
where $`\gamma _E=0.57721556\mathrm{}`$ is the Euler-Mascheroni constant. This formula predicts a square root singularity at $`r=\pi `$, preventing a smooth interpolating flow to the far infrared, but we can at least match against the NLIE for values of $`r`$ out to this point. As table 7 shows, our numerical results agree well with the exact formula (4.27).
## 5 A remark on $`\varphi _{13}`$ perturbations
In this short section we comment on some features of massless $`\varphi _{13}`$ perturbations which are revealed when similar methods are used. In numerical results were only presented for the zero twist ($`c_{\mathrm{eff}}(0)=c_{\mathrm{eff}}(\mathrm{})=1`$) case, but the same equation can describe a minimal model $`_{p,q}`$ perturbed by $`\varphi _{13}`$. One simply has to set the kernel parameter $`\zeta `$ ($`p`$ in ) to $`p/(qp)`$, and the twist $`\alpha `$ in that paper to $`\pi /p`$.
Defining an index $`J=qp`$ , $`\zeta =p/J`$ and the flow from $`\zeta `$ to $`\zeta 1`$ is
$$_{p,p+J}+\varphi _{13}_{pJ,p}\text{attracted via }\varphi _{31}\text{ .}$$
(5.1)
Again, the number of sequences for each $`J`$ is given by the Euler $`\phi `$-function, $`\phi (J)`$ . As in the $`a_2^{(2)}`$ case, we cannot access all possible minimal models via this equation. The kernel $`\varphi (\theta )`$ in the massless equation of has a pole at $`i\pi (\zeta 1)`$ which crosses the real $`\theta `$ axis as $`\zeta `$ falls below $`1`$. As before, this could probably be overcome using analytic continuation but for now we choose $`\zeta >1`$, requiring $`2p>q`$, to prevent such problems occuring. We also find a set of models, $`_{p,2p1}+\varphi _{13}`$, where $`\alpha ^{}=1`$ and the would-be IR model is $`_{1,p}`$. As in the $`a_2^{(2)}`$ case, $`c_{\mathrm{eff}}`$ suffers some sort of discontinuity at an intermediate scale for these flows.
Solving the massless equation in for the unitary cases ($`J=1`$), we find that results from the TBA equations of are matched to high accuracy<sup>\**</sup><sup>\**</sup>\**we understand that this has already been checked by Alyosha Zamolodchikov . The results for $`J>1`$ are perhaps more interesting, since TBA systems for these nonunitary flows are not known. Just as for the $`\varphi _{21}/\varphi _{15}`$ flows with $`|I|>1`$, the monotonicity property is lost. Typical cases are illustrated in figures 7 and 8.
## 6 Conclusions
We hope to have demonstrated the rich structure of flows that can be studied by means of a relatively simple nonlinear integral equation. A number of unexpected features have emerged, most notably the consistent failure of the effective central charges to be monotonic functions of the system size. Indeed, looking at the full set of minimal models, both unitary and nonunitary, we see that for massless flows a monotonic behaviour of $`c_{\mathrm{eff}}(r)`$ is very much the exception rather than the rule.
As for future work, the following directions seem natural:
* A more detailed study of both the massive and the massless $`a_2^{(2)}`$-related nonlinear integral equations is warranted. This should include a more detailed look at their analytic properties, and further comparisons with perturbed conformal field theories. In addition, the equations should be generalised to describe excited states.
* The type II conjecture remains a curious observed fact, and a deep understanding seems still to be lacking. In cases suffering from this ambiguity, a knowledge of $`c_{\mathrm{eff}}(r)`$ will never suffice to disentangle the possible destinations of the flows. It would be reassuring to confirm that these models behave as conjectured by some other method.
* The work of was based in part on a massless S-matrix. Here we avoided this aspect, but it would be interesting to find an S-matrix description of the new flows. One might also hope to find integrable lattice models which would yield the massless flows in their continuum limits.
* As discussed in §3, we were unable to treat some flows, due to singularities in the kernel $`\varphi (\theta )`$ crossing the integration contour. It would be worthwhile to study the NLIE in these zones where $`\xi `$ falls below $`1`$, perhaps by analytic continuation. It would also be interesting to study the NLIE of for $`\zeta <1`$, where a similar difficulty is encountered.
* The method described in §2 should enable a massless nonlinear integral equation to be obtained from any of the (massive) ADE-related systems described in . It would be interesting to know whether these share the same strange properties that we have observed here in the examples related to $`a_2^{(2)}`$ and $`a_1^{(1)}`$.
* As mentioned in §4, the initial discovery of the flows (1.1) and (1.2) in came via a staircase model. The new flows also fit naturally into staircase-like patterns, and it is natural to ask whether the existing set of known staircase models could be enlarged so as to include these cases.
* Finally, it remains an open question to find TBA equations describing the non-unitary $`\varphi _{13}`$ flows, or any of the new $`\varphi _{21}`$ or $`\varphi _{15}`$ flows. Why is there no TBA for any non-monotonic case? And what has this to do with the monstron<sup>††</sup><sup>††</sup>††see , if anything?
Acknowledgements – We are grateful to Changrim Ahn, Philippe Di Francesco, Francesco Ravanini, Tom Wynter, Alyosha Zamolodchikov and Jean-Bernard Zuber for useful discussions. PED and TCD thank the EPSRC for an Advanced Fellowship and a Research Studentship respectively, and RT thanks the Universiteit van Amsterdam for a post-doctoral fellowship. PED and TCD are grateful to SPhT Saclay, and RT to Durham University and the APCTP, for hospitality during various stages of this project. The work was supported in part by a TMR grant of the European Commission, reference ERBFMRXCT960012. |
warning/0001/nucl-th0001024.html | ar5iv | text | # Excitation functions of hadronic observables from SIS to RHIC energiessupported by BMBF and GSI Darmstadt
## 1 Introduction
The ultimate goal of relativistic nucleus-nucleus collisions is to reanalyze the early ’big-bang’ under laboratory conditions and to find the ’smoking gun’ for a phase transition from the initial quark-gluon plasma (QGP) to a phase characterized by an interacting hadron gas . Any theoretical approach might describe such a phase transition starting from the partonic side with strongly interacting quarks and gluons or from the hadronic side by involving hadronic degrees of freedom , i.e. hadrons with proper self-energies or spectral functions at high baryon density or temperature. It remains to be seen which approach will prove to be more successful, economic and transparent.
Nucleus-nucleus collisions with initial energies per nucleon of $`\sqrt{s}`$ = 200 GeV or $``$21.5 A$``$TeV will be available soon at the Relativistic-Heavy-Ion-Collider (RHIC) in Brookhaven. In central collisions of Au + Au here energy densities above 5 GeV/fm<sup>3</sup> are expected such that the critical energy density for a QGP phase should be overcome in considerable space-time volumes where the relevant degrees of freedom are partons (quarks and gluons). Parton cascade calculations have been used so far to estimate the energy densities and particle production yields in violent reactions at $`\sqrt{s}`$ = 200 GeV, an order of magnitude higher than at SPS energies ($`\sqrt{s}`$ 20 GeV). Intuitively one expects that the initial nonequilibrium phase of a nucleus-nucleus collision at RHIC energies should be described by parton degrees of freedom whereas hadrons are only formed (by ’condensation’) at a later stage of the reaction which might be a couple of fm/c from the initial contact of the heavy ions. Thus parton cascade calculations – including transitions rates from perturbative QCD – should be adequate for all initial reactions involving a large 4-momentum transfer between the consituents since QCD is well tested in its short distance properties. The question, however, remains to which extent the parton calculations can be extrapolated to low $`Q^2`$ where hadronic scales become important. As a rough estimate one can employ here the average mass of vector mesons, the nucleon and its first excited state, which gives $`Q_{crit}^2`$ 1 GeV<sup>2</sup>. On the other hand, using the uncertainty relation this implies time scales of $`\mathrm{\Delta }t<`$ 0.2 fm/c = $`Q_{crit}^1`$ or relative separations of partons $`\mathrm{\Delta }r<`$ 0.2 fm, which are small compared to the hadronic size or average life time of the $`\rho ,\mathrm{\Delta },etc.`$ in free space or the formation time of hadrons $`t_F0.70.8`$ fm/c as used in the HSD transport approach .
Turning the argument around, a nonequilibrium hadronic approach involving a time scale of 0.7-0.8 fm/c cannot tell anything about shorter times because the uncertainty relation does not allow to distinguish states which are separated in mass by less than $``$ 300 MeV = $`t_F^1`$, which is the $`N\mathrm{\Delta }`$ mass difference. Thus one faces the problem that neither the parton description nor a nonequilibrium hadronic model should be valid for times 0.2 fm/c $`t`$ 0.7-0.8 fm/c in individual hadronic reactions, which corresponds to the nonperturbative formation time of the hadronic wavefunction. This regime of the ’soft’ QCD physics is presently not well understood and appropriate dynamical models are urgently needed. In the HSD approach these intermediate times are described by color neutral strings, where the leading quarks and ’diquarks’ in a baryonic string (or quarks and antiquarks in a mesonic string etc.) are allowed to rescatter again with hadronic cross sections divided by the number of constituent quarks and antiquarks in the hadrons, respectively .
Furthermore, the question of chiral symmetry restoration at high baryon density and/or high temperature is of fundamental interest, too . Whereas lattice QCD calculations at zero baryon chemical potential indicate that a restoration of chiral symmetry goes along with the deconfinement phase transition at some critical temperature $`T_c`$, the situation is less clear for finite baryon density where QCD sum rule studies indicate a linear decrease of the scalar quark condensate $`<\overline{q}q>`$ – which is nonvanishing in the vacuum due to a spontaneous breaking of chiral symmetry – with baryon density $`\rho _B`$ towards a chiral symmetric phase characterized by $`<\overline{q}q>`$ = 0. This decrease of the scalar condensate is expected to lead to a change of the hadron properties with density and temperature, i.e. in a chirally restored phase the hadrons might become approximately massless or at least vector and axial vector currents should become equal ; the latter implies that the $`\rho `$ and $`a_1`$ spectral functions should become identical. Since the scalar condensate $`<q\overline{q}>`$ is not a direct observable, its manifestations should be found in different hadronic abundancies and spectra.
Nowadays, our knowledge about the hadron properties at high temperature or baryon density is based on heavy-ion experiments from BEVALAC/SIS to SPS energies where hot and dense nuclear systems are produced on a timescale of a few fm/c. As mentioned above, the information from ultrarelativistic nucleus-nucleus collisions at RHIC, i.e. initial $`\sqrt{s}=`$ 200 GeV per nucleon, will be available soon . However, any conclusions about the properties of hadrons in the nuclear environment are based on the comparison of experimental data with nonequilibrium kinetic transport theory . As a genuine feature of transport theories there are two essential ingredients: i.e. the baryon (and meson) scalar and vector self-energies as well as in-medium elastic and inelastic cross sections for all hadrons involved. Whereas in the low-energy regime these ’transport coefficients’ can be calculated in the Dirac-Brueckner approach starting from the bare nucleon-nucleon interaction this is no longer possible at high baryon density $`(\rho _B2`$-$`3\rho _0)`$ and high temperature, since the number of independent hadronic degrees of freedom increases drastically and the interacting hadronic system should enter a phase with $`<q\overline{q}>`$ 0 as mentioned before. As a consequence the hadron self-energies or spectral functions in the nuclear medium will change substantially especially close to the chiral phase transition and transport theoretical studies should include the generic properties of QCD in line with nonperturbative computations on the lattice .
In this work we concentrate on excitation functions of hadronic observables from SIS to RHIC energies with the aim to find out optimal experimental conditions to search for ’traditional’ phenomena such as strangeness enhancement (Section 3), low mass dilepton enhancement (Section 4) or charmonium suppression in nucleus-nucleus collisions (Section 5). Our studies are performed within the HSD transport approach that has been described in Refs. and been tested for $`p+p`$, $`p+A`$ and $`A+A`$ collisions from the SIS to SPS energy regime . Actual predictions for hadron rapidity spectra and $`J/\mathrm{\Psi }`$ suppression will be presented in Section 6 while Section 7 concludes the study with a summary.
## 2 Theoretical considerations
In this Section we briefly recall the ingredients of the covariant transport theory, that has been denoted as Hadron-String-Dynamics (HSD) , which formally can be written as a coupled set of transport equations for the phase-space distributions $`f_h(x,p)`$ of hadron $`h`$ , i.e.
$`\left\{\left(\mathrm{\Pi }_\mu \mathrm{\Pi }_\nu _\mu ^pU_h^\nu M_h^{}_\mu ^pU_h^S\right)_x^\mu +\left(\mathrm{\Pi }_\nu _\mu ^xU_h^\nu +M_h^{}_\mu ^xU_h^S\right)_p^\mu \right\}f_h(x,p)`$
$`={\displaystyle \underset{h_2h_3h_4\mathrm{}}{}}{\displaystyle d2d3d4\mathrm{}[G^{}G]_{1234\mathrm{}}\delta _\mathrm{\Gamma }^4(\mathrm{\Pi }+\mathrm{\Pi }_2\mathrm{\Pi }_3\mathrm{\Pi }_4\mathrm{})}`$
$`\times \{f_{h_3}(x,p_3)f_{h_4}(x,p_4)\overline{f}_h(x,p)\overline{f}_{h_2}(x,p_2)`$
$`f_h(x,p)f_{h_2}(x,p_2)\overline{f}_{h_3}(x,p_3)\overline{f}_{h_4}(x,p_4)\}\mathrm{}.`$ (1)
In Eq. (1) $`U_h^S(x,p)`$ and $`U_h^\mu (x,p)`$ denote the real part of the scalar and vector hadron self-energies, respectively, while $`[G^+G]_{1234\mathrm{}}\delta _\mathrm{\Gamma }^4(\mathrm{\Pi }+\mathrm{\Pi }_2\mathrm{\Pi }_3\mathrm{\Pi }_4\mathrm{})`$ is the ’transition rate’ for the process $`1+23+4+\mathrm{}`$ which is taken to be on-shell in the semiclassical limit adopted. The hadron quasi-particle properties in (1) are defined via the mass-shell constraint ,
$$\delta (\mathrm{\Pi }_\mu \mathrm{\Pi }^\mu M_h^2),$$
(2)
with effective masses and momenta (for a hadron of bare mass $`M_h`$ and momentum $`p^\mu `$) given by
$`M_h^{}(x,p)`$ $`=`$ $`M_h+U_h^S(x,p)`$
$`\mathrm{\Pi }^\mu (x,p)`$ $`=`$ $`p^\mu U_h^\mu (x,p),`$ (3)
while the phase-space factors
$$\overline{f}_h(x,p)=1\pm f_h(x,p)$$
(4)
are responsible for fermion Pauli-blocking or Bose enhancement, respectively, depending on the type of hadron in the final/initial channel. The dots in Eq. (1) stand for further contributions to the collision term with more than two hadrons in the final/initial channels. The transport approach (1) is fully specified by $`U_h^S(x,p)`$ and $`U_h^\mu (x,p)`$ $`(\mu =0,1,2,3)`$, which determine the mean-field propagation of the hadrons, and by the transition rates $`G^{}G\delta ^4(\mathrm{})`$ in the collision term, that describe the scattering and hadron production/absorption rates.
In Ref. the scalar and vector mean-fields $`U_h^S`$ and $`U_h^\mu `$ have been determined in the mean-field limit from an effective hadronic Lagrangian density $`_H`$ that has been fitted to the equation of state of nucleonic matter as resulting from the Nambu-Jona-Lasinio (NJL) model. Without going through the arguments again we show in Fig. 1 the nucleon scalar ($`U_S`$) and negative vector potential ($`U_0`$) as a function of the nuclear density $`\rho `$ and relative momentum $`𝐩`$ of the nucleon with respect to the nuclear matter rest frame. Whereas the vector potential increases practically linearly with density (at low momentum $`𝐩`$) the scalar potential saturates with density such that the nucleon effective mass $`M^{}=M_0+U_S`$ almost drops to zero for $`\rho `$ 0.6 fm<sup>-3</sup>. Both potentials decrease rather fast in magnitude with momentum $`𝐩`$ and practically vanish above a few GeV/c.
In Fig. 2 the real part of the potential
$`U_{SEP}`$ $`=`$ $`U_0(\rho _0,𝐩)+\sqrt{𝐩^2+(M_N+U_S)^2}\sqrt{𝐩^2+M_N^2}`$ (5)
is shown again as a function of $`\rho `$ and $`𝐩`$. Whereas we see a net attraction for momenta $`|𝐩|`$ 0.5 GeV/c up to densities of $``$ 0.3 fm<sup>-3</sup>, the net potential becomes repulsive for higher momenta, reaches a maximum repulsion at $`|𝐩|`$ 1 GeV/c and then drops again with $`|𝐩|`$. We mention that at density $`\rho _0`$ the potential $`U_{SEP}`$ compares well with the potential from the data analysis of Hama et al. as well as Dirac-Brueckner computations from up to a kinetic energy $`E_{kin}`$ of 1 GeV . The formula (5) reduces to the familiar expression for the Schroedinger equivalent potential (Eq. (3.16) of Ref. ) in the low density limit.
In our transport calculations we include nucleons, $`\mathrm{\Delta }`$’s, N(1440), N(1535), $`\mathrm{\Lambda }`$, $`\mathrm{\Sigma }`$ and $`\mathrm{\Sigma }^{}`$ hyperons, $`\mathrm{\Xi }`$’s and $`\mathrm{\Omega }`$’s as well as their antiparticles. In a first approximation we assume here that all baryons (made out of light $`(u,d)`$ quarks) have the same scalar and vector self-energies as the nucleons while the hyperons pick up a factor 2/3 according to the light quark content and $`\mathrm{\Xi }`$’s a factor of 1/3, respectively.
In the HSD approach the high energy inelastic hadron-hadron collisions are described by the FRITIOF model , where two incoming hadrons emerge the reaction as two excited color singlet states, i.e. ’strings’. According to the Lund model a string is characterized by the leading constituent quarks of the incoming hadron and a tube of color flux is supposed to be formed connecting the rapidly receding string-ends. In the HSD approach baryonic ($`qqq`$) and mesonic ($`q\overline{q}`$) strings are considered with different flavors ($`q=u,d,s`$). In the uniform color field of the strings virtual $`q\overline{q}`$ or $`qq\overline{q}\overline{q}`$ pairs are produced causing the tube to fission and thus to create mesons or baryon-antibaryon pairs. The production probability $`P`$ of massive $`s\overline{s}`$ or $`qq\overline{q}\overline{q}`$ pairs is suppressed in comparison to light flavor production ($`u\overline{u}`$, $`d\overline{d}`$) according to a Schwinger-like formula
$`{\displaystyle \frac{P(s\overline{s})}{P(u\overline{u})}}=\gamma _s=\mathrm{exp}(\pi {\displaystyle \frac{m_s^2m_q^2}{2\kappa }},)`$ (6)
with $`\kappa 1`$ GeV/fm denoting the string tension. Thus in the Lund string picture the production of strangeness and baryon-antibaryon pairs is controlled by the constituent quark and diquark masses. Inserting the constituent quark masses $`m_u=0.3`$ GeV and $`m_s=0.5`$ GeV a value of $`\gamma _s0.3`$ is obtained. While the strangeness production in proton-proton collisions at SPS energies is reasonably well reproduced with this value, the strangeness yield for p + Be collisions at AGS energies is underestimated by roughly 30% (cf. ). For that reason the relative factors used in the HSD model are
$`u:d:s:uu=\{\begin{array}{cc}1:1:0.3:0.07\hfill & ,\text{at SPS to RHIC energies}\hfill \\ 1:1:0.4:0.07\hfill & ,\text{at AGS energies},\hfill \end{array}`$ (9)
with a linear transition of the strangeness suppression factor $`\gamma _s=s:u`$ as a function of $`\sqrt{s}`$ in between.
Additionally a fragmentation function $`f(x,m_t)`$ has to be specified, which is the probability distribution for hadrons with transverse mass $`m_t`$ to acquire the energy-momentum fraction $`x`$ from the fragmenting string,
$`f(x,m_t){\displaystyle \frac{1}{x}}(1x)^a\mathrm{exp}\left(bm_t^2/x\right),`$ (10)
with $`a=0.23`$ and $`b=0.34`$ GeV<sup>-2</sup> .
We recall that the LUND model includes partonic diffractive scattering and mini-jet production as well . The latter phenomena are not important at SPS energies and below, however, become appreciable at RHIC energies. In this respect the HSD approach dynamically also includes the hard partonic processes as far as quarks and antiquarks are involved. However, it does not employ hard gluon-gluon processes beyond the level of ’string phenomenology’. This has to be kept in mind with respect to the predictive power of the model at RHIC energies and beyond.
The medium modifications due to the hadron self-energies, furthermore, require to introduce some conserving approximations in the collision terms in line with the modified quasi-particle properties. Since these in-medium modifications – related to ’low momentum physics’ – are not of primary interest in this study we discard an explicit discussion here and refer the reader to Ref. .
### 2.1 The scalar condensate
The scalar quark condensate $`<q\overline{q}>`$ is viewed as an order parameter for the restoration of chiral symmetry at high baryon density and temperature. A model independent relation for the scalar quark condensate at finite (but small) baryon density and temperature has been given by Drukarev and Levin , i.e
$$\frac{<q\overline{q}>}{<q\overline{q}>_V}=1\frac{\rho _S}{f_\pi ^2m_\pi ^2}\left[\mathrm{\Sigma }_\pi +m\frac{d}{dm}\left(\frac{E(\rho )}{A}\right)\right],$$
(11)
where $`<q\overline{q}>_V`$ denotes the vacuum condensate, $`\mathrm{\Sigma }_\pi `$ 45 MeV is the pion-nucleon $`\mathrm{\Sigma }`$-term, $`f_\pi `$ and $`m_\pi `$ the pion decay constant and pion mass, respectively, while $`E(\rho )/A`$ is the binding energy per nucleon. Since for low densities the scalar density $`\rho _S`$ in (11) may be replaced by the baryon density $`\rho _B`$, the change in the scalar quark condensate starts linearly with $`\rho _B`$ and is reduced by a factor 1/3 at saturation density $`\rho _0`$. A simple linear extrapolation then would indicate that at $`\rho _B3\rho _0`$ a restoration of chiral symmetry might be achieved in heavy-ion collisions.
A reasonable estimate for the scalar condensate in dynamical calculations has been suggested by Friman et al. ,
$$\frac{<q\overline{q}>}{<q\overline{q}>_V}=1\frac{\mathrm{\Sigma }_\pi }{f_\pi ^2m_\pi ^2}\rho _S\underset{h}{}\frac{\sigma _h\rho _S^h}{f_\pi ^2m_\pi ^2},$$
(12)
where $`\sigma _h`$ denotes the $`\sigma `$-commutator of the relevant mesons $`h`$. For pions and mesons made out of light quarks and antiquarks we use $`\sigma _h=m_\pi /2`$ whereas for mesons with a strange (antistrange) quark we adopt $`\sigma _h=m_\pi /4`$ according to the light quark content. Within the same spirit the $`\mathrm{\Sigma }`$-term for hyperons is taken as 2/3 $`\mathrm{\Sigma }_\pi `$ 30 MeV while for $`\mathrm{\Xi }`$’s we use 1/3 $`\mathrm{\Sigma }_\pi `$ 15 MeV.
The scalar density of mesons (of type $`h`$) is given by
$$\rho _S^h=\frac{(2s+1)(2t+1)}{(2\pi )^3}d^3𝐩\frac{m_h}{\sqrt{𝐩^2+m_h^2}}f_h(𝐫,𝐩;t),$$
(13)
with $`f_h`$ denoting the meson phase-space distribution of species $`h`$. In (13) $`s,t`$ denote the spin and isospin quantum numbers, respectively. We note that the scalar density of baryons is calculated in line with (13) by replacing the mass and momentum by the effective quantities in (3).
The actual numerical result for the space-time dependence of the scalar condensate (12) is shown in Figs. 3 and 4 for central Au + Au collisions at 6 A$``$GeV and 20 A$``$GeV, respectively. Here the condensate is divided by the vacuum condensate $`<q\overline{q}>_V`$ such that the nonperturbative vacuum is characterized by a value of 1. Furthermore, the $`z`$-direction has been stretched by the Lorentz-factors $`\gamma _{cm}`$ to compensate for Lorentz contraction, while negative numerical values for the condensate have been suppressed. According to (11), (12) the scalar condensate is reduced inside the approaching nuclei by about 35%; this reduction becomes more pronounced when the nuclei achieve full overlap. As seen from Fig. 3 already at 6 A$``$GeV there is a substantial space-time region of vanishing scalar condensate for 6 fm/c $`t13`$ fm/c, where the conventional hadronic picture is not expected to hold anymore. The space-time volume of vanishing quark condensate slightly increases for 20 A$``$GeV (Fig. 4), however, the increase is only moderate and resembles very much the situation for the space-time integral of high density baryon matter (cf. Fig. 1 in Ref. ).
We mention that at higher bombarding energies the 4-volume ($`x=(t,𝐫)`$)
$`V(\alpha )={\displaystyle d^3𝐫𝑑t\mathrm{\Theta }(\alpha \frac{<q\overline{q}>_x}{<q\overline{q}>_V})},`$ (14)
that counts the fraction of the scalar condensate below the value of $`\alpha `$ (0 $`\alpha `$ 0.3) is essentially determined by the pion density whereas below about 20 A$``$GeV the dominant contribution stems from the baryons (cf. also Ref. ). Since the pion density can be considered as a measure of the vacuum ’temperature’, a chiral order transition below about 20 A$``$GeV is dominated by the baryon density while especially at SPS or even RHIC energies a chiral order transition is due to ’temperature’.
The question now arises, if there are proper experimental observables that allow to trace down such type of phase transition (or cross over). When gating on central collisions of Au + Au (or Pb + Pb) such phenomena should show up in their excitation functions. We note that in a pure hadronic transport approach we expect a smooth behaviour of practically all observables with bombarding energy due to an increase of thermal excitation energy. This is not so obvious for the HSD approach where a gradual transition from hadronic excitations to strings and quark (or diquark) degrees of freedom is involved. In fact, in Ref. it has been claimed that the excitation function of transverse and elliptic proton flow together with the transverse $`p_t`$ spectra of protons suggest a transition from hadron to string matter at about 5 A$``$GeV. Here, however, we concentrate on meson abundancies and spectra and refer the reader to Ref. for the collective dynamical aspects and to Ref. for the thermal properties of the theory that involves a limiting (’Hagedorn’) temperature of $`T_S`$ 150 MeV due to the continuum string excitations.
## 3 Meson production
To present a general overview on the meson abundancies in nucleus-nucleus collisions we show in Fig. 5 the meson multiplicities for central collisions of Au + Au from SIS to RHIC energies. All multiplicities for $`\pi ^+,\eta ,K^+,K^{},\varphi `$ as well as $`J/\mathrm{\Psi }`$ mesons show a monotonic increase with bombarding energy which is only very steep at ’subthreshold’ energies, i.e. at bombarding energies per nucleon below the threshold in free space for $`NN`$ collisions. At higher bombarding energies the meson abundancies group according to their quark content, i.e. the multiplicities are reduced (relative to $`\pi ^+`$) by about a factor of 4–5 for a strange quark, a factor of $``$ 50 for $`s\overline{s}\varphi `$ and a factor of $`210^5`$ for $`c\overline{c}J/\mathrm{\Psi }`$. We mention that in these calculations the meson rescattering and reabsorption processes have been taken into account; this reduces the $`J/\mathrm{\Psi }`$ cross section at RHIC energies by about a factor of 10 (cf. Section 5). At ’subthreshold’ energies the $`\varphi `$ multiplicity turns out to be almost the same as the antikaon multiplicity, but then rises less steeply with bombarding energy. Apart from the $`\varphi `$ excitation function – that still has to be controlled experimentally – we thus find no pronounced change in the shape of the meson abundancies up to RHIC energies where data are expected to come up soon.
We recall that our investigations on strangeness production up to 2 A$``$GeV have given some evidence for attractive antikaon potentials in the medium while for kaons only a very moderate repulsive potential was suggested ; $`\eta `$ mesons apparently do not show sensible in-medium effects according to the studies in Ref. in comparison to the available experimental spectra. At AGS and SPS energies, on the other hand, the potential effects on kaon and antikaon abundancies have been found to be only very low such that meson potentials have been discarded in Fig. 5 for the overview.
Since the meson abundancies show no sudden change in the excitation function, we turn to particle ratios. Here strangeness enhancement has been suggested for more than 2 decades to possibly qualify as a sensible probe for a QGP phase (cf. Ref. for a recent overview). Here we examine the $`K^+/\pi ^+`$ ratio at midrapidity in central Au + Au (or Pb + Pb) collisions where experimental data are now available from SIS to SPS energies . We recall that detailed comparisons of pion and kaon rapidity distributions and transverse momentum spectra to the available data have been presented in Refs. such that we directly can continue with the corresponding excitation function.
In order to discuss the strangeness production over the complete energy range from SIS to RHIC energies we show in Fig. 6 the calculated $`K^+/\pi ^+`$ ratio at midrapidity (-0.5 $`y_{cm}`$ 0.5) for central Au + Au collisions in comparison to the experimental data. This ratio experimentally is substantially lower at SPS energies ($`13.5\%`$) compared to AGS energies ($`19\%`$). At SPS energies this ratio is only enhanced by a factor 1.75 for central Pb + Pb collisions relative to p + p reactions and has to be compared to the factor $``$ 3 at AGS. Such a decrease of the scaled kaon yield from AGS to SPS energies is not described by the HSD transport model (without kaon self-energies) which shows a monotonic increase with bombarding energy similar to $`pp`$ reactions (open circles). Furthermore, the higher temperatures and particle densities at SPS energies tend to enhance the $`K^+/\pi ^+`$ yield closer to its thermal equilibrium value of $`2025\%`$ at chemical freezout and temperatures of $`T150`$ MeV.
Our findings have to be compared to results obtained by other microscopic approaches. Here only the RQMD model very recently provides a study partly comparable to that of Ref. , however, excluding p + A reactions to control the amount of $`K^+`$ production by rescattering. Whereas kaon production in p + p reactions is comparable to our results in (input of the transport model), the RQMD model yields a higher $`K^+/\pi ^+`$ ratio in Au + Au, Pb + Pb collisions at all energies due to a higher kaon yield from rescattering. The latter can be attributed to ’high mass strange resonances’ that have been incorporated to describe the low energy kaon production via resonance production and decay (s-channels). Since these ’high mass resonances’ can be repopulated in resonance-resonance scattering, a rather high $`\sqrt{s}`$ is concentrated in a single degree of freedom for a short time. Such ’hot spots’ then lead to a higher $`K^+/\pi ^+`$ ratio in A + A reactions especially at AGS energies. While the $`K^+/\pi ^+`$ ratio can be reasonably described at AGS energies in Au + Au reactions, this ratio is overestimated significantly at SPS energies (cf. Fig. 4 of Ref. ). Thus also the RQMD model does not describe the relative decrease of the $`K^+/\pi ^+`$ ratio from 11 A$``$GeV to 160 A$``$GeV. The results of the RQMD calculations at SIS energies are not known to the authors.
## 4 Dilepton production
Electromagnetic decays to virtual photons ($`e^+e^{}`$ or $`\mu ^+\mu ^{}`$ pairs) have been suggested long ago to serve as a possible signature for a phase transition to the QGP or to be an ideal probe for vector meson spectroscopy in the nuclear medium. As pointed out in Refs. the isovector current-current correlation function is proportional to the imaginary part of the $`\rho `$-meson propagator and also to the dilepton invariant mass spectra. Dileptons are particularly well suited for an investigation of the violent phases of a high-energy heavy-ion collision because they can leave the reaction volume essentially undistorted by final-state interactions. Indeed, the dilepton studies in heavy-ion collisions by the DLS Collaboration at the BEVALAC and by the CERES , HELIOS , NA38 and NA50 Collaborations at SPS energies have found a vivid interest in the nuclear physics community.
We recall that the question of chiral symmetry restoration does not necessarily imply that vector meson masses have to drop with baryon density or temperature . Actually, chiral symmetry restoration (ChSR) only demands that the isovector current-current correlation function and the axial vector current-current correlation function (dominated by the chiral partner of the $`\rho `$, the $`a_1`$-meson) should become identical at high $`\rho _B`$ or temperature $`T`$, respectively, because there should be no more differences between left- and right-handed particles or equivalently vector and axial vector currents . Thus also a strong broadening of the $`\rho `$\- as well as the $`a_1`$-spectral function and their mixing in the medium can be considered as a signature for ChSR which, however, is not easy to detect experimentally.
In Ref. we have demonstrated that the present experimental data on the low mass dilepton enhancement at SPS energies can be described equally well within the ’dropping $`\rho `$-mass’ scenario as well as within the ’melting’ $`\rho `$ picture, which implies a large spreading in mass of the $`\rho `$ spectral function due to its couplings to baryons and/or mesons. The situation at SIS/BEVALAC energies is more ’puzzling’ since here the low mass dilepton enhancement is neither described in the dynamical spectral function approach nor in the ’dropping mass’ scheme , though the $`pp`$ dilepton data from the DLS collaboration are reproduced within the known sources rather well from 1 – 5 GeV bombarding energy . In short, the dynamical origin of the low mass dilepton enhancement is not yet understood.
Here we propose to investigate the excitation function of low mass dilepton enhancement in central Au + Au collisions. As discussed in Refs. this excess of dileptons is most probably due to the isovector ($`\rho `$) channel, i.e. the imaginary part of the isovector current-current correlation function which, however, mixes with the axial vector current-current correlation function at finite temperature and baryon density .
Since the $`\rho `$-meson spectral function is of primary interest, it is important to have some information about the actual baryon densities that the $`\rho `$-meson experiences during its propagation and decay in central Au + Au collisions. This information is displayed in Fig. 7 – as resulting from the HSD transport calculation – for central reactions at 2, 5, 10, 20, 40 and 160 A$``$GeV. Here the meson-baryon production channels and baryon-baryon production channels are summed up by the solid lines and denoted as ($`\pi B\rho ,BB\rho `$). The meson-meson production channels such as $`\pi \pi \rho ,a_1\pi \rho `$ etc. are summed up in the dashed histograms indicated as $`\pi \pi \rho `$ according to the dominating channel. We find that especially the initial $`BB`$ production channels produce $`\rho `$-mesons at very high densities close to 2 $`\rho _0\gamma _{cm}`$ where $`\gamma _{cm}`$ is the Lorentz factor in the cms. However, these production channels are much less abundant than the meson-meson channels which extend over a larger time span (for the heavy system Au + Au) and essentially occur at much lower baryon density, respectively. This effect becomes even more pronounced with increasing bombarding energy, i.e. 40 – 160 A$``$GeV, where most of the $`\rho `$-mesons are produced at baryon densities below 2 $`\rho _0`$. Since in Fig. 7 the relative abundancy of $`\rho `$-mesons is displayed as a function of the baryon density at the production (formation) point, the time averaged value of the density is even lower due to a fast expansion of the hadronic fireball. Thus to probe on average high baryon densities by $`\rho `$-mesons one should step down in energy to 2 – 5 A$``$GeV in order to optimize effects due to the coupling to baryons.
A general overview of dilepton mass spectra in central ($`b`$=2 fm) Au + Au collisions from 2 A$``$GeV to 21.5 A$``$TeV is presented in Fig. 8, where the upper part corresponds to the case of vacuum spectral functions for all mesons (cf. Fig. 8.20 of Ref. ), while the lower part is calculated by employing the $`\rho `$ spectral function from Rapp et al. . Whereas for the free meson spectral function one observes essentially an increase of the dilepton production channels with bombarding energy without any substantial change in the spectral shape (except for an increasing peak from the $`\varphi `$; cf. Fig. 5), the in-medium calculations yield almost exponential mass spectra above $`M`$ 0.4 GeV with small peaks from vacuum $`\omega `$ and $`\varphi `$ decays at $`M0.78`$ GeV and 1.02 GeV.
We mention that for the vacuum spectral functions (upper part of Fig. 8) the shape of the dilepton spectra (for $`M_{e^+e^{}}`$ 0.15 GeV) is due to the superposition of $`\eta ,\eta ^{},\omega `$ and $`a_1`$ Dalitz decays and direct vector meson decays ($`\rho ,\omega ,\varphi `$), where all mesons may also be produced in meson-baryon and meson-meson collisions. The increase in dilepton yield (for $`M_{e^+e^{}}`$ 0.15 GeV) with bombarding energy thus is due to an enhanced production of $`\eta ,\eta ^{},\rho ,\omega ,\varphi ,a_1`$ mesons. Their relative abundance from SIS to RHIC energies does not scale directly with the charged particle multiplicity which is dominated by protons, $`\pi ^\pm `$ and $`K^\pm `$ (cf. Fig. 5). However, above about 50 – 100 A$``$GeV the meson ratios do not change very much (cf. Figs. 5,6) while the charged particle multiplicity becomes dominated by $`\pi ^\pm `$ and $`K^\pm `$. Thus from SPS to RHIC energies the low mass dilepton yield should approximately be proportional to the charged particle multiplicity, too.
The relative change in the dilepton spectra is quantitatively displayed in Fig. 9 – again for central ($`b`$= 2 fm) collisions of Au + Au – for the case of free meson spectral functions (solid lines) and the spectral function from Rapp et al. (dashed lines). In line with the discussion above the most prominent spectral changes are expected at rather low bombarding energies from 2–10 A$``$GeV, where the enhancement in the invariant mass range from 0.3–0.6 GeV is about a factor of 4. Note that the $`\omega `$ and $`\varphi `$ vacuum decays show clear peaks on top of the ’exponential’ continuum, which will have to be identified experimentally. We mention that in all our calculations we have incorporated an experimental (’optimistic’) mass resolution of $`\mathrm{\Delta }M`$ = 10 MeV for the dilepton invariant mass.
## 5 Charmonium production and suppression
Matsui and Satz have proposed that a suppression of the $`J/\mathrm{\Psi }`$ yield in ultra-relativistic heavy-ion collisions is a plausible signature for the formation of the quark-gluon plasma because the $`J/\mathrm{\Psi }`$ should dissolve in the QGP due to color screening. This suggestion has stimulated a number of heavy-ion experiments at CERN SPS to measure the $`J/\mathrm{\Psi }`$ via its dimuon decay. Indeed, these experiments have shown a significant reduction of the $`J/\mathrm{\Psi }`$ yield when going from proton-nucleus to nucleus-nucleus collisions . Especially for Pb + Pb at 160 A$``$GeV an even more dramatic reduction of $`J/\mathrm{\Psi }`$ has been reported by the NA50 Collaboration .
To interpret the experimental results, various models based on $`J/\mathrm{\Psi }`$ absorption by hadrons have been also proposed (cf. Refs. for recent reviews) that do not involve the assumption of a QGP phase transition. The role of comover dissociation is presently again heavily debated especially since theoretical calculations for $`J/\mathrm{\Psi }`$-meson dissociation cross sections differ by up to a factor of 50 . The problem is even more complicated since the $`J/\mathrm{\Psi }`$ meson is not created instantly as a hadronic state and there is also a substantial ($``$ 35 %) feeding from the $`\chi _c\gamma `$ decays. Moreover, it is expected that the $`c\overline{c}`$ pair is first produced in a color-octet state together with a gluon (’pre-resonance state’) and that this more extended configuration has a larger interaction cross section with baryons and mesons before the $`J/\mathrm{\Psi }`$ singlet state, $`\mathrm{\Psi }^{}`$ or $`\chi _c`$ finally emerges after some formation time $`\tau _c`$. Additionally, the dissociation on mesons of formed $`J/\mathrm{\Psi }`$’s will differ from $`\chi _c`$ due to their different thresholds with respect to the $`D\overline{D}`$ channel as well as for the $`\mathrm{\Psi }^{}`$. Since experimental information on the various charmonium-meson cross sections – especially at low relative momenta – will be hard to obtain, the excitation function of charmonium suppression in central nucleus-nucleus collisions might be exploited to obtain additional information on the absorption scenarios. One expects that quite below the bombarding energy necessary for the formation of a QGP phase the charmonium absorption should be entirely due to dissociation with hadrons; any additional suppression due to color screening then will show up in a more rapid suppression with the incident energy or with the energy density achieved .
In this Section we calculate the excitation function for $`J/\mathrm{\Psi }`$ suppression within two absorption scenarios, i.e. the ’early’ and ’late’ comover dissociation models, which have been explored in detail by our group before at SPS energies . The ’early’ comover absorption scenario is based on the idea, that a $`c\overline{c}`$ pair is created in the initial ’hard’ phase of the nucleus-nucleus collision, where the string density is very high, and the $`c\overline{c}`$ pair is dissolved in the color electric field of neighboring strings. Since in the HSD approach the information on the string density as well as the string space-time extension is available, the absorption model has only a single parameter, i.e. the average transverse dimension of an extending string. In Ref. a string radius of 0.2 fm was found to describe simultaneously the data for p + A and S + U at 200 A$``$GeV from NA38 and for Pb + Pb at 160 A$``$GeV from NA50 when adopting the conventional charmonium-nucleon dissociation cross section of 6 mb. It has been also speculated that the overlap of strings due to percolation might describe the phase transition to a QGP .
In the ’late’ comover scenario the additional charmonium suppression is due to charmonium-meson scattering to $`D\overline{D}`$ with an average charmonium-meson cross section of $``$ 3 mb . Cross sections of this order have been calculated by Haglin in Ref. within a meson-exchange model and thus might appear not unrealistic. In our present calculation we refer to the model II of Ref. including a ’pre-resonance’ charmonium life time of 0.3–0.5 fm/c which is supported (within the errorbars) by a more recent analysis of charmonium suppression as a function of the Feynman variable $`x_F`$ in p + A reactions from He et al. and Kharzeev et al. . We recall that within the model II of Ref. the $`J/\mathrm{\Psi }`$ suppression data for p + A and S + U at 200 A$``$GeV from NA38 and for Pb + Pb at 160 A$``$GeV from NA50 have been described very well when adopting a ’pre-resonance’-nucleon cross section of 6 mb and a $`J/\mathrm{\Psi }`$-nucleon cross section of 3–4 mb in line with the data on $`J/\mathrm{\Psi }`$ photoproduction . Independent dynamical studies on charmonium suppression within the UrQMD model later on have lead to very similar conclusions.
Since the comparison of our calculations at SPS energies has been performed to data taken in 1995 and before we show in Fig. 10 a comparison of the ’early’ comover model (dashed line) and the ’late’ comover model II with the more recent data from NA50 for Pb + Pb at 160 A$``$GeV using the explicit numbers from Refs. <sup>1</sup><sup>1</sup>1This comparison is necessary since the 1995 data have been rescaled in and our calculations are reported inconsistently in the more recent presentations of this topic .. The ’early’ comover absorption model here gives a little too low suppression at high $`E_T`$ whereas the ’late’ comover absorption model is still in line with the more recent data from 1996 and 1998 with minimum bias (open triangles and open circles). The data from 1996 (full squares), that show a (much debated) two-step behaviour, do not agree with the explicit $`E_T`$ dependence from our calculations; however, the 1996 minimum bias data (open triangles) well match for $`E_T`$ 40 GeV whereas the $`J/\mathrm{\Psi }`$ suppression is slightly overestimated at lower $`E_T`$. We do not comment on the highest $`E_T`$ data points from 1998 with minimum bias since our earlier analysis did not extend to these specific events.
The question now arises, if the excitation function for $`J/\mathrm{\Psi }`$ suppression in central Au + Au collisions might show some unusual behaviour within the two scenarios discussed above or how they might be disentangled. In order to achieve the same suppression factor in central collisions within the ’early’ comover model we have increased the string absorption radius from $`r_s`$ = 0.2 fm to 0.22 fm to get the same value for the $`J/\mathrm{\Psi }`$ survival factor $`S_{J/\mathrm{\Psi }}`$ at 160 A$``$GeV. We note that the total $`J/\mathrm{\Psi }`$ multiplicity shown in Fig. 5 has been calculated within the ’late’ comover model. It drops by almost 3 orders of magnitude when decreasing the bombarding energy from 160 A$``$GeV to 20 A$``$GeV. At the lowest energy considered here the experimental $`J/\mathrm{\Psi }`$ signal will be very hard to measure; it is hopeless within the present experimental setups. Nevertheless, it is worth exploring theoretically if some unusual excitation function might be found. Our results are displayed in Fig. 11 (l.h.s.) for the ’early’ absorption model (open circles) and for the ’late’ comover model (full circles); both models practically do not differ in their excitation functions and show a very smooth decrease of $`S_{J/\mathrm{\Psi }}`$ from 0.4 to 0.3 with increasing bombarding energy. The net absorption by baryons is dominant in both scenarios, however, differs in magnitude due to the simple fact that in the ’early’ (string) absorption model there is less suppression by baryons since the absorption by strings competes at early times. In the ’late’ comover model there is more absorption by baryons because the mesons are formed at later stages and not competitive in the early phase; their relative contribution is lower as for strings accordingly. However, these individual contributions cannot be distinguished experimentally and thus are ’irrelevant’.
The r.h.s. of Fig. 11 shows the survival probability $`S_{J/\mathrm{\Psi }}`$ in the ’late’ comover model for central Au + Au collisions from 0.160 A$``$TeV to 21.5 A$``$TeV, respectively, where we have gated on $`J/\mathrm{\Psi }`$’s in the rapidity interval -1 $`y_{cm}`$ 1. The solid line stands for the total $`J/\mathrm{\Psi }`$ survival probability while the dashed line displays the relative absorption on baryons and the dotted line the relative dissociation on mesons. Whereas the dissociation on baryons is practically constant with bombarding energy, the absorption on mesons increases in line with the higher meson densities achieved with increasing $`E/A`$. We note that the ’early’ comover model leads to a similar total absorption within the numerical accuracy.
We have to mention that neither the ’early’ nor the ’late’ comover model might be realized in nature exclusively and both absorption processes should occur within the same reaction with probably different weights. Since we do not find a substantial difference for both scenarios also a linear combination of both absorption models, i.e. decreasing $`r_s`$ as well as the charmonium-meson cross section accordingly, will lead to a similar excitation function. This also holds for the relative suppression on the transverse energy $`E_T`$ (cf. Fig. 10).
Inspite of the rather disappointing perspectives to disentangle the ’late’ and ’early’ comover models experimentally at the full range of SPS energies, the excitation function of $`J/\mathrm{\Psi }`$ still might show some discontinuity in $`E/A`$ experimentally, which could rule out the two models studied here and indicate a transition to a QGP phase. This subject is taken up in the next Section again with respect to the dependence of $`S_{J/\mathrm{\Psi }}`$ on the transverse energy $`E_T`$ produced in Au + Au collisions at RHIC energies.
## 6 Predictions for RHIC energies
As noted in the introduction one expects that the initial nonequilibrium phase of a nucleus-nucleus collision at RHIC energies should be described by parton degrees of freedom, whereas hadrons are only formed (by ’condensation’) at a later stage of the reaction and interact until freeze out. Thus parton cascade calculations should be adequate for all initial reactions involving a large 4-momentum transfer between the constituents, while hadron cascades should be appropriate in the final hadronic expansion phase. We suggest that the dynamics in between the partonic and hadronic phase might be described by quarks (diquarks) and strings as e.g. implemented in the HSD approach.
The practical question is, however, if nonequilibrium partonic and hadron/string models can be distinguished at all, i.e. do they lead to different predictions for experimental observables? In fact, first applications of the parton cascade model developed by Geiger to nucleus-nucleus collisions at SPS energies have suggested that a reasonable description of the meson and baryon rapidity distributions can also be achieved on the basis of partonic degrees of freedom. However, in the latter calculations the extrapolation of the strong coupling constant to low $`Q^2`$ has been overestimated as discovered recently by Bass and Müller . This finding invalidates the detailed predictions and comparisons within the hybrid models VNI+UrQMD or VNI+HSD presented in Ref. that have been tailored to describe the dynamics at RHIC or even LHC energies.
We start with $`pp`$ collisions at $`\sqrt{s}`$ = 200 GeV. The calculated results for the proton, $`\pi ^+`$ and $`K^+`$ rapidity distributions in the cms are shown in Fig. 12 (upper part) for both models, which are denoted individually by the labels VNI and HSD in obvious notation. On the level of $`pp`$ collisions we find only minor differences between the two kinetic models. The parton cascade shows a slightly higher amount of proton stopping as the HSD model (l.h.s.) and as a consequence a slightly higher production of $`\pi ^+`$ and $`K^+`$ mesons (r.h.s.), because the energy taken from the relative motion of the leading baryons is converted to the production of mesons. It is presently unclear which of the two approaches will be closer to experiment; a proper description of $`pp`$ data will be a necessary step before performing reliable extrapolations to nucleus-nucleus collisions.
Inspite of this missing experimental information we directly step towards central collisions ($`b`$ 1.5 fm) for Au + Au at $`\sqrt{s}`$ = 200 GeV. The calculated results for the net proton (here $`p\overline{p}`$), antiproton, $`\pi ^+`$ and $`K^+`$ rapidity distributions in the cms are shown in Fig. 12 (lower part). In the HSD scenario essentially ’comover’ scattering occurs with a low change of the meson rapidity distribution. Thus the meson rapidity distributions are roughly the same as for $`pp`$ collisions. Also note that at midrapidity the net baryon density $`N_pN_{\overline{p}}`$ is practically zero, however, even at midrapidity at lot of baryons appear that are produced together with antibaryons. Thus also mesons (especially $`c\overline{c}`$ pairs) encounter a lot of baryons and antibaryons on their way to the continuum. Whereas the HSD approach predicts a vanishing net baryon density at midrapidity, other recent models – that combine high and low energy transport concepts – suggest a sizeable net proton density for $`y_{cm}`$ 0 .
The amount of higher order hadronic rescattering processes at RHIC energies is depicted in Fig. 13 (lower right part) as emerging from the HSD calculation, where the number of baryon-baryon ($`BB`$) and meson-baryon collisions ($`mB`$) is shown as a function of the invariant energy $`\sqrt{s}`$. We mention that quark-baryon and diquark-baryon collisions are counted here as $`mB`$ or $`BB`$ collisions, respectively. Apart from the initial small peak at $`\sqrt{s}`$ = 200 GeV a substantial amount of intermediate and low energy rescattering processes with maxima at 2.5 GeV and 1.8 GeV are found, which essentially stand for flavor exchange processes, multiple pion production in $`mB`$ and $`BB`$ collisions as well as secondary strangeness production channels. For comparison the corresponding $`\sqrt{s}`$ distributions are also displayed for bombarding energies of 2, 11 and 160 A$``$GeV, respectively, showing a dominance of low energy $`BB`$ and $`mB`$ collisions. The latter reactions occur at energy scales where perturbative QCD is no longer applicable. This has to be kept in mind additionally when comparing to $`pp`$ and $`pA`$ reactions at $`\sqrt{s}`$ = 200 GeV.
We return to the question of charmonium suppression at RHIC energies since it is expected that one might probe increasing energy densities also with increasing centrality of the collision, where the latter can be correlated with the transverse energy $`E_T`$ produced in a collision event. As argued e.g. by Satz the survival factor $`S_{J/\mathrm{\Psi }}`$ then should show steps as a function of $`E_T`$ due to the melting of first the $`\chi _c`$ and then the $`J/\mathrm{\Psi }`$ in a QGP phase. As seen from Fig. 10 there are no pronounced steps in the $`E_T`$ dependence of $`J/\mathrm{\Psi }`$ suppression in the data for Pb + Pb at SPS energies according to the authors point of view; this situation might change at RHIC energies.
Using the ’late’ comover model described in Section 5 we have calculated the $`J/\mathrm{\Psi }`$ survival factor $`S_{J/\mathrm{\Psi }}`$ as a function of the transverse energy $`E_T`$ in the cms rapidity window \[-1,1\] for Au + Au at $`\sqrt{s}`$ = 200 GeV on an event by event basis covering all impact parameters from $`b`$ = 0 to 13 fm. The resulting correlation of $`S_{J/\mathrm{\Psi }}`$ with $`E_T`$ is shown in Fig. 14 and indicates a smooth decrease with centrality (or increasing $`E_T`$) reaching an average survival probability of $``$ 0.1 for the most central events (cf. Fig. 11, r.h.s.). This result can be understood as follows: According to our calculations the net $`J/\mathrm{\Psi }`$ dissociation by mesons in central Au + Au collisions at the SPS is $`16\%`$ (cf. Fig. 11) while the rapidity distribution of negatively charged particles ($`h^{}`$) at midrapidity here is about 180. At the RHIC energy we get a corresponding $`h^{}`$ rapidity density at midrapidity of $``$ 450 (cf. Fig. 12) which is higher by a factor of 2.5. Simply multiplying the $`J/\mathrm{\Psi }`$ meson absorption at the SPS of $`16\%`$ by the factor 2.5 we obtain about $`40\%`$ for central collisions at RHIC energies, which together with $`52\%`$ of absorption on baryons gives a survival probability of $`8\%`$. The actual numerical results in Fig. 14 indicate that this simple estimate works quite well. On the other hand, if the $`h^{}`$ rapidity density is found to be lower (higher) experimentally, we expect corresponding changes in the $`J/\mathrm{\Psi }`$ suppression for central events if the ’late’ comover absorption model holds true. This dependence might well be tested experimentally in the near future to possibly falsify the comover dissociation model.
## 7 Summary
In this work we have performed a systematic analysis of hadron production in central Au + Au collisions from SIS to RHIC energies within the HSD transport approach. We have concentrated here on the ’classical’ signatures, i.e. strangeness and low mass dilepton enhancement as well as charmonium suppression. For all observables our calculations give a monotonic increase (for the ratio $`K^+/\pi ^+`$ and charmonium suppression) or decrease (for the low mass dilepton enhancement) with bombarding energy, respectively. So far, experimental data are available only in a limited range of bombarding energies or at a single energy, respectively. We have pointed out that the relative maximum indicated by the experimental data in the $`K^+/\pi ^+`$ ratio at about 10 A$``$GeV (or higher?) is not reproduced within the transport approach that is based on quark, diquark, string and hadronic degrees of freedom. We speculate that at AGS energies this failure might be attributed to a restoration of chiral symmetry in a sufficiently large space-time volume (cf. Figs. 3 and 4).
The enhancement of low mass dileptons in the range 0.3 GeV $`M_{e^+e^{}}`$ 0.6 GeV is most pronounced at lower bombarding energies of 2–5 A$``$GeV within our calculations since here the space-time volume for densities above 2 $`\rho _0`$ is very large such that a majority of $`\rho `$-mesons probes the high density phase of the reaction (cf. Fig. 5). With increasing bombarding energy the average density – which a $`\rho `$-meson experiences – drops substantially such that high energy nucleus-nucleus collisions are not well suited for in-medium $`\rho `$ spectroscopy.
The suppression of charmonium (here $`J/\mathrm{\Psi }`$) increases smoothly with bombarding energy and centrality of the reaction within the ’early’ and ’late’ comover absorption scenarios. Unfortunately, both scenarios cannot be distinguished by means of the excitation function since they give approximately the same survival probability $`S_{J/\mathrm{\Psi }}`$ with bombarding energy (cf. Fig. 11). In the transport approach the smooth increase of charmonium absorption with bombarding energy is easy to understand: a major fraction of $`J/\mathrm{\Psi }`$’s is anyhow dissociated by baryons which basically are of the same number at all energies considered here; only the relative collisional energy changes. The additional absorption by ’early’ strings or ’late’ hadrons increases smoothly with bombarding energy since the string and hadron density increases accordingly. At RHIC energies this additional suppression mechanism leads to a $`J/\mathrm{\Psi }`$ suppression of about 90% in central Au + Au collisions even without employing an explicit formation of a QGP. We note, however, that the charmonium suppression shows a smooth dependence on the transverse energy $`E_T`$ (cf. Fig. 14); any gradual steps of $`S_{J/\mathrm{\Psi }}`$ with $`E_T`$ due to a melting of the $`\chi _c`$ or the $`J/\mathrm{\Psi }`$ at higher energy density would indicate a new suppression mechanism which might be attributed to color screening in a QGP phase .
The authors acknowledge inspiring discussions with J. Aichelin, S. A. Bass, G. E. Brown, C. Greiner, M. Gyulassy, C. M. Ko, U. Mosel, R. Rapp, H. Satz, H. Sorge, H. Stöcker, J. Wambach, X.-N. Wang and K. Werner. |
warning/0001/gr-qc0001098.html | ar5iv | text | # Stochastic semiclassical fluctuations in Minkowski spacetime
## I Introduction
It has been pointed out that the semiclassical theory of gravity cannot provide a correct description of the dynamics of the gravitational field in situations where the quantum stress-energy fluctuations are important . In such situations, these fluctuations may have relevant back-reaction effects in the form of induced gravitational fluctuations which, in a certain regime, are expected to be described as classical stochastic fluctuations. A generalization of the semiclassical theory is thus necessary to account for these effects. In two previous papers, Refs. and , we have shown how a stochastic semiclassical theory of gravity can be formulated to improve the description of the gravitational field when stress-energy fluctuations are relevant.
In Ref. , we adopted an axiomatic approach to construct a perturbative generalization of semiclassical gravity which incorporates the back reaction of the lowest order stress-energy fluctuations in the form of a stochastic correction. We started noting that, for a given solution of semiclassical gravity, the lowest order matter stress-energy fluctuations can be associated to a classical stochastic tensor. We then sought a consistent equation in which this stochastic tensor was the source of linear perturbations of the semiclassical metric. The equation obtained is the so-called semiclassical Einstein-Langevin equation.
In Ref. , we followed the idea, first proposed by Hu in the context of back reaction in semiclassical gravity, of viewing the metric field as the “system” of interest and the matter fields (modeled in that paper by a single scalar field) as being part of its “environment.” We then showed that the semiclassical Einstein-Langevin equation introduced in Ref. can be formally derived by a method based on the influence functional of Feynman and Vernon (see also Ref. ). That derivation shed light into the physical meaning of the semiclassical Langevin-type equations around specific backgrounds previously obtained with the same functional approach , since the stochastic source term was shown to be closely linked to the matter stress-energy fluctuations. We also developed a method to compute the semiclassical Einstein-Langevin equation using dimensional regularization, which provides an alternative and more direct way of computing this equation with respect to previous calculations.
This paper is intended to be a first application of the full stochastic semiclassical theory of gravity, where we evaluate the stochastic gravitational fluctuations in a Minkowski background. In order to do so, we first use the method developed in Ref. to derive the semiclassical Einstein-Langevin equation around a class of trivial solutions of semiclassical gravity consisting of Minkowski spacetime and a linear real scalar field in its vacuum state, which may be considered the ground state of semiclassical gravity. Although the Minkowski vacuum is an eigenstate of the total four-momentum operator of a field in Minkowski spacetime, it is not an eigenstate of the stress-energy operator. Hence, even for these solutions of semiclassical gravity, for which the expectation value of the stress-energy operator can always be chosen to be zero, the fluctuations of this operator are non-vanishing. This fact leads to consider the stochastic corrections to these solutions described by the semiclassical Einstein-Langevin equation.
We then solve the Einstein-Langevin equation for the linearized Einstein tensor and compute the associated two-point correlation functions. Even though, in this case, we expect to have negligibly small values for these correlation functions at the domain of validity of the theory, i.e., for points separated by lengths larger than the Planck length, there are several reasons why we think that it is worth carrying out this calculation.
On the one hand, these are, to our knowledge, the first solutions obtained to the full semiclassical Einstein-Langevin equation. We are only aware of analogous solutions to a “reduced” version of this equation inspired in a “mini-superspace” model . There is also a previous attempt to obtain a solution to the Einstein-Langevin equation in Ref. , but, there, the non-local terms in the Einstein-Langevin equation were neglected.
The Einstein-Langevin equations computed in this paper are simple enough to be explicitly solved and, at least for the case of a conformal field, the expressions obtained for the correlation functions can be explicitly evaluated in terms of elementary functions. Thus, our calculation can serve as a testing ground for the solutions of the Einstein-Langevin equation in more complex situations of physical interest (for instance, for a Robertson-Walker background and a field in a thermal state).
On the other hand, the results of this calculation, which confirm our expectations that gravitational fluctuations are negligible at length scales larger than the Planck length, can be considered as a first check that stochastic semiclassical gravity predicts reasonable results.
In addition, we can extract conclusions on the possible qualitative behavior of the solutions to the Einstein-Langevin equation. Thus, it is interesting to note that the correlation functions are characterized by correlation lengths of the order of the Planck length; furthermore, such correlation lengths enter in a non-analytic way in the correlation functions. This kind of non-analytic behavior is actually quite common in the solutions to Langevin-type equations with dissipative terms and hints at the possibility that correlation functions for other solutions to the Einstein-Langevin equation are also non-analytic in their characteristic correlation lengths.
The plan of the paper is the following. In Sec. II, we give a brief overview of the method developed in Ref. to compute the semiclassical Einstein-Langevin equation. We then consider the background solutions of semiclassical gravity consisting of a Minkowski spacetime and a real scalar field in the Minkowski vacuum. In Sec. III, we compute the kernels which appear in the Einstein-Langevin equation. In Sec. IV, we derive the Einstein-Langevin equation for metric perturbations around Minkowski spacetime. As a side result, we obtain some semiclassical results, which include the expectation value of the stress-energy tensor of a scalar field with arbitrary mass and arbitrary coupling parameter to linear order in the metric perturbations, and also some results concerning the production of particles by metric perturbations: the probability of particle creation and the number and energy of created particles. In Sec. V, we solve this equation for the components of the linearized Einstein tensor and compute the corresponding two-point correlation functions. For the case of a conformal field and spacelike separated points, explicit calculations show that the correlation functions are characterized by correlation lengths of the order of the Planck length. We conclude in Sec. VI with a discussion of our results. We also include some appendices with technical details used in the calculations.
Throughout this paper we use the $`(+++)`$ sign conventions and the abstract index notation of Ref. , and we work with units in which $`c=\mathrm{}=1`$.
## II Overview
In this section, we give a very brief summary of the main results of Refs. and which are relevant for the computations in the present paper. One starts with a solution of semiclassical gravity consisting of a globally hyperbolic spacetime $`(,g_{ab})`$, a linear real scalar field quantized on it and some physically reasonable state for this field (we work in the Heisenberg picture). According to the stochastic semiclassical theory of gravity , quantum fluctuations in the stress-energy tensor of matter induce stochastic linear perturbations $`h_{ab}`$ to the semiclassical metric $`g_{ab}`$. The dynamics of these perturbations is described by a stochastic equation called the semiclassical Einstein-Langevin equation.
Assuming that our semiclassical gravity solution allows the use of dimensional analytic continuation to define regularized matrix elements of the stress-energy “operator,” we shall write the equations in dimensional regularization, that is, assuming an arbitrary dimension $`n`$ of the spacetime. Using this regularization method, we use a notation in which a subindex $`n`$ is attached to those quantities that have different physical dimensions from the corresponding physical quantities. The $`n`$-dimensional spacetime $`(,g_{ab})`$ has to be a solution of the semiclassical Einstein equation in dimensional regularization
$$\frac{1}{8\pi G_B}\left(G^{ab}[g]+\mathrm{\Lambda }_Bg^{ab}\right)(\frac{4}{3}\alpha _BD^{ab}+2\beta _BB^{ab})[g]=\mu ^{(n4)}\widehat{T}_n^{ab}[g],$$
(1)
where $`G_B`$, $`\mathrm{\Lambda }_B`$, $`\alpha _B`$ and $`\beta _B`$ are bare coupling constants and $`G_{ab}`$ is the Einstein tensor. The tensors $`D^{ab}`$ and $`B^{ab}`$ are obtained by functional derivation with respect to the metric of the action terms corresponding to the Lagrangian densities $`R_{abcd}R^{abcd}R_{ab}R^{ab}`$ and $`R^2`$, respectively, where $`R_{abcd}`$ is the Riemann tensor, $`R_{ab}`$ is the Ricci tensor and $`R`$ is the scalar curvature (see Ref. for the explicit expressions for the tensors $`D^{ab}`$ and $`B^{ab}`$). In the last equation, $`\widehat{T}_n^{ab}`$ is the stress-energy “operator” in dimensional regularization and the expectation value is taken in some state for the scalar field in the $`n`$-dimensional spacetime. Writing the bare coupling constants in Eq. (1) as renormalized coupling constants plus some counterterms which absorb the ultraviolet divergencies of the right hand side, one can take the limit $`n4`$, which leads to the physical semiclassical Einstein equation.
Assuming that $`g_{ab}`$ is a solution of Eq. (1), the semiclassical Einstein-Langevin equation can be similarly written in dimensional regularization as
$$\frac{1}{8\pi G_B}\left(G^{ab}[g+h]+\mathrm{\Lambda }_B\left(g^{ab}h^{ab}\right)\right)(\frac{4}{3}\alpha _BD^{ab}+2\beta _BB^{ab})[g+h]=\mu ^{(n4)}\widehat{T}_n^{ab}[g+h]+2\mu ^{(n4)}\xi _n^{ab},$$
(2)
where $`h_{ab}`$ is a linear stochastic perturbation to $`g_{ab}`$, and $`h^{ab}g^{ac}g^{bd}h_{cd}`$. In this last equation, $`\xi _n^{ab}`$ is a Gaussian stochastic tensor characterized by the correlators
$$\xi _n^{ab}(x)_c=0,\xi _n^{ab}(x)\xi _n^{cd}(y)_c=N_n^{abcd}(x,y),$$
(3)
where $`8N_n^{abcd}(x,y)\{\widehat{t}_n^{ab}(x),\widehat{t}_n^{cd}(y)\}[g]`$, with $`\widehat{t}_n^{ab}\widehat{T}_n^{ab}\widehat{T}_n^{ab}`$; here, $`_c`$ means statistical average and $`\{,\}`$ denotes an anticommutator. As we pointed out in Ref. , the noise kernel $`N_n^{abcd}(x,y)`$ is free of ultraviolet divergencies in the limit $`n4`$. Therefore, in the semiclassical Einstein-Langevin equation (2), one can perform exactly the same renormalization procedure as the one for the semiclassical Einstein equation (1), and Eq. (2) yields the physical semiclassical Einstein-Langevin equation in four spacetime dimensions.
In Ref. , we used a method based on the CTP functional technique applied to a system-environment interaction, more specifically, on the influence action formalism of Feynman and Vernon, to obtain an explicit expression for the expansion of $`\widehat{T}_n^{ab}[g+h]`$ up to first order in $`h_{cd}`$. In this way, we can write the Einstein-Langevin equation (2) in a more explicit form. This expansion involves the kernel $`H_n^{abcd}(x,y)H_{\mathrm{S}_n}^{abcd}(x,y)+H_{\mathrm{A}_n}^{abcd}(x,y)`$, with
$$H_{\mathrm{S}_n}^{abcd}(x,y)\frac{1}{4}\mathrm{Im}\mathrm{T}^{}\left(\widehat{T}_n^{ab}(x)\widehat{T}_n^{cd}(y)\right)[g],H_{\mathrm{A}_n}^{abcd}(x,y)\frac{i}{8}[\widehat{T}_n^{ab}(x),\widehat{T}_n^{cd}(y)][g],$$
(4)
where $`[,]`$ means a commutator, and we use the symbol $`\mathrm{T}^{}`$ to denote that we have to time order the field operators $`\widehat{\mathrm{\Phi }}_n`$ first and then to apply the derivative operators which appear in each term of the product $`T^{ab}(x)T^{cd}(y)`$, where $`T^{ab}`$ is the classical stress-energy tensor; see Ref. for more details. In Eq. (2), all the ultraviolet divergencies in the limit $`n4`$, which shall be removed by renormalization of the coupling constants, are in some terms containing $`\widehat{\mathrm{\Phi }}_n^2(x)`$ and in $`H_{\mathrm{S}_n}^{abcd}(x,y)`$, whereas the kernels $`N_n^{abcd}(x,y)`$ and $`H_{\mathrm{A}_n}^{abcd}(x,y)`$ are free of ultraviolet divergencies. These two last kernels can be related to the real and imaginary parts of $`\widehat{t}_n^{ab}(x)\widehat{t}_n^{cd}(y)`$ by
$$N_n^{abcd}(x,y)=\frac{1}{4}\mathrm{Re}\widehat{t}_n^{ab}(x)\widehat{t}_n^{cd}(y),H_{\mathrm{A}_n}^{abcd}(x,y)=\frac{1}{4}\mathrm{Im}\widehat{t}_n^{ab}(x)\widehat{t}_n^{cd}(y).$$
(5)
We now consider the case in which we start with a vacuum state $`|0`$ for the field quantized in spacetime $`(,g_{ab})`$. In this case, it was shown in Ref. that all the expectation values entering the Einstein-Langevin equation (2) can be written in terms of the Wightman and Feynman functions, defined as
$$G_n^+(x,y)0|\widehat{\mathrm{\Phi }}_n(x)\widehat{\mathrm{\Phi }}_n(y)|0[g],iG_{F_n}(x,y)0|\mathrm{T}\left(\widehat{\mathrm{\Phi }}_n(x)\widehat{\mathrm{\Phi }}_n(y)\right)|0[g].$$
(6)
For instance, we can write $`\widehat{\mathrm{\Phi }}_n^2(x)=iG_{F_n}(x,x)=G_n^+(x,x)`$. The expressions for the kernels, which shall be used in our calculations, can be found in Appendix A.
### A Perturbations around Minkowski spacetime
An interesting case to be analyzed in the framework of the semiclassical stochastic theory of gravity is that of a Minkowski spacetime solution of semiclassical gravity. The flat metric $`\eta _{ab}`$ in a manifold $`\mathrm{IR}^4`$ (topologically) and the usual Minkowski vacuum, denoted as $`|0`$, give the class of simplest solutions to the semiclassical Einstein equation \[note that each possible value of the parameters $`(m^2,\xi )`$ leads to a different solution\], the so called trivial solutions of semiclassical gravity . In fact, we can always choose a renormalization scheme in which the renormalized expectation value $`0|\widehat{T}_R^{ab}|0[\eta ]=0`$. Thus, Minkowski spacetime $`(\mathrm{IR}^4,\eta _{ab})`$ and the vacuum state $`|0`$ are a solution to the semiclassical Einstein equation with renormalized cosmological constant $`\mathrm{\Lambda }=0`$. The fact that the vacuum expectation value of the renormalized stress-energy operator in Minkowski spacetime should vanish was originally proposed by Wald and it may be understood as a renormalization convention . There are other possible renormalization prescriptions (see, for instance, Ref. ) in which such vacuum expectation value is proportional to $`\eta ^{ab}`$, and this would determine the value of the cosmological constant $`\mathrm{\Lambda }`$ in the semiclassical equation. Of course, all these renormalization schemes give physically equivalent results: the total effective cosmological constant, i.e., the constant of proportionality in the sum of all the terms proportional to the metric in the semiclassical Einstein and Einstein-Langevin equations, has to be zero.
Although the vacuum $`|0`$ is an eigenstate of the total four-momentum operator in Minkowski spacetime, this state is not an eigenstate of $`\widehat{T}_{ab}^R[\eta ]`$. Hence, even in these trivial solutions of semiclassical gravity, there are quantum fluctuations in the stress-energy tensor of matter and, as a result, the noise kernel does not vanish. This fact leads to consider the stochastic corrections to this class of trivial solutions of semiclassical gravity. Since, in this case, the Wightman and Feynman functions (6), their values in the two-point coincidence limit, and the products of derivatives of two of such functions appearing in expressions (A3) and (LABEL:Feynman\_expression\_3) (Appendix A) are known in dimensional regularization, we can compute the semiclassical Einstein-Langevin equation using the method outlined above.
In order to perform the calculations, it is convenient to work in a global inertial coordinate system $`\{x^\mu \}`$ and in the associated basis, in which the components of the flat metric are simply $`\eta _{\mu \nu }=\mathrm{diag}(1,1,\mathrm{},1)`$. In Minkowski spacetime, the components of the classical stress-energy tensor functional reduce to
$$T^{\mu \nu }[\eta ,\mathrm{\Phi }]=^\mu \mathrm{\Phi }^\nu \mathrm{\Phi }\frac{1}{2}\eta ^{\mu \nu }^\rho \mathrm{\Phi }_\rho \mathrm{\Phi }\frac{1}{2}\eta ^{\mu \nu }m^2\mathrm{\Phi }^2+\xi \left(\eta ^{\mu \nu }\mathrm{}^\mu ^\nu \right)\mathrm{\Phi }^2,$$
(7)
where $`\mathrm{}_\mu ^\mu `$, and the formal expression for the components of the corresponding “operator” in dimensional regularization is
$$\widehat{T}_n^{\mu \nu }[\eta ]=\frac{1}{2}\{^\mu \widehat{\mathrm{\Phi }}_n,^\nu \widehat{\mathrm{\Phi }}_n\}+𝒟^{\mu \nu }\widehat{\mathrm{\Phi }}_n^2,$$
(8)
where $`𝒟^{\mu \nu }`$ are the differential operators $`𝒟_x^{\mu \nu }(\xi 1/4)\eta ^{\mu \nu }\mathrm{}_x\xi _x^\mu _x^\nu `$ and $`\widehat{\mathrm{\Phi }}_n(x)`$ is the field operator in the Heisenberg picture in an $`n`$-dimensional Minkowski spacetime, which satisfies the Klein-Gordon equation $`(\mathrm{}m^2)\widehat{\mathrm{\Phi }}_n=0`$.
Notice, from (8), that the stress-energy tensor depends on the coupling parameter $`\xi `$ of the scalar field to the scalar curvature even in the limit of a flat spacetime. Therefore, that tensor differs in general from the canonical stress-energy tensor in flat spacetime, which corresponds to the value $`\xi =0`$. Nevertheless, it is easy to see that the $`n`$-momentum density components $`\widehat{T}_n^{0\mu }_{_{\left(\xi \right)}}[\eta ]`$ (we temporary use this notation to indicate the dependence on the parameter $`\xi `$) and $`\widehat{T}_n^{0\mu }_{_{\left(\xi =0\right)}}[\eta ]`$ differ in a space divergence and, hence, dropping surface terms, they both yield the same $`n`$-momentum operator:
$$\widehat{P}^\mu d^{n1}𝐱:\widehat{T}_n^{0\mu }_{_{\left(\xi \right)}}[\eta ]:=d^{n1}𝐱:\widehat{T}_n^{0\mu }_{_{\left(\xi =0\right)}}[\eta ]:,$$
(9)
where the integration is on a hypersurface $`x^0=\mathrm{constant}`$ ($`\widehat{P}^\mu `$ is actually independent of the value of $`x^0`$) and we use the notation for coordinates $`x^\mu (x^0,𝐱)`$, i.e., $`𝐱`$ are space coordinates on each of the hypersurfaces $`x^0=\mathrm{constant}`$. The symbol $`::`$ in Eq. (9) means normal ordering of the creation and annihilation operators on the Fock space built on the Minkowski vacuum $`|0`$ (in $`n`$ spacetime dimensions), which is an eigenstate with zero eigenvalue of the operators (9).
The Wightman and Feynman functions (6) in Minkowski spacetime are well known:
$`G_n^+(x,y)0|\widehat{\mathrm{\Phi }}_n(x)\widehat{\mathrm{\Phi }}_n(y)|0[\eta ]=i\mathrm{\Delta }_n^+(xy),`$ (10)
$`G_{F_n}(x,y)i0|\mathrm{T}\left(\widehat{\mathrm{\Phi }}_n(x)\widehat{\mathrm{\Phi }}_n(y)\right)|0[\eta ]=\mathrm{\Delta }_{F_n}(xy),`$ (11)
with
$`\mathrm{\Delta }_n^+(x)=2\pi i{\displaystyle \frac{d^nk}{(2\pi )^n}e^{ikx}\delta (k^2+m^2)\theta (k^0)},`$ (12)
$`\mathrm{\Delta }_{F_n}(x)={\displaystyle \frac{d^nk}{(2\pi )^n}\frac{e^{ikx}}{k^2+m^2iϵ}},ϵ0^+,`$ (13)
where $`k^2\eta _{\mu \nu }k^\mu k^\nu `$ and $`kx\eta _{\mu \nu }k^\mu x^\nu `$. Note that the derivatives of these functions satisfy $`_\mu ^x\mathrm{\Delta }_n^+(xy)=_\mu \mathrm{\Delta }_n^+(xy)`$ and $`_\mu ^y\mathrm{\Delta }_n^+(xy)=_\mu \mathrm{\Delta }_n^+(xy)`$, and similarly for the Feynman propagator $`\mathrm{\Delta }_{F_n}(xy)`$.
To write down the semiclassical Einstein equation (1) for this case, we need to compute the vacuum expectation value of the stress-energy operator components (8). Since, from (11), we have that $`0|\widehat{\mathrm{\Phi }}_n^2(x)|0=i\mathrm{\Delta }_{F_n}(0)=i\mathrm{\Delta }_n^+(0)`$, which is a constant (independent of $`x`$), we have simply
$$0|\widehat{T}_n^{\mu \nu }|0[\eta ]=\frac{1}{2}0|\{^\mu \widehat{\mathrm{\Phi }}_n,^\nu \widehat{\mathrm{\Phi }}_n\}|0[\eta ]=i(^\mu ^\nu \mathrm{\Delta }_{F_n})(0)=i\frac{d^nk}{(2\pi )^n}\frac{k^\mu k^\nu }{k^2+m^2iϵ}=\frac{\eta ^{\mu \nu }}{2}\left(\frac{m^2}{4\pi }\right)^{n/2}\mathrm{\Gamma }\left(\frac{n}{2}\right),$$
(14)
where the integrals in dimensional regularization have been computed in the standard way (see Appendix B) and where $`\mathrm{\Gamma }(z)`$ is the Euler’s gamma function. The semiclassical Einstein equation (1), which now reduces to
$$\frac{\mathrm{\Lambda }_B}{8\pi G_B}\eta ^{\mu \nu }=\mu ^{(n4)}0|\widehat{T}_n^{\mu \nu }|0[\eta ],$$
(15)
simply sets the value of the bare coupling constant $`\mathrm{\Lambda }_B/G_B`$. Note, from (14), that in order to have $`0|\widehat{T}_R^{ab}|0[\eta ]=0`$, the renormalized (and regularized) stress-energy tensor “operator” for a scalar field in Minkowski spacetime has to be defined as
$$\widehat{T}_R^{ab}[\eta ]=\mu ^{(n4)}\widehat{T}_n^{ab}[\eta ]\frac{\eta ^{ab}}{2}\frac{m^4}{(4\pi )^2}\left(\frac{m^2}{4\pi \mu ^2}\right)^{_{\frac{n4}{2}}}\mathrm{\Gamma }\left(\frac{n}{2}\right),$$
(16)
which corresponds to a renormalization of the cosmological constant
$$\frac{\mathrm{\Lambda }_B}{G_B}=\frac{\mathrm{\Lambda }}{G}\frac{2}{\pi }\frac{m^4}{n(n2)}\kappa _n+O(n4),$$
(17)
where
$$\kappa _n\frac{1}{(n4)}\left(\frac{e^\gamma m^2}{4\pi \mu ^2}\right)^{_{\frac{n4}{2}}}=\frac{1}{n4}+\frac{1}{2}\mathrm{ln}\left(\frac{e^\gamma m^2}{4\pi \mu ^2}\right)+O(n4),$$
(18)
being $`\gamma `$ the Euler’s constant. In the case of a massless scalar field, $`m^2=0`$, one simply has $`\mathrm{\Lambda }_B/G_B=\mathrm{\Lambda }/G`$. Introducing this renormalized coupling constant into Eq. (15), we can take the limit $`n4`$. We find again that, for $`(\mathrm{IR}^4,\eta _{ab},|0)`$ to satisfy the semiclassical Einstein equation, we must take $`\mathrm{\Lambda }=0`$.
We are now in the position to write down the Einstein-Langevin equations for the components $`h_{\mu \nu }`$ of the stochastic metric perturbation in dimensional regularization. In our case, using $`0|\widehat{\mathrm{\Phi }}_n^2(x)|0=i\mathrm{\Delta }_{F_n}(0)`$ and the explicit expression for Eq. (2) found in Ref. , we obtain that this equation reduces to
$`{\displaystyle \frac{1}{8\pi G_B}}\left[G^{\left(1\right)\mu \nu }+\mathrm{\Lambda }_B\left(h^{\mu \nu }{\displaystyle \frac{1}{2}}\eta ^{\mu \nu }h\right)\right](x){\displaystyle \frac{4}{3}}\alpha _BD^{\left(1\right)\mu \nu }(x)2\beta _BB^{\left(1\right)\mu \nu }(x)`$ (19)
$`\xi G^{\left(1\right)\mu \nu }(x)\mu ^{(n4)}i\mathrm{\Delta }_{F_n}(0)+2{\displaystyle d^ny\mu ^{(n4)}H_n^{\mu \nu \alpha \beta }(x,y)h_{\alpha \beta }(y)}=2\xi ^{\mu \nu }(x),`$ (20)
where $`\xi ^{\mu \nu }`$ are the components of a Gaussian stochastic tensor of zero average and
$$\xi ^{\mu \nu }(x)\xi ^{\alpha \beta }(y)_c=\mu ^{2(n4)}N_n^{\mu \nu \alpha \beta }(x,y),$$
(21)
and where indices are raised in $`h_{\mu \nu }`$ with the flat metric and $`hh_\rho ^\rho `$. We use a superindex $`\left(1\right)`$ to denote the components of a tensor linearized around the flat metric. In the last expressions, $`N_n^{\mu \nu \alpha \beta }(x,y)`$ and $`H_n^{\mu \nu \alpha \beta }(x,y)`$ are the components of the kernels defined above. In Eq. (20), we have made use of the explicit expression for $`G^{\left(1\right)\mu \nu }`$. This expression and those for $`D^{\left(1\right)\mu \nu }`$ and $`B^{\left(1\right)\mu \nu }`$ are given in Appendix E; the last two can also be written as
$$D^{\left(1\right)\mu \nu }(x)=\frac{1}{2}(3_x^{\mu \alpha }_x^{\nu \beta }_x^{\mu \nu }_x^{\alpha \beta })h_{\alpha \beta }(x),B^{\left(1\right)\mu \nu }(x)=2_x^{\mu \nu }_x^{\alpha \beta }h_{\alpha \beta }(x),$$
(22)
where $`_x^{\mu \nu }`$ is the differential operator $`_x^{\mu \nu }\eta ^{\mu \nu }\mathrm{}_x_x^\mu _x^\nu `$.
## III The kernels for a Minkowski background
The kernels $`N_n^{\mu \nu \alpha \beta }(x,y)`$ and $`H_n^{\mu \nu \alpha \beta }(x,y)=H_{\mathrm{S}_n}^{\mu \nu \alpha \beta }(x,y)+H_{\mathrm{A}_n}^{\mu \nu \alpha \beta }(x,y)`$ can now be computed using (5) and the expressions (A3) and (LABEL:Feynman\_expression\_3). In Ref. , we have shown that the kernel $`H_{\mathrm{A}_n}^{\mu \nu \alpha \beta }(x,y)`$ plays the role of a dissipation kernel, since it is related to the noise kernel, $`N_n^{\mu \nu \alpha \beta }(x,y)`$, by a fluctuation-dissipation relation. From the definitions (4) and the fact that the Minkowski vacuum $`|0`$ is an eigenstate of the operator $`\widehat{P}^\mu `$, given by (9), these kernels satisfy
$$d^{n1}𝐱N_n^{0\mu \alpha \beta }(x,y)=d^{n1}𝐱H_{\mathrm{A}_n}^{0\mu \alpha \beta }(x,y)=0.$$
(23)
### A The noise and dissipation kernels
Since the two kernels (5) are free of ultraviolet divergencies in the limit $`n4`$, we can deal directly with
$$M^{\mu \nu \alpha \beta }(xy)\underset{n4}{lim}\mu ^{2(n4)}0|\widehat{t}_n^{\mu \nu }(x)\widehat{t}_n^{\alpha \beta }(y)|0[\eta ].$$
(24)
The kernels $`4N^{\mu \nu \alpha \beta }(x,y)=\mathrm{Re}M^{\mu \nu \alpha \beta }(xy)`$ and $`4H_\mathrm{A}^{\mu \nu \alpha \beta }(x,y)=\mathrm{Im}M^{\mu \nu \alpha \beta }(xy)`$ are actually the components of the “physical” noise and dissipation kernels that will appear in the Einstein-Langevin equations once the renormalization procedure has been carried out. Note that, in the renormalization scheme in which $`\widehat{T}_R^{ab}[\eta ]`$ is given by (16), we can write $`M^{\mu \nu \alpha \beta }(xy)=0|\widehat{T}_R^{\mu \nu }(x)\widehat{T}_R^{\alpha \beta }(y)|0[\eta ]`$, where the limit $`n4`$ is understood. This kernel can be expressed in terms of the Wightman function in four spacetime dimensions,
$$\mathrm{\Delta }^+(x)=2\pi i\frac{d^4k}{(2\pi )^4}e^{ikx}\delta (k^2+m^2)\theta (k^0),$$
(25)
in the following way:
$`M^{\mu \nu \alpha \beta }(x)=2`$ $`[^\mu ^{(\alpha }\mathrm{\Delta }^+(x)^{\beta )}^\nu \mathrm{\Delta }^+(x)+𝒟^{\mu \nu }\left(^\alpha \mathrm{\Delta }^+(x)^\beta \mathrm{\Delta }^+(x)\right)`$ (27)
$`+𝒟^{\alpha \beta }\left(^\mu \mathrm{\Delta }^+(x)^\nu \mathrm{\Delta }^+(x)\right)+𝒟^{\mu \nu }𝒟^{\alpha \beta }\left(\mathrm{\Delta }^{+2}(x)\right)].`$
The different terms in Eq. (27) can be easily computed using the integrals
$`I(p){\displaystyle \frac{d^4k}{(2\pi )^4}\delta (k^2+m^2)\theta (k^0)\delta [(kp)^2+m^2]\theta (k^0p^0)},`$ (28)
$`I^{\mu _1\mathrm{}\mu _r}(p){\displaystyle \frac{d^4k}{(2\pi )^4}k^{\mu _1}\mathrm{}k^{\mu _r}\delta (k^2+m^2)\theta (k^0)\delta [(kp)^2+m^2]\theta (k^0p^0)},`$ (29)
with $`r=1,2,3,4`$, given in Appendix B; all of them can be expressed in terms of $`I(p)`$. We obtain expressions (C1)-(LABEL:Wightman\_3). It is convenient to separate $`I(p)`$ in its even and odd parts with respect to the variables $`p^\mu `$ as
$$I(p)=I_\mathrm{S}(p)+I_\mathrm{A}(p),$$
(30)
where $`I_\mathrm{S}(p)=I_\mathrm{S}(p)`$ and $`I_\mathrm{A}(p)=I_\mathrm{A}(p)`$. These two functions are explicitly given by
$`I_\mathrm{S}(p)={\displaystyle \frac{1}{8(2\pi )^3}}\theta (p^24m^2)\sqrt{1+4{\displaystyle \frac{m^2}{p^2}}},`$ (31)
$`I_\mathrm{A}(p)={\displaystyle \frac{1}{8(2\pi )^3}}\mathrm{sign}p^0\theta (p^24m^2)\sqrt{1+4{\displaystyle \frac{m^2}{p^2}}}.`$ (32)
Using the results of Appendix B, we obtain expressions (LABEL:Wightman\_4)-(LABEL:Wightman\_6) and, after some calculations, we find
$`M^{\mu \nu \alpha \beta }(x)=`$ $`{\displaystyle \frac{\pi ^2}{45}}(3_x^{\mu (\alpha }_x^{\beta )\nu }_x^{\mu \nu }_x^{\alpha \beta }){\displaystyle \frac{d^4p}{(2\pi )^4}e^{ipx}\left(1+4\frac{m^2}{p^2}\right)^2I(p)}`$ (34)
$`+{\displaystyle \frac{8\pi ^2}{9}}_x^{\mu \nu }_x^{\alpha \beta }{\displaystyle \frac{d^4p}{(2\pi )^4}e^{ipx}\left(3\mathrm{\Delta }\xi +\frac{m^2}{p^2}\right)^2I(p)},`$
where $`\mathrm{\Delta }\xi \xi 1/6`$. The real and imaginary parts of the last expression, which yield the noise and dissipation kernels, are easily recognized as the terms containing $`I_\mathrm{S}(p)`$ and $`I_\mathrm{A}(p)`$, respectively. To write them explicitly, it is useful to introduce the new kernels
$`N_\mathrm{A}(x;m^2){\displaystyle \frac{1}{1920\pi }}{\displaystyle \frac{d^4p}{(2\pi )^4}e^{ipx}\theta (p^24m^2)\sqrt{1+4\frac{m^2}{p^2}}\left(1+4\frac{m^2}{p^2}\right)^2},`$ (35)
$`N_\mathrm{B}(x;m^2,\mathrm{\Delta }\xi ){\displaystyle \frac{1}{288\pi }}{\displaystyle \frac{d^4p}{(2\pi )^4}e^{ipx}\theta (p^24m^2)\sqrt{1+4\frac{m^2}{p^2}}\left(3\mathrm{\Delta }\xi +\frac{m^2}{p^2}\right)^2},`$ (36)
$`D_\mathrm{A}(x;m^2){\displaystyle \frac{i}{1920\pi }}{\displaystyle \frac{d^4p}{(2\pi )^4}e^{ipx}\mathrm{sign}p^0\theta (p^24m^2)\sqrt{1+4\frac{m^2}{p^2}}\left(1+4\frac{m^2}{p^2}\right)^2},`$ (37)
$`D_\mathrm{B}(x;m^2,\mathrm{\Delta }\xi ){\displaystyle \frac{i}{288\pi }}{\displaystyle \frac{d^4p}{(2\pi )^4}e^{ipx}\mathrm{sign}p^0\theta (p^24m^2)\sqrt{1+4\frac{m^2}{p^2}}\left(3\mathrm{\Delta }\xi +\frac{m^2}{p^2}\right)^2},`$ (38)
and we finally get
$`N^{\mu \nu \alpha \beta }(x,y)={\displaystyle \frac{1}{6}}(3_x^{\mu (\alpha }_x^{\beta )\nu }_x^{\mu \nu }_x^{\alpha \beta })N_\mathrm{A}(xy;m^2)+_x^{\mu \nu }_x^{\alpha \beta }N_\mathrm{B}(xy;m^2,\mathrm{\Delta }\xi ),`$ (39)
$`H_\mathrm{A}^{\mu \nu \alpha \beta }(x,y)={\displaystyle \frac{1}{6}}(3_x^{\mu (\alpha }_x^{\beta )\nu }_x^{\mu \nu }_x^{\alpha \beta })D_\mathrm{A}(xy;m^2)+_x^{\mu \nu }_x^{\alpha \beta }D_\mathrm{B}(xy;m^2,\mathrm{\Delta }\xi ).`$ (40)
Notice that the noise and dissipation kernels defined in (38) are actually real because, for the noise kernels, only the $`\mathrm{cos}px`$ terms of the exponentials $`e^{ipx}`$ contribute to the integrals, and, for the dissipation kernels, the only contribution of such exponentials comes from the $`i\mathrm{sin}px`$ terms.
We can now evaluate the contribution of the dissipation kernel components $`H_\mathrm{A}^{\mu \nu \alpha \beta }(x,y)`$ to the Einstein-Langevin equations (20) \[after taking the limit $`n4`$\]. From (40), integrating by parts, and using (22) and the fact that, in four spacetime dimensions, $`D^{\left(1\right)\mu \nu }(x)=(3/2)A^{\left(1\right)\mu \nu }(x)`$ \[the tensor $`A^{ab}`$ is obtained from the derivative with respect to the metric of an action term corresponding to the Lagrangian density $`C_{abcd}C^{abcd}`$, where $`C_{abcd}`$ is the Weyl tensor, see Ref. for details\], it is easy to see that
$$2d^4yH_\mathrm{A}^{\mu \nu \alpha \beta }(x,y)h_{\alpha \beta }(y)=d^4y\left[D_\mathrm{A}(xy;m^2)A^{\left(1\right)\mu \nu }(y)+D_\mathrm{B}(xy;m^2,\mathrm{\Delta }\xi )B^{\left(1\right)\mu \nu }(y)\right].$$
(41)
These non-local terms in the semiclassical Einstein-Langevin equations can actually be identified as being part of $`\widehat{T}_R^{\mu \nu }[\eta +h]`$.
### B The kernel $`H_{\mathrm{S}_n}^{\mu \nu \alpha \beta }(x,y)`$
The evaluation of the kernel components $`H_{\mathrm{S}_n}^{\mu \nu \alpha \beta }(x,y)`$ is a much more cumbersome task. Since these quantities contain divergencies in the limit $`n4`$, we shall compute them using dimensional regularization. Using Eq. (LABEL:Feynman\_expression\_3), these components can be written in terms of the Feynman propagator (13) as
$$\mu ^{(n4)}H_{\mathrm{S}_n}^{\mu \nu \alpha \beta }(x,y)=\frac{1}{4}\mathrm{Im}K^{\mu \nu \alpha \beta }(xy),$$
(42)
where
$`K^{\mu \nu \alpha \beta }(x)\mu ^{(n4)}\{2^\mu ^{(\alpha }\mathrm{\Delta }_{F_n}(x)^{\beta )}^\nu \mathrm{\Delta }_{F_n}(x)+2𝒟^{\mu \nu }\left(^\alpha \mathrm{\Delta }_{F_n}(x)^\beta \mathrm{\Delta }_{F_n}(x)\right)`$ (43)
$`+\mathrm{\hspace{0.17em}2}𝒟^{\alpha \beta }\left(^\mu \mathrm{\Delta }_{F_n}(x)^\nu \mathrm{\Delta }_{F_n}(x)\right)+2𝒟^{\mu \nu }𝒟^{\alpha \beta }\left(\mathrm{\Delta }_{F_n}^2(x)\right)+[\eta ^{\mu \nu }^{(\alpha }\mathrm{\Delta }_{F_n}(x)^{\beta )}+\eta ^{\alpha \beta }^{(\mu }\mathrm{\Delta }_{F_n}(x)^{\nu )}`$ (44)
$`+\mathrm{\Delta }_{F_n}(0)(\eta ^{\mu \nu }𝒟^{\alpha \beta }+\eta ^{\alpha \beta }𝒟^{\mu \nu })+{\displaystyle \frac{1}{4}}\eta ^{\mu \nu }\eta ^{\alpha \beta }(\mathrm{\Delta }_{F_n}(x)\mathrm{}m^2\mathrm{\Delta }_{F_n}(0))]\delta ^n(x)\}.`$ (45)
Let us define the integrals
$`J_n(p)\mu ^{(n4)}{\displaystyle \frac{d^nk}{(2\pi )^n}\frac{1}{(k^2+m^2iϵ)[(kp)^2+m^2iϵ]}},`$ (46)
$`J_n^{\mu _1\mathrm{}\mu _r}(p)\mu ^{(n4)}{\displaystyle \frac{d^nk}{(2\pi )^n}\frac{k^{\mu _1}\mathrm{}k^{\mu _r}}{(k^2+m^2iϵ)[(kp)^2+m^2iϵ]}},`$ (47)
with $`r=1,2,3,4`$, and
$`I_{0_n}\mu ^{(n4)}{\displaystyle \frac{d^nk}{(2\pi )^n}\frac{1}{(k^2+m^2iϵ)}},`$ (48)
$`I_{0_n}^{\mu _1\mathrm{}\mu _r}\mu ^{(n4)}{\displaystyle \frac{d^nk}{(2\pi )^n}\frac{k^{\mu _1}\mathrm{}k^{\mu _r}}{(k^2+m^2iϵ)}},`$ (49)
with $`r=1,2`$, where a limit $`ϵ0^+`$ is understood in all these expressions. Then, the different terms in Eq. (45) can be computed using Eqs. (D1)-(D7). The results for the expansions of the integrals (47) and (49) around $`n=4`$ are given in Appendix B. In fact, $`I_{0_n}^\mu =0`$ and the remaining integrals can be written in terms of $`I_{0_n}`$ and $`J_n(p)`$ given in Eqs. (B1) and (B4). Using the results of Appendix B, we obtain Eqs. (LABEL:Feynman\_7) and (D17) and, from Eqs. (D5)-(D7), we get
$`\mu ^{(n4)}\left[\eta ^{\mu \nu }^{(\alpha }\mathrm{\Delta }_{F_n}(x)^{\beta )}+\eta ^{\alpha \beta }^{(\mu }\mathrm{\Delta }_{F_n}(x)^{\nu )}\right]\delta ^n(x)=2\eta ^{\mu \nu }\eta ^{\alpha \beta }{\displaystyle \frac{m^2}{n}}I_{0_n}\delta ^n(x),`$ (50)
$`\mu ^{(n4)}\left(\mathrm{\Delta }_{F_n}(x)\mathrm{}m^2\mathrm{\Delta }_{F_n}(0)\right)\delta ^n(x)=I_{0_n}\mathrm{}\delta ^n(x).`$ (51)
We are now in the position to work out the explicit expression for $`K^{\mu \nu \alpha \beta }(x)`$, defined in (45). We use Eqs. (51), the results (D1), (D5), (LABEL:Feynman\_7) and (D17), the identities $`\delta ^n(x)=(2\pi )^nd^npe^{ipx}`$, $`_x^{\mu \nu }d^npe^{ipx}f(p)=d^npe^{ipx}f(p)p^2P^{\mu \nu }`$ and $`_x^\mu _x^\nu d^npe^{ipx}f(p)=d^npe^{ipx}f(p)p^\mu p^\nu `$, where $`f(p)`$ is an arbitrary function of $`p^\mu `$ and $`P^{\mu \nu }`$ is the projector orthogonal to $`p^\mu `$ defined as $`p^2P^{\mu \nu }\eta ^{\mu \nu }p^2p^\mu p^\nu `$, and the expansions in (B1) and (B4) for $`J_n(p)`$ and $`I_{0_n}`$. After a rather long but straightforward calculation, we get, expanding around $`n=4`$,
$`K^{\mu \nu \alpha \beta }(x)={\displaystyle \frac{i}{(4\pi )^2}}\{\kappa _n[{\displaystyle \frac{1}{90}}(3_x^{\mu (\alpha }_x^{\beta )\nu }_x^{\mu \nu }_x^{\alpha \beta })\delta ^n(x)+4\mathrm{\Delta }\xi ^2_x^{\mu \nu }_x^{\alpha \beta }\delta ^n(x)`$ (52)
$`+{\displaystyle \frac{2}{3}}{\displaystyle \frac{m^2}{(n2)}}(\eta ^{\mu \nu }\eta ^{\alpha \beta }\mathrm{}_x\eta ^{\mu (\alpha }\eta ^{\beta )\nu }\mathrm{}_x+\eta ^{\mu (\alpha }^{\beta )}_x^\nu _x+\eta ^{\nu (\alpha }^{\beta )}_x^\mu _x\eta ^{\mu \nu }^\alpha _x^\beta _x\eta ^{\alpha \beta }^\mu _x^\nu _x)\delta ^n(x)`$ (53)
$`+{\displaystyle \frac{4m^4}{n(n2)}}(2\eta ^{\mu (\alpha }\eta ^{\beta )\nu }\eta ^{\mu \nu }\eta ^{\alpha \beta })\delta ^n(x)]+{\displaystyle \frac{1}{180}}(3_x^{\mu (\alpha }_x^{\beta )\nu }_x^{\mu \nu }_x^{\alpha \beta })`$ (54)
$`\times {\displaystyle }{\displaystyle \frac{d^np}{(2\pi )^n}}e^{ipx}(1+4{\displaystyle \frac{m^2}{p^2}})^2\varphi (p^2)+{\displaystyle \frac{2}{9}}_x^{\mu \nu }_x^{\alpha \beta }{\displaystyle }{\displaystyle \frac{d^np}{(2\pi )^n}}e^{ipx}(3\mathrm{\Delta }\xi +{\displaystyle \frac{m^2}{p^2}})^2\varphi (p^2)`$ (55)
$`\left[{\displaystyle \frac{4}{675}}(3_x^{\mu (\alpha }_x^{\beta )\nu }_x^{\mu \nu }_x^{\alpha \beta })+{\displaystyle \frac{1}{270}}(60\xi 11)_x^{\mu \nu }_x^{\alpha \beta }\right]\delta ^n(x)`$ (56)
$`m^2[{\displaystyle \frac{2}{135}}(3_x^{\mu (\alpha }_x^{\beta )\nu }_x^{\mu \nu }_x^{\alpha \beta })+{\displaystyle \frac{1}{27}}_x^{\mu \nu }_x^{\alpha \beta }]\mathrm{\Delta }_n(x)\}+O(n4),`$ (57)
where $`\kappa _n`$ and $`\varphi (p^2)`$ have been defined in (18) and (B16), and $`\mathrm{\Delta }_n(x)`$ is given by
$$\mathrm{\Delta }_n(x)\frac{d^np}{(2\pi )^n}e^{ipx}\frac{1}{p^2}.$$
(58)
The imaginary part of (57) \[which, using (42), gives the kernel components $`\mu ^{(n4)}H_{\mathrm{S}_n}^{\mu \nu \alpha \beta }(x,y)`$\] can be easily obtained multiplying this expression by $`i`$ and retaining only the real part, $`\phi (p^2)`$, of the function $`\varphi (p^2)`$. Making use of this result, it is easy to compute the contribution of these kernel components to the Einstein-Langevin equations (20). Integrating by parts, using Eqs. (E1)-(E5) and Eq. (22), and taking into account that, from Eqs. (14) and (15),
$$\frac{\mathrm{\Lambda }_B}{8\pi G_B}=\frac{1}{4\pi ^2}\frac{m^4}{n(n2)}\kappa _n+O(n4),$$
(59)
we finally find
$`2{\displaystyle }d^ny\mu ^{(n4)}H_{\mathrm{S}_n}^{\mu \nu \alpha \beta }(x,y)h_{\alpha \beta }(y)={\displaystyle \frac{\mathrm{\Lambda }_B}{8\pi G_B}}[h^{\mu \nu }{\displaystyle \frac{1}{2}}\eta ^{\mu \nu }h](x)+{\displaystyle \frac{\kappa _n}{(4\pi )^2}}[{\displaystyle \frac{2}{3}}{\displaystyle \frac{m^2}{(n2)}}G^{\left(1\right)\mu \nu }`$ (60)
$`+{\displaystyle \frac{1}{90}}D^{\left(1\right)\mu \nu }+\mathrm{\Delta }\xi ^2B^{\left(1\right)\mu \nu }](x)+{\displaystyle \frac{1}{2880\pi ^2}}\{{\displaystyle \frac{16}{15}}D^{\left(1\right)\mu \nu }(x)+({\displaystyle \frac{1}{6}}10\mathrm{\Delta }\xi )B^{\left(1\right)\mu \nu }(x)`$ (61)
$`+{\displaystyle d^ny\frac{d^np}{(2\pi )^n}e^{ip(xy)}\phi (p^2)\left[\left(1+4\frac{m^2}{p^2}\right)^2D^{\left(1\right)\mu \nu }(y)+10\left(3\mathrm{\Delta }\xi +\frac{m^2}{p^2}\right)^2B^{\left(1\right)\mu \nu }(y)\right]}`$ (62)
$`{\displaystyle \frac{m^2}{3}}{\displaystyle }d^ny\mathrm{\Delta }_n(xy)(8D^{\left(1\right)\mu \nu }(y)+5B^{\left(1\right)\mu \nu }(y))\}+O(n4).`$ (63)
### C Fluctuation-dissipation relation
From expressions (40) and (38) it is easy to check that there exists a relation between the noise and dissipation kernels in the form of a fluctuation-dissipation relation which was derived in Ref. in a more general context. Introducing the Fourier transforms in the time coordinates of these kernels as
$$N^{\mu \nu \alpha \beta }(x,y)=_{\mathrm{}}^{\mathrm{}}\frac{dp^0}{2\pi }e^{ip^0(x^0y^0)}\overline{N}^{\mu \nu \alpha \beta }(p^0;𝐱,𝐲),$$
(64)
and similarly for the dissipation kernel, this relation can be written as
$$\overline{H}_\mathrm{A}^{\mu \nu \alpha \beta }(p^0;𝐱,𝐲)=i\mathrm{sign}p^0\overline{N}^{\mu \nu \alpha \beta }(p^0;𝐱,𝐲),$$
(65)
or, equivalently, as
$$H_\mathrm{A}^{\mu \nu \alpha \beta }(x^0,𝐱;y^0,𝐲)=\frac{1}{\pi }_{\mathrm{}}^{\mathrm{}}𝑑z^0\mathrm{P}\left(\frac{1}{x^0z^0}\right)N^{\mu \nu \alpha \beta }(z^0,𝐱;y^0,𝐲),$$
(66)
where $`\mathrm{P}(1/x^0)`$ denotes the principal value distribution.
From (23), taking the limit $`n4`$, we see that the noise and dissipation kernels must satisfy
$$d^3𝐱N^{0\mu \alpha \beta }(x,y)=d^3𝐱H_\mathrm{A}^{0\mu \alpha \beta }(x,y)=0.$$
(67)
In order to check the last relations, it is useful to write the $`_x^{\mu \nu }`$ derivatives in expressions (40) using $`_x^{\mu \nu }d^4pe^{ip(xy)}f(p)=d^4pe^{ip(xy)}f(p)p^2P^{\mu \nu }`$, where $`f(p)`$ is any function of $`p^\mu `$ and $`P^{\mu \nu }`$ is the projector orthogonal to $`p^\mu `$ defined above. The identities (67) follow by noting that $`p^2P^{00}=p^ip_i`$ and $`p^2P^{0i}=p^0p^i`$, where we use the index $`i=1,2,3`$ to denote the space components, and that $`d^3𝐱\mathrm{exp}(ip_ix^i)=(2\pi )^3_{i=1}^3\delta (p^i)`$. It is also easy to check that the noise kernel satisfies $`_\mu ^^xN^{\mu \nu \alpha \beta }(x,y)=0`$ and, hence, the stochastic source in the Einstein-Langevin equations will be conserved up to first order in perturbation theory.
## IV The semiclassical Einstein-Langevin equations
The results of the previous section are now ready to be introduced into the Einstein-Langevin equations (20). In fact, substituting expression (63) in such equations, and using Eqs. (D5) and (B1) for the $`\mu ^{(n4)}\mathrm{\Delta }_{F_n}(0)`$ term, we get
$`{\displaystyle \frac{1}{8\pi G_B}}G^{\left(1\right)\mu \nu }(x){\displaystyle \frac{4}{3}}\alpha _BD^{\left(1\right)\mu \nu }(x)2\beta _BB^{\left(1\right)\mu \nu }(x)+{\displaystyle \frac{\kappa _n}{(4\pi )^2}}[4\mathrm{\Delta }\xi {\displaystyle \frac{m^2}{(n2)}}G^{\left(1\right)\mu \nu }+{\displaystyle \frac{1}{90}}D^{\left(1\right)\mu \nu }`$ (68)
$`+\mathrm{\Delta }\xi ^2B^{\left(1\right)\mu \nu }](x)+{\displaystyle \frac{1}{2880\pi ^2}}\{{\displaystyle \frac{16}{15}}D^{\left(1\right)\mu \nu }(x)+({\displaystyle \frac{1}{6}}10\mathrm{\Delta }\xi )B^{\left(1\right)\mu \nu }(x)`$ (69)
$`+{\displaystyle d^ny\frac{d^np}{(2\pi )^n}e^{ip(xy)}\phi (p^2)\left[\left(1+4\frac{m^2}{p^2}\right)^2D^{\left(1\right)\mu \nu }(y)+10\left(3\mathrm{\Delta }\xi +\frac{m^2}{p^2}\right)^2B^{\left(1\right)\mu \nu }(y)\right]}`$ (70)
$`{\displaystyle \frac{m^2}{3}}{\displaystyle }d^ny\mathrm{\Delta }_n(xy)(8D^{\left(1\right)\mu \nu }+5B^{\left(1\right)\mu \nu })(y)\}+2{\displaystyle }d^ny\mu ^{(n4)}H_{\mathrm{A}_n}^{\mu \nu \alpha \beta }(x,y)h_{\alpha \beta }(y)+O(n4)`$ (71)
$`=2\xi ^{\mu \nu }(x).`$ (72)
Notice that the terms containing the bare cosmological constant have canceled. These equations can now be renormalized, that is, we can now write the bare coupling constants as renormalized coupling constants plus some suitably chosen counterterms and take the limit $`n4`$. In order to carry out such a procedure, it is convenient to distinguish between massive and massless scalar fields. We shall evaluate these two cases in different subsections.
### A Massive field ($`m0`$)
In the case of a scalar field with mass $`m0`$, we can use, as we have done in Eq. (17) for the cosmological constant, a renormalization scheme consisting on the subtraction of terms proportional to $`\kappa _n`$. More specifically, we may introduce the renormalized coupling constants $`1/G`$, $`\alpha `$ and $`\beta `$ as
$`{\displaystyle \frac{1}{G_B}}={\displaystyle \frac{1}{G}}+{\displaystyle \frac{2}{\pi }}\mathrm{\Delta }\xi {\displaystyle \frac{m^2}{(n2)}}\kappa _n+O(n4),`$ (73)
$`\alpha _B=\alpha +{\displaystyle \frac{1}{(4\pi )^2}}{\displaystyle \frac{1}{120}}\kappa _n+O(n4),`$ (74)
$`\beta _B=\beta +{\displaystyle \frac{\mathrm{\Delta }\xi ^2}{32\pi ^2}}\kappa _n+O(n4).`$ (75)
Note that for conformal coupling, $`\mathrm{\Delta }\xi =0`$, one has $`1/G_B=1/G`$ and $`\beta _B=\beta `$, that is, only the coupling constant $`\alpha `$ and the cosmological constant need renormalization. Substituting the above expressions into Eq. (72), we can now take the limit $`n4`$, using Eqs. (58), (41) and the fact that, for $`n=4`$, $`D^{\left(1\right)\mu \nu }(x)=(3/2)A^{\left(1\right)\mu \nu }(x)`$. We obtain the semiclassical Einstein-Langevin equations for the physical stochastic perturbations $`h_{\mu \nu }`$ in the four-dimensional manifold $`\mathrm{IR}^4`$. Introducing the two new kernels
$`H_\mathrm{A}(x;m^2){\displaystyle \frac{1}{1920\pi ^2}}{\displaystyle }{\displaystyle \frac{d^4p}{(2\pi )^4}}e^{ipx}\{(1+4{\displaystyle \frac{m^2}{p^2}})^2[i\pi \mathrm{sign}p^0\theta (p^24m^2)\sqrt{1+4{\displaystyle \frac{m^2}{p^2}}}`$ (76)
$`+\phi (p^2)]{\displaystyle \frac{8}{3}}{\displaystyle \frac{m^2}{p^2}}\},`$ (77)
$`H_\mathrm{B}(x;m^2,\mathrm{\Delta }\xi ){\displaystyle \frac{1}{288\pi ^2}}{\displaystyle }{\displaystyle \frac{d^4p}{(2\pi )^4}}e^{ipx}\{(3\mathrm{\Delta }\xi +{\displaystyle \frac{m^2}{p^2}})^2[i\pi \mathrm{sign}p^0\theta (p^24m^2)\sqrt{1+4{\displaystyle \frac{m^2}{p^2}}}`$ (78)
$`+\phi (p^2)]{\displaystyle \frac{1}{6}}{\displaystyle \frac{m^2}{p^2}}\},`$ (79)
where $`\phi (p^2)`$ is given by the restriction to $`n=4`$ of expression (B18), these Einstein-Langevin equations can be written as
$`{\displaystyle \frac{1}{8\pi G}}G^{\left(1\right)\mu \nu }(x)2\left(\alpha A^{\left(1\right)\mu \nu }(x)+\beta B^{\left(1\right)\mu \nu }(x)\right)+{\displaystyle \frac{1}{2880\pi ^2}}\left[{\displaystyle \frac{8}{5}}A^{\left(1\right)\mu \nu }(x)+({\displaystyle \frac{1}{6}}10\mathrm{\Delta }\xi )B^{\left(1\right)\mu \nu }(x)\right]`$ (80)
$`+{\displaystyle d^4y\left[H_\mathrm{A}(xy;m^2)A^{\left(1\right)\mu \nu }(y)+H_\mathrm{B}(xy;m^2,\mathrm{\Delta }\xi )B^{\left(1\right)\mu \nu }(y)\right]}=2\xi ^{\mu \nu }(x),`$ (81)
where $`\xi ^{\mu \nu }`$ are the components of a Gaussian stochastic tensor of vanishing mean value and two-point correlation function $`\xi ^{\mu \nu }(x)\xi ^{\alpha \beta }(y)_c=N^{\mu \nu \alpha \beta }(x,y)`$, given in (40). Note that the two kernels defined in (79) are real and can be split into an even part and an odd part with respect to the variables $`x^\mu `$, with the odd terms being the dissipation kernels $`D_\mathrm{A}(x;m^2)`$ and $`D_\mathrm{B}(x;m^2,\mathrm{\Delta }\xi )`$ defined in (38). In spite of appearances, one can show that the Fourier transforms of the even parts of these kernels are finite in the limit $`p^20`$ and, hence, the kernels $`H_\mathrm{A}`$ and $`H_\mathrm{B}`$ are well defined distributions.
We should mention that, in a previous work in Ref. , the same Einstein-Langevin equations were calculated using rather different methods. The way in which the result is written makes difficult a direct comparison with our equations (81). For instance, it is not obvious that in those previously derived equations there is some analog of the dissipation kernels related to the noise kernels by a fluctuation-dissipation relation of the form (65) or (66).
### B Massless field ($`m=0`$)
In this subsection, we consider the limit $`m0`$ of equations (72). The renormalization scheme used in the previous subsection becomes singular in the massless limit because the expressions (75) for $`\alpha _B`$ and $`\beta _B`$ diverge when $`m0`$. Therefore, a different renormalization scheme is needed in this case. First, note that we may separate $`\kappa _n`$ in (18) as $`\kappa _n=\stackrel{~}{\kappa }_n+\frac{1}{2}\mathrm{ln}(m^2/\mu ^2)+O(n4)`$, where
$$\stackrel{~}{\kappa }_n\frac{1}{(n4)}\left(\frac{e^\gamma }{4\pi }\right)^{_{\frac{n4}{2}}}=\frac{1}{n4}+\frac{1}{2}\mathrm{ln}\left(\frac{e^\gamma }{4\pi }\right)+O(n4),$$
(82)
and that \[see Eq. (B18)\]
$$\underset{m^20}{lim}\left[\phi (p^2)+\mathrm{ln}(m^2/\mu ^2)\right]=2+\mathrm{ln}\left|\frac{p^2}{\mu ^2}\right|.$$
(83)
Hence, in the massless limit, equations (72) reduce to
$`{\displaystyle \frac{1}{8\pi G_B}}G^{\left(1\right)\mu \nu }(x){\displaystyle \frac{4}{3}}\alpha _BD^{\left(1\right)\mu \nu }(x)2\beta _BB^{\left(1\right)\mu \nu }(x)+{\displaystyle \frac{1}{(4\pi )^2}}(\stackrel{~}{\kappa }_n1)[{\displaystyle \frac{1}{90}}D^{\left(1\right)\mu \nu }+\mathrm{\Delta }\xi ^2B^{\left(1\right)\mu \nu }](x)`$ (84)
$`+{\displaystyle \frac{1}{2880\pi ^2}}\{{\displaystyle \frac{16}{15}}D^{\left(1\right)\mu \nu }(x)+({\displaystyle \frac{1}{6}}10\mathrm{\Delta }\xi )B^{\left(1\right)\mu \nu }(x)+{\displaystyle }d^ny{\displaystyle }{\displaystyle \frac{d^np}{(2\pi )^n}}e^{ip(xy)}\mathrm{ln}|{\displaystyle \frac{p^2}{\mu ^2}}|[D^{\left(1\right)\mu \nu }(y)`$ (85)
$`+\mathrm{\hspace{0.17em}90}\mathrm{\Delta }\xi ^2B^{\left(1\right)\mu \nu }(y)]\}+lim_{m^20}2{\displaystyle }d^ny\mu ^{(n4)}H_{\mathrm{A}_n}^{\mu \nu \alpha \beta }(x,y)h_{\alpha \beta }(y)+O(n4)=2\xi ^{\mu \nu }(x).`$ (86)
These equations can be renormalized by introducing the renormalized coupling constants $`1/G`$, $`\alpha `$ and $`\beta `$ as
$$\frac{1}{G_B}=\frac{1}{G},\alpha _B=\alpha +\frac{1}{(4\pi )^2}\frac{1}{120}(\stackrel{~}{\kappa }_n1)+O(n4),\beta _B=\beta +\frac{\mathrm{\Delta }\xi ^2}{32\pi ^2}(\stackrel{~}{\kappa }_n1)+O(n4).$$
(87)
Thus, in the massless limit, the Newtonian gravitational constant is not renormalized and, in the conformal coupling case, $`\mathrm{\Delta }\xi =0`$, we have again that $`\beta _B=\beta `$. Introducing the last expressions into Eq. (86), we can take the limit $`n4`$. Note that, by making $`m=0`$ in (38), the noise and dissipation kernels can be written as
$`N_\mathrm{A}(x;m^2=0)=N(x),N_\mathrm{B}(x;m^2=0,\mathrm{\Delta }\xi )=60\mathrm{\Delta }\xi ^2N(x),`$ (88)
$`D_\mathrm{A}(x;m^2=0)=D(x),D_\mathrm{B}(x;m^2=0,\mathrm{\Delta }\xi )=60\mathrm{\Delta }\xi ^2D(x),`$ (89)
where
$$N(x)\frac{1}{1920\pi }\frac{d^4p}{(2\pi )^4}e^{ipx}\theta (p^2),D(x)\frac{i}{1920\pi }\frac{d^4p}{(2\pi )^4}e^{ipx}\mathrm{sign}p^0\theta (p^2).$$
(90)
It is now convenient to introduce the new kernel
$`H(x;\mu ^2)`$ $``$ $`{\displaystyle \frac{1}{1920\pi ^2}}{\displaystyle \frac{d^4p}{(2\pi )^4}e^{ipx}\left[\mathrm{ln}\left|\frac{p^2}{\mu ^2}\right|i\pi \mathrm{sign}p^0\theta (p^2)\right]}`$ (91)
$`=`$ $`{\displaystyle \frac{1}{1920\pi ^2}}\underset{ϵ0^+}{lim}{\displaystyle \frac{d^4p}{(2\pi )^4}e^{ipx}\mathrm{ln}\left(\frac{(p^0+iϵ)^2+p^ip_i}{\mu ^2}\right)}.`$ (92)
Again, this kernel is real and can be written as the sum of an even part and an odd part in the variables $`x^\mu `$, where the odd part is the dissipation kernel $`D(x)`$. The Fourier transforms (90) and (92) can actually be computed and, thus, in this case, we have explicit expressions for the kernels in position space. For $`N(x)`$ and $`D(x)`$, we get (see, for instance, Ref. )
$$N(x)=\frac{1}{1920\pi }\left[\frac{1}{\pi ^3}𝒫f\left(\frac{1}{(x^2)^2}\right)+\delta ^4(x)\right],D(x)=\frac{1}{1920\pi ^3}\mathrm{sign}x^0\frac{d}{d(x^2)}\delta (x^2),$$
(93)
where $`𝒫f`$ denotes a distribution generated by the Hadamard finite part of a divergent integral (see Refs. for the definition of these distributions). The expression for the kernel $`H(x;\mu ^2)`$ can be found in Refs. and it is given by
$`H(x;\mu ^2)`$ $`=`$ $`{\displaystyle \frac{1}{960\pi ^2}}\left\{𝒫f\left({\displaystyle \frac{1}{\pi }}\theta (x^0){\displaystyle \frac{d}{d(x^2)}}\delta (x^2)\right)+\left(1\gamma \mathrm{ln}\mu \right)\delta ^4(x)\right\}`$ (94)
$`=`$ $`{\displaystyle \frac{1}{960\pi ^2}}\underset{\lambda 0^+}{lim}\left\{{\displaystyle \frac{1}{\pi }}\theta (x^0)\theta (|𝐱|\lambda ){\displaystyle \frac{d}{d(x^2)}}\delta (x^2)+\left[1\gamma \mathrm{ln}(\mu \lambda )\right]\delta ^4(x)\right\}.`$ (95)
See Ref. for the details on how this last distribution acts on a test function. Finally, the semiclassical Einstein-Langevin equations for the physical stochastic perturbations $`h_{\mu \nu }`$ in the massless case are
$`{\displaystyle \frac{1}{8\pi G}}G^{\left(1\right)\mu \nu }(x)2\left(\alpha A^{\left(1\right)\mu \nu }(x)+\beta B^{\left(1\right)\mu \nu }(x)\right)+{\displaystyle \frac{1}{2880\pi ^2}}\left[{\displaystyle \frac{8}{5}}A^{\left(1\right)\mu \nu }(x)+({\displaystyle \frac{1}{6}}10\mathrm{\Delta }\xi )B^{\left(1\right)\mu \nu }(x)\right]`$ (96)
$`+{\displaystyle d^4yH(xy;\mu ^2)\left[A^{\left(1\right)\mu \nu }(y)+60\mathrm{\Delta }\xi ^2B^{\left(1\right)\mu \nu }(y)\right]}=2\xi ^{\mu \nu }(x),`$ (97)
where the Gaussian stochastic source components $`\xi ^{\mu \nu }`$ have zero mean value and
$$\xi ^{\mu \nu }(x)\xi ^{\alpha \beta }(y)_c=\underset{m0}{lim}N^{\mu \nu \alpha \beta }(x,y)=[\frac{1}{6}(3_x^{\mu (\alpha }_x^{\beta )\nu }_x^{\mu \nu }_x^{\alpha \beta })+60\mathrm{\Delta }\xi ^2_x^{\mu \nu }_x^{\alpha \beta }]N(xy).$$
(98)
It is interesting to consider the conformally coupled scalar field, i.e., the case $`\mathrm{\Delta }\xi =0`$, of particular interest because of its similarities with the electromagnetic field. It was shown in Refs. that, for this field, the stochastic source tensor must be “traceless” (up to first order in perturbation theory around semiclassical gravity), in the sense that the stochastic variable $`\xi _\mu ^\mu \eta _{\mu \nu }\xi ^{\mu \nu }`$ behaves deterministically as a vanishing scalar field. This can be easily checked by noticing, from Eq. (98), that, when $`\mathrm{\Delta }\xi =0`$, one has $`\xi _\mu ^\mu (x)\xi ^{\alpha \beta }(y)_c=0`$, since $`_\mu ^\mu =3\mathrm{}`$ and $`^{\mu \alpha }_\mu ^\beta =\mathrm{}^{\alpha \beta }`$. The semiclassical Einstein-Langevin equations for this particular case \[and generalized to a spatially flat Robertson-Walker (RW) background\] were first obtained in Ref. (in this reference, the coupling constant $`\beta `$ was set to zero). In order to compare with this previous result, it is worth noticing that the description of the stochastic source in terms of a symmetric and “traceless” tensor, with nine independent components $`\xi ^{\mu \nu }`$, is equivalent to a description in terms of a Gaussian stochastic tensor with the same symmetry properties as the Weyl tensor, with components $`\xi _c^{\mu \nu \alpha \beta }`$, defined as $`\xi ^{\mu \nu }=2_\alpha _\beta \xi _c^{\mu \alpha \nu \beta }`$; this tensor is used in Ref. . The symmetry properties of the $`\xi _c^{\mu \nu \alpha \beta }`$ ensure that there are also nine independent components in $`2_\alpha _\beta \xi _c^{\mu \alpha \nu \beta }`$. It is easy to show that, for this combination to satisfy the correlation relation (98) with $`\mathrm{\Delta }\xi =0`$, the relevant correlators for the new stochastic tensor must be
$$\xi _c^{\mu \nu \alpha \beta }(x)\xi _c^{\rho \sigma \lambda \theta }(y)_{\xi _c}=T^{\mu \nu \alpha \beta \rho \sigma \lambda \theta }N(xy),$$
(99)
where $`T^{\mu \nu \alpha \beta \rho \sigma \lambda \theta }`$ is a linear combination of terms like $`\eta ^{\mu \rho }\eta ^{\nu \sigma }\eta ^{\alpha \lambda }\eta ^{\beta \theta }`$ in such a way that it has the same symmetries as the product of two Weyl tensor components $`C^{\mu \nu \alpha \beta }C^{\rho \sigma \lambda \theta }`$, its explicit expression is given in Ref. . Thus, after a redefinition of the arbitrary mass scale $`\mu `$ in Eq. (97) to absorb the constants of proportionality of the local terms with $`A^{\left(1\right)\mu \nu }(x)`$, one can see that the resulting equations for the $`\mathrm{\Delta }\xi =0`$ case are actually equivalent to those found in Ref. .
### C Expectation value of the stress-energy tensor
From the above equations one may extract the expectation value of the renormalized stress-energy tensor for a scalar field in a spacetime $`(\mathrm{IR}^4,\eta _{ab}+h_{ab})`$, computed up to first order in perturbation theory around the trivial solution of semiclassical gravity. Such an expectation value can be obtained by identification of Eqs. (81) and (97) with the components of the physical Einstein-Langevin equation, which in our particular case simply reads
$$\frac{1}{8\pi G}G^{\left(1\right)\mu \nu }2\left(\alpha A^{\left(1\right)\mu \nu }+\beta B^{\left(1\right)\mu \nu }\right)=\widehat{T}_R^{\mu \nu }[\eta +h]+2\xi ^{\mu \nu }.$$
(100)
By comparison of Eqs. (81) and (97) with the last equation, we can identify
$`\widehat{T}_R^{\mu \nu }(x)[\eta +h]={\displaystyle \frac{1}{2880\pi ^2}}\left[{\displaystyle \frac{8}{5}}A^{\left(1\right)\mu \nu }(x)({\displaystyle \frac{1}{6}}10\mathrm{\Delta }\xi )B^{\left(1\right)\mu \nu }(x)\right]`$ (101)
$`{\displaystyle d^4y\left[H_\mathrm{A}(xy;m^2)A^{\left(1\right)\mu \nu }(y)+H_\mathrm{B}(xy;m^2,\mathrm{\Delta }\xi )B^{\left(1\right)\mu \nu }(y)\right]}+O(h^2),`$ (102)
for a massive scalar field, $`m0`$, and
$`\widehat{T}_R^{\mu \nu }(x)[\eta +h]=`$ $`{\displaystyle \frac{1}{2880\pi ^2}}\left[{\displaystyle \frac{8}{5}}A^{\left(1\right)\mu \nu }(x)({\displaystyle \frac{1}{6}}10\mathrm{\Delta }\xi )B^{\left(1\right)\mu \nu }(x)\right]`$ (104)
$`{\displaystyle d^4yH(xy;\mu ^2)\left[A^{\left(1\right)\mu \nu }(y)+60\mathrm{\Delta }\xi ^2B^{\left(1\right)\mu \nu }(y)\right]}+O(h^2),`$
for a massless scalar field, $`m=0`$. Notice that in the massive case we have chosen, as usual, a renormalization scheme such that the expectation value of the renormalized stress-energy tensor does not have local terms proportional to the metric and the Einstein tensor . The result (104) agrees with the general form found by Horowitz using an axiomatic approach and coincides with that given in Ref. . The particular cases of conformal coupling, $`\mathrm{\Delta }\xi =0`$, and minimal coupling, $`\mathrm{\Delta }\xi =1/6`$, are also in agreement with the results for this cases given in Refs. (modulo local terms proportional to $`A^{\left(1\right)\mu \nu }`$ and $`B^{\left(1\right)\mu \nu }`$ due to different choices of the renormalization scheme). For the case of a massive minimally coupled scalar field, $`\mathrm{\Delta }\xi =1/6`$, our result (102) is equivalent to that of Ref. .
As it was pointed out above, in the case of conformal coupling, both for massive and massless scalar fields, one has $`\beta _B=\beta `$. This means that, in these cases, the terms proportional to $`B^{\left(1\right)\mu \nu }`$ in the above expectation values of the stress-energy tensor are actually independent of the renormalization scheme chosen. Due to the conformal invariance of $`d^4x\sqrt{g}C_{cabd}C^{cabd}`$, the tensor $`A^{ab}`$ is traceless and we have $`A^{\left(1\right)}_\mu ^\mu =0`$. Therefore, the terms with $`B^{\left(1\right)\mu \nu }`$ are precisely those which give trace to the expectation value of the stress-energy tensor in (102) and (104). In the massless conformally coupled case, $`m=0`$ and $`\mathrm{\Delta }\xi =0`$, such terms give the trace anomaly up to first order in $`h_{\mu \nu }`$:
$$\widehat{T}_R^\mu _\mu (x)[\eta +h]=\frac{1}{2880\pi ^2}\frac{1}{6}B^{\left(1\right)}_\mu ^\mu +O(h^2)=\frac{1}{2880\pi ^2}\mathrm{}R^{\left(1\right)}+O(h^2),$$
(105)
where we have used expression (E3) for $`B^{\left(1\right)\mu \nu }`$.
### D Particle creation
We can also use the result (40) for the noise kernel to evaluate the total probability of particle creation and the number of created particles for a real scalar field in a spacetime $`(\mathrm{IR}^4,\eta _{ab}+h_{ab})`$. The metric perturbation $`h_{ab}`$ (here an arbitrary perturbation) is assumed to vanish, either in an exact way or “asymptotically,” in the “remote past” and in the “far future,” so that the scalar field has well defined “in” and “out” many particle states. In that case, the absolute value of the logarithm of the vacuum persistence probability $`|0,\mathrm{out}|0,\mathrm{in}|^2`$, where $`|0,\mathrm{in}`$ and $`|0,\mathrm{out}`$ are, respectively, the “in” and “out” vacua in the Heisenberg picture, gives a measure of the total probability of particle creation. On the other hand, the number of created particles can be defined as the expectation value in the “in” vacuum of the number operator for “out” particles. As it was shown in Ref. , the total probability of particle creation and one half of the number of created particles coincide to lowest non-trivial order in the metric perturbation, these are
$$P[h]=d^4xd^4yh_{\mu \nu }(x)N^{\mu \nu \alpha \beta }(x,y)h_{\alpha \beta }(y)+0(h^3),$$
(106)
where $`N^{\mu \nu \alpha \beta }(x,y)`$ is the noise kernel given in (40), which in the massless case reduces to (98). The above expression for the total probability of pair creation by metric perturbations about Minkowski spacetime was first derived in Ref. . Using (40), we can write $`P[h]=P_\mathrm{A}[h]+P_\mathrm{B}[h]+0(h^3)`$, where
$`P_\mathrm{A}[h]{\displaystyle \frac{1}{6}}{\displaystyle d^4xd^4y(3_x^{\mu \alpha }_x^{\nu \beta }_x^{\mu \nu }_x^{\alpha \beta })N_\mathrm{A}(xy;m^2)h_{\mu \nu }(x)h_{\alpha \beta }(y)},`$ (107)
$`P_\mathrm{B}[h]{\displaystyle d^4xd^4y_x^{\mu \nu }_x^{\alpha \beta }N_\mathrm{B}(xy;m^2,\mathrm{\Delta }\xi )h_{\mu \nu }(x)h_{\alpha \beta }(y)}.`$ (108)
Integrating by parts (we always neglect surface terms), using expression (E5) for $`R^{\left(1\right)}`$, which can also be written as $`R^{\left(1\right)}=^{\mu \nu }h_{\mu \nu }`$, we find
$$P_\mathrm{B}[h]=d^4xd^4yR^{\left(1\right)}(x)N_\mathrm{B}(xy;m^2,\mathrm{\Delta }\xi )R^{\left(1\right)}(y).$$
(109)
In order to work out $`P_\mathrm{A}[h]`$, it is useful to take into account that, using the symmetry properties of the Weyl and Riemann tensors and the expression (E6) for $`R^{\left(1\right)\rho \sigma \lambda \tau }`$, one can write
$$C_{\rho \sigma \lambda \tau }^{\left(1\right)}(x)C^{\left(1\right)\rho \sigma \lambda \tau }(y)=C_{\rho \sigma \lambda \tau }^{\left(1\right)}(x)R^{\left(1\right)\rho \sigma \lambda \tau }(y)=2C^{\left(1\right)\rho \sigma \lambda \tau }(x)\delta _\rho ^\alpha \delta _\lambda ^\beta _\sigma _\tau h_{\alpha \beta }(y).$$
(110)
Using the last identity, the expression (E11) for $`C^{\left(1\right)\rho \sigma \lambda \tau }`$ and integrating by parts the first expression in (108) we get
$$P_\mathrm{A}[h]=d^4xd^4yC_{\mu \nu \alpha \beta }^{\left(1\right)}(x)N_\mathrm{A}(xy;m^2)C^{\left(1\right)\mu \nu \alpha \beta }(y).$$
(111)
Thus, $`P_\mathrm{A}[h]`$ and $`P_\mathrm{B}[h]`$ depend, respectively, on the Weyl tensor and the scalar curvature to first order in the metric perturbation. The result for the massless case, $`m=0`$, can be easily obtained from the above expressions, using Eqs. (89). If, in addition, we make $`\mathrm{\Delta }\xi =0`$, i.e., conformal coupling, we have $`P_\mathrm{B}[h]=0`$. Hence, for a conformal scalar field, particle creation is due to the breaking of conformal flatness in the spacetime, which implies a non-zero Weyl tensor.
In order to compare with previously obtained results, it is useful to introduce the Fourier transform of a field $`f(x)`$ as $`\stackrel{~}{f}(p)d^4xe^{ipx}f(x)`$. Note that, if $`f(x)`$ is real, then $`\stackrel{~}{f}(p)=\stackrel{~}{f}^{}(p)`$. Using the expressions (38) for the kernels $`N_\mathrm{A}`$ and $`N_\mathrm{B}`$, the above result for the total probability of particle creation and the number of particles created can also be written as
$`P[h]={\displaystyle \frac{1}{1920\pi }}{\displaystyle \frac{d^4p}{(2\pi )^4}\theta (p^24m^2)\sqrt{1+4\frac{m^2}{p^2}}}`$ $`[\stackrel{~}{C}_{\mu \nu \alpha \beta }^{\left(1\right)}(p)\stackrel{~}{C}^{\left(1\right)\mu \nu \alpha \beta }(p)(1+4{\displaystyle \frac{m^2}{p^2}})^2`$ (113)
$`+{\displaystyle \frac{20}{3}}\left|\stackrel{~}{R}^{\left(1\right)}(p)|^2(3\mathrm{\Delta }\xi +{\displaystyle \frac{m^2}{p^2}})^2\right]+O(h^3),`$
in agreement with the results of Ref. (except for a sign in the coefficient of the term with $`|\stackrel{~}{R}^{\left(1\right)}(p)|^2`$). It is also easy to see that the above result is equivalent to that found in Ref. if we take into account that, for integrals of the form $`Id^4p\stackrel{~}{f}_{a_1\mathrm{}a_r}(p)G(p^2)\stackrel{~}{f}^{a_1\mathrm{}a_r}(p)`$, where $`f_{a_1\mathrm{}a_r}(x)`$ is any real tensor field in Minkowski spacetime and $`G(p^2)`$ is any scalar function of $`p^2`$, one has that
$$I=2d^4p\theta (p^0)\stackrel{~}{f}_{a_1\mathrm{}a_r}(p)G(p^2)\stackrel{~}{f}^{a_1\mathrm{}a_r}(p)=2d^4p\theta (p^0)\stackrel{~}{f}_{a_1\mathrm{}a_r}(p)G(p^2)\stackrel{~}{f}^{a_1\mathrm{}a_r}(p).$$
(114)
In the massless conformally coupled case, $`m=0`$ and $`\mathrm{\Delta }\xi =0`$, the result (113) reduces to that found in Ref. .
The energy of the created particles, $`E[h]`$, defined as the expectation value of the “out” energy operator in the “in” vacuum can be computed using the expressions derived in Ref. . We find that this energy is given by an expression like (113), but with a factor $`2p^0\theta (p^0)`$ inserted in the integrand . Since the kernels $`N_\mathrm{A}`$ and $`D_\mathrm{A}`$ are related by the fluctuation-dissipation relation (65), and the same holds for $`N_\mathrm{B}`$ and $`D_\mathrm{B}`$, it is easy to see \[similarly to (114)\] that
$$E[h]=i\frac{d^4p}{(2\pi )^4}p^0\left[\stackrel{~}{C}_{\mu \nu \alpha \beta }^{\left(1\right)}(p)\stackrel{~}{C}^{\left(1\right)\mu \nu \alpha \beta }(p)\stackrel{~}{D}_\mathrm{A}(p)+\left|\stackrel{~}{R}^{\left(1\right)}(p)\right|^2\stackrel{~}{D}_\mathrm{B}(p)\right]+O(h^3),$$
(115)
where $`\stackrel{~}{D}_\mathrm{A}(p)`$ and $`\stackrel{~}{D}_\mathrm{B}(p)`$ are the Fourier transforms of the dissipation kernels defined in (38). For perturbations of a spatially flat RW spacetime (i.e., $`h_{\mu \nu }=2\mathrm{\Delta }a(\eta )\eta _{\mu \nu }`$, where $`x^0\eta `$ is the conformal time and $`\mathrm{\Delta }a(\eta )`$ is the perturbation of the scale factor), this last expression agrees with that of Ref. , see also Ref. .
So far in this subsection the metric perturbations are arbitrary. We may also be interested in the particles created by the back reaction on the metric due to the stress-energy fluctuations. Then we would have to use the solutions of the Einstein-Langevin equations (81) and (97) in the above results. However, to be consistent, one should look for solutions whose moments vanish asymptotically in the “remote past” and in the “far future.” These conditions are generally too strong, since they would break the time translation invariance in the correlation functions. In fact, the solutions that we find in the next section do not satisfy these conditions.
## V Correlation functions for gravitational perturbations
In this section, we solve the semiclassical Einstein-Langevin equations (81) and (97) for the components $`G^{\left(1\right)\mu \nu }`$ of the linearized Einstein tensor. In subsection V A we use these solutions to compute the corresponding two-point correlation functions, which give a measure of the gravitational fluctuations predicted by the stochastic semiclassical theory of gravity in the present case. Since the linearized Einstein tensor is invariant under gauge transformations of the metric perturbations, these two-point correlation functions are also gauge invariant. Once we have computed the two-point correlation functions for the linearized Einstein tensor, we find solutions for the metric perturbations in subsection V C and we show how the associated two-point correlation functions can be computed. This procedure to solve the Einstein-Langevin equations is similar to the one used by Horowitz , see also Ref. , to analyze the stability of Minkowski spacetime in semiclassical gravity.
From expressions (E2) and (E3) restricted to $`n=4`$, it is easy to see that $`A^{\left(1\right)\mu \nu }`$ and $`B^{\left(1\right)\mu \nu }`$ can be written in terms of $`G^{\left(1\right)\mu \nu }`$ as
$$A^{\left(1\right)\mu \nu }=\frac{2}{3}(^{\mu \nu }G^{\left(1\right)}_\alpha ^\alpha _\alpha ^\alpha G^{\left(1\right)\mu \nu }),B^{\left(1\right)\mu \nu }=2^{\mu \nu }G^{\left(1\right)}_\alpha ^\alpha ,$$
(116)
where we have used that $`3\mathrm{}=_\alpha ^\alpha `$. Therefore, the Einstein-Langevin equations (81) and (97) can be seen as linear integro-differential stochastic equations for the components $`G^{\left(1\right)\mu \nu }`$. Such equations can be written in both cases, $`m0`$ and $`m=0`$, as
$$\frac{1}{8\pi G}G^{\left(1\right)\mu \nu }(x)2\left(\overline{\alpha }A^{\left(1\right)\mu \nu }(x)+\overline{\beta }B^{\left(1\right)\mu \nu }(x)\right)+d^4y\left[H_\mathrm{A}(xy)A^{\left(1\right)\mu \nu }(y)+H_\mathrm{B}(xy)B^{\left(1\right)\mu \nu }(y)\right]=2\xi ^{\mu \nu }(x),$$
(117)
where the new constants $`\overline{\alpha }`$ and $`\overline{\beta }`$, and the kernels $`H_\mathrm{A}(x)`$ and $`H_\mathrm{B}(x)`$ can be identified in each case by comparison of this last equation with Eqs. (81) and (97). For instance, when $`m=0`$, we have $`H_\mathrm{A}(x)=H(x;\mu ^2)`$ and $`H_\mathrm{B}(x)=60\mathrm{\Delta }\xi ^2H(x;\mu ^2)`$. In this case, we can use the arbitrariness of the mass scale $`\mu `$ to eliminate one of the parameters $`\overline{\alpha }`$ or $`\overline{\beta }`$.
In order to find solutions to these equations, it is convenient to Fourier transform them. Introducing Fourier transforms as in subsection IV D, one finds, from (116),
$$\stackrel{~}{A}^{\left(1\right)\mu \nu }(p)=2p^2\stackrel{~}{G}^{\left(1\right)\mu \nu }(p)\frac{2}{3}p^2P^{\mu \nu }\stackrel{~}{G}^{\left(1\right)}_\alpha ^\alpha (p),\stackrel{~}{B}^{\left(1\right)\mu \nu }(p)=2p^2P^{\mu \nu }\stackrel{~}{G}^{\left(1\right)}_\alpha ^\alpha (p).$$
(118)
Using these relations, the Fourier transform of Eq. (117) reads
$$F_{\alpha \beta }^{\mu \nu }(p)\stackrel{~}{G}^{\left(1\right)\alpha \beta }(p)=16\pi G\stackrel{~}{\xi }^{\mu \nu }(p),$$
(119)
where
$$F_{\alpha \beta }^{\mu \nu }(p)F_1(p)\delta _{(\alpha }^\mu \delta _{\beta )}^\nu +F_2(p)p^2P^{\mu \nu }\eta _{\alpha \beta },$$
(120)
with
$$F_1(p)1+16\pi Gp^2\left[\stackrel{~}{H}_\mathrm{A}(p)2\overline{\alpha }\right],F_2(p)\frac{16}{3}\pi G\left[\stackrel{~}{H}_\mathrm{A}(p)+3\stackrel{~}{H}_\mathrm{B}(p)2\overline{\alpha }6\overline{\beta }\right].$$
(121)
In Eq. (119), $`\stackrel{~}{\xi }^{\mu \nu }(p)`$, the Fourier transform of $`\xi ^{\mu \nu }(x)`$, is a Gaussian stochastic source of zero average and
$$\stackrel{~}{\xi }^{\mu \nu }(p)\stackrel{~}{\xi }^{\alpha \beta }(p^{})_c=(2\pi )^4\delta ^4(p+p^{})\stackrel{~}{N}^{\mu \nu \alpha \beta }(p),$$
(122)
where we have introduced the Fourier transform of the noise kernel. The explicit expression for $`\stackrel{~}{N}^{\mu \nu \alpha \beta }(p)`$ is found from (40) and (38) to be
$`\stackrel{~}{N}^{\mu \nu \alpha \beta }(p)={\displaystyle \frac{1}{2880\pi }}\theta (p^24m^2)\sqrt{1+4{\displaystyle \frac{m^2}{p^2}}}`$ $`[{\displaystyle \frac{1}{4}}(1+4{\displaystyle \frac{m^2}{p^2}})^2(p^2)^2(3P^{\mu (\alpha }P^{\beta )\nu }P^{\mu \nu }P^{\alpha \beta })`$ (124)
$`+\mathrm{\hspace{0.17em}10}(3\mathrm{\Delta }\xi +{\displaystyle \frac{m^2}{p^2}})^2(p^2)^2P^{\mu \nu }P^{\alpha \beta }],`$
which in the massless case reduces to
$$\underset{m0}{lim}\stackrel{~}{N}^{\mu \nu \alpha \beta }(p)=\frac{1}{1920\pi }\theta (p^2)\left[\frac{1}{6}(p^2)^2\left(3P^{\mu (\alpha }P^{\beta )\nu }P^{\mu \nu }P^{\alpha \beta }\right)+60\mathrm{\Delta }\xi ^2(p^2)^2P^{\mu \nu }P^{\alpha \beta }\right].$$
(125)
### A Correlation functions for the linearized Einstein tensor
In general, we can write $`G^{\left(1\right)\mu \nu }=G^{\left(1\right)\mu \nu }_c+G_\mathrm{f}^{\left(1\right)\mu \nu }`$, where $`G_\mathrm{f}^{\left(1\right)\mu \nu }`$ is a solution to Eqs. (117) \[or, in the Fourier transformed version, (119)\] with zero average. The averages $`G^{\left(1\right)\mu \nu }_c`$ must be a solution of the linearized semiclassical Einstein equations obtained by averaging Eqs. (117) \[or (119)\]. Solutions to these equations (specially in the massless case, $`m=0`$) have been studied by several authors , particularly in connection with the issue of the stability of the trivial solutions of semiclassical gravity. The two-point correlation functions for the linearized Einstein tensor are given by
$$𝒢^{\mu \nu \alpha \beta }(x,x^{})G^{\left(1\right)\mu \nu }(x)G^{\left(1\right)\alpha \beta }(x^{})_cG^{\left(1\right)\mu \nu }(x)_cG^{\left(1\right)\alpha \beta }(x^{})_c=G_\mathrm{f}^{\left(1\right)\mu \nu }(x)G_\mathrm{f}^{\left(1\right)\alpha \beta }(x^{})_c.$$
(126)
Next, we shall seek the family of solutions to the Einstein-Langevin equations which can be written as a linear functional of the stochastic source and whose Fourier transform, $`\stackrel{~}{G}^{\left(1\right)\mu \nu }(p)`$, depends locally on $`\stackrel{~}{\xi }^{\alpha \beta }(p)`$. Each of such solutions is a Gaussian stochastic field and, thus, it can be completely characterized by the averages $`G^{\left(1\right)\mu \nu }_c`$ and the two-point correlation functions (126). For such a family of solutions, $`\stackrel{~}{G}_\mathrm{f}^{\left(1\right)\mu \nu }(p)`$ is the most general solution to Eq. (119) which is linear, homogeneous and local in $`\stackrel{~}{\xi }^{\alpha \beta }(p)`$. It can be written as
$$\stackrel{~}{G}_\mathrm{f}^{\left(1\right)\mu \nu }(p)=16\pi GD_{\alpha \beta }^{\mu \nu }(p)\stackrel{~}{\xi }^{\alpha \beta }(p),$$
(127)
where $`D_{\alpha \beta }^{\mu \nu }(p)`$ are the components of a Lorentz invariant tensor field distribution in Minkowski spacetime \[by “Lorentz invariant” we mean invariant under the transformations of the orthochronous Lorentz subgroup; see Ref. for more details on the definition and properties of these tensor distributions\], symmetric under the interchanges $`\alpha \beta `$ and $`\mu \nu `$, which is the most general solution of
$$F_{\rho \sigma }^{\mu \nu }(p)D_{\alpha \beta }^{\rho \sigma }(p)=\delta _{(\alpha }^\mu \delta _{\beta )}^\nu .$$
(128)
In addition, we must impose the conservation condition to the solutions: $`p_\nu \stackrel{~}{G}_\mathrm{f}^{\left(1\right)\mu \nu }(p)=0`$, where this zero must be understood as a stochastic variable which behaves deterministically as a zero vector field. We can write $`D_{\alpha \beta }^{\mu \nu }(p)=D_{p\alpha \beta }^{\mu \nu }(p)+D_{h\alpha \beta }^{\mu \nu }(p)`$, where $`D_{p\alpha \beta }^{\mu \nu }(p)`$ is a particular solution to Eq. (128) and $`D_{h\alpha \beta }^{\mu \nu }(p)`$ is the most general solution to the corresponding homogeneous equation. Correspondingly \[see Eq. (127)\], we can write $`\stackrel{~}{G}_\mathrm{f}^{\left(1\right)\mu \nu }(p)=\stackrel{~}{G}_p^{\left(1\right)\mu \nu }(p)+\stackrel{~}{G}_h^{\left(1\right)\mu \nu }(p)`$. To find the particular solution, we try an ansatz of the form
$$D_{p\alpha \beta }^{\mu \nu }(p)=d_1(p)\delta _{(\alpha }^\mu \delta _{\beta )}^\nu +d_2(p)p^2P^{\mu \nu }\eta _{\alpha \beta }.$$
(129)
Substituting this ansatz into Eqs. (128), it is easy to see that it solves these equations if we take
$$d_1(p)=\left[\frac{1}{F_1(p)}\right]_r,d_2(p)=\left[\frac{F_2(p)}{F_1(p)F_3(p)}\right]_r,$$
(130)
with
$$F_3(p)F_1(p)+3p^2F_2(p)=148\pi Gp^2\left[\stackrel{~}{H}_\mathrm{B}(p)2\overline{\beta }\right],$$
(131)
and where the notation $`[]_r`$ means that the zeros of the denominators are regulated with appropriate prescriptions in such a way that $`d_1(p)`$ and $`d_2(p)`$ are well defined Lorentz invariant scalar distributions. This yields a particular solution to the Einstein-Langevin equations:
$$\stackrel{~}{G}_p^{\left(1\right)\mu \nu }(p)=16\pi GD_{p\alpha \beta }^{\mu \nu }(p)\stackrel{~}{\xi }^{\alpha \beta }(p),$$
(132)
which, since the stochastic source is conserved, satisfies the conservation condition. Note that, in the case of a massless scalar field, $`m=0`$, the above solution has a functional form analogous to that of the solutions of linearized semiclassical gravity found in the Appendix of Ref. . Notice also that, for a massless conformally coupled field, $`m=0`$ and $`\mathrm{\Delta }\xi =0`$, the second term in the right hand side of Eq. (129) will not contribute in the correlation functions (126), since, as we have pointed out in Sec. IV B, in this case the stochastic source is “traceless.”
Next, we can work out the general form for $`D_{h\alpha \beta }^{\mu \nu }(p)`$, which is a linear combination of terms consisting of a Lorentz invariant scalar distribution times one of the products $`\delta _{(\alpha }^\mu \delta _{\beta )}^\nu `$, $`p^2P^{\mu \nu }\eta _{\alpha \beta }`$, $`\eta ^{\mu \nu }\eta _{\alpha \beta }`$, $`\eta ^{\mu \nu }p^2P_{\alpha \beta }`$, $`\delta _{(\alpha }^{(\mu }p^2P_{\beta )}^{\nu )}`$ and $`p^2P^{\mu \nu }p^2P_{\alpha \beta }`$. However, taking into account that the stochastic source is conserved, we can omit some terms in $`D_{h\alpha \beta }^{\mu \nu }(p)`$ and simply write
$$\stackrel{~}{G}_h^{\left(1\right)\mu \nu }(p)=16\pi GD_{h\alpha \beta }^{\mu \nu }(p)\stackrel{~}{\xi }^{\alpha \beta }(p),$$
(133)
with
$$D_{h\alpha \beta }^{\mu \nu }(p)=h_1(p)\delta _{(\alpha }^\mu \delta _{\beta )}^\nu +h_2(p)p^2P^{\mu \nu }\eta _{\alpha \beta }+h_3(p)\eta ^{\mu \nu }\eta _{\alpha \beta },$$
(134)
where $`h_1(p)`$, $`h_2(p)`$ and $`h_3(p)`$ are Lorentz invariant scalar distributions. From the fact that $`D_{h\alpha \beta }^{\mu \nu }(p)`$ must satisfy the homogeneous equation corresponding to Eq. (128), we find that $`h_1(p)`$ and $`h_3(p)`$ have support on the set of points $`\{p^\mu \}`$ for which $`F_1(p)=0`$, and that $`h_2(p)`$ has support on the set of points $`\{p^\mu \}`$ for which $`F_1(p)=0`$ or $`F_3(p)=0`$. Moreover, the conservation condition for $`\stackrel{~}{G}_h^{\left(1\right)\mu \nu }(p)`$ implies that the term with $`h_3(p)`$ is only allowed in the case of a massless conformally coupled field, $`m=0`$ and $`\mathrm{\Delta }\xi =0`$. From (122), we get
$$\stackrel{~}{G}_h^{\left(1\right)\mu \nu }(p)\stackrel{~}{\xi }^{\alpha \beta }(p^{})_c=(2\pi )^4\mathrm{\hspace{0.17em}16}\pi G\delta ^4(p+p^{})D_{h\rho \sigma }^{\mu \nu }(p)\stackrel{~}{N}^{\rho \sigma \alpha \beta }(p).$$
(135)
Note, from expressions (124) and (125), that the support of $`\stackrel{~}{N}^{\mu \nu \alpha \beta }(p)`$ is on the set of points $`\{p^\mu \}`$ for which $`p^20`$ when $`m=0`$, and for which $`p^24m^2>0`$ when $`m0`$. At such points, using expressions (121), (131), (92) and (79), it is easy to see that $`F_1(p)`$ is always different from zero, and that $`F_3(p)`$ is also always different from zero, except for some particular values of $`\mathrm{\Delta }\xi `$ and $`\overline{\beta }`$:
* when $`m=0`$, $`\mathrm{\Delta }\xi =0`$ and $`\overline{\beta }>0`$;
* when $`m0`$, $`0<\mathrm{\Delta }\xi <(1/12)`$ and $`\overline{\beta }=(\mathrm{\Delta }\xi /32\pi ^2)[\pi /(Gm^2)+1/36]`$.
In the case a), $`F_3(p)=0`$ for the set of points $`\{p^\mu \}`$ satisfying $`p^2=1/(96\pi G\overline{\beta })`$; in the case b), $`F_3(p)=0`$ for $`\{p^\mu \}`$ such that $`p^2=m^2/(3\mathrm{\Delta }\xi )`$. Hence, except for the above cases a) and b), the intersection of the supports of $`\stackrel{~}{N}^{\mu \nu \alpha \beta }(p)`$ and $`D_{h\lambda \gamma }^{\rho \sigma }(p)`$ is an empty set and, thus, the correlation function (135) is zero. In the cases a) and b), we can have a contribution to (135) coming from the term with $`h_2(p)`$ in (134) of the form $`D_{h\rho \sigma }^{\mu \nu }(p)\stackrel{~}{N}^{\rho \sigma \alpha \beta }(p)=H_3(p;\{C\})p^2P^{\mu \nu }\stackrel{~}{N}_\rho ^{\alpha \beta \rho }(p)`$, where $`H_3(p;\{C\})`$ is the most general Lorentz invariant distribution satisfying $`F_3(p)H_3(p;\{C\})=0`$, which depends on a set of arbitrary parameters represented as $`\{C\}`$. However, from (124), we see that $`\stackrel{~}{N}_\rho ^{\alpha \beta \rho }(p)`$ is proportional to $`\theta (p^24m^2)(1+4m^2/p^2)^{(1/2)}(3\mathrm{\Delta }\xi +m^2/p^2)^2`$. Thus, in the case a), we have $`\stackrel{~}{N}_\rho ^{\alpha \beta \rho }(p)=0`$ and, in the case b), the intersection of the supports of $`\stackrel{~}{N}_\rho ^{\alpha \beta \rho }(p)`$ and of $`H_3(p;\{C\})`$ is an empty set. Therefore, from the above analysis, we conclude that $`\stackrel{~}{G}_h^{\left(1\right)\mu \nu }(p)`$ gives no contribution to the correlation functions (126), since $`\stackrel{~}{G}_h^{\left(1\right)\mu \nu }(p)\stackrel{~}{\xi }^{\alpha \beta }(p^{})_c=0`$, and we have simply $`𝒢^{\mu \nu \alpha \beta }(x,x^{})=G_p^{\left(1\right)\mu \nu }(x)G_p^{\left(1\right)\alpha \beta }(x^{})_c`$, where $`G_p^{\left(1\right)\mu \nu }(x)`$ is the inverse Fourier transform of (132).
The correlation functions (126) can then be computed from
$$\stackrel{~}{G}_p^{\left(1\right)\mu \nu }(p)\stackrel{~}{G}_p^{\left(1\right)\alpha \beta }(p^{})_c=64(2\pi )^6G^2\delta ^4(p+p^{})D_{p\rho \sigma }^{\mu \nu }(p)D_{p\lambda \gamma }^{\alpha \beta }(p)\stackrel{~}{N}^{\rho \sigma \lambda \gamma }(p).$$
(136)
It is easy to see from the above analysis that the prescriptions $`[]_r`$ in the factors $`D_p`$ are irrelevant in the last expression and, thus, they can be suppressed. Taking into account that $`F_l(p)=F_l^{}(p)`$, with $`l=1,2,3`$, we get from Eqs. (129) and (130)
$`\stackrel{~}{G}_p^{\left(1\right)\mu \nu }(p)\stackrel{~}{G}_p^{\left(1\right)\alpha \beta }(p^{})_c=`$ $`64(2\pi )^6G^2{\displaystyle \frac{\delta ^4(p+p^{})}{\left|F_1(p)\right|^2}}[\stackrel{~}{N}^{\mu \nu \alpha \beta }(p){\displaystyle \frac{F_2(p)}{F_3(p)}}p^2P^{\mu \nu }\stackrel{~}{N}_\rho ^{\alpha \beta \rho }(p)`$ (138)
$`{\displaystyle \frac{F_2^{}(p)}{F_3^{}(p)}}p^2P^{\alpha \beta }\stackrel{~}{N}_\rho ^{\mu \nu \rho }(p)+{\displaystyle \frac{\left|F_2(p)\right|^2}{\left|F_3(p)\right|^2}}p^2P^{\mu \nu }p^2P^{\alpha \beta }\stackrel{~}{N}_{\rho \sigma }^{\rho \sigma }(p)].`$
This last expression is well defined as a bi-distribution and can be easily evaluated using Eq. (124). We find
$`\stackrel{~}{G}_p^{\left(1\right)\mu \nu }(p)\stackrel{~}{G}_p^{\left(1\right)\alpha \beta }(p^{})_c={\displaystyle \frac{2}{45}}`$ $`(2\pi )^5G^2{\displaystyle \frac{\delta ^4(p+p^{})}{\left|F_1(p)\right|^2}}\theta (p^24m^2)\sqrt{1+4{\displaystyle \frac{m^2}{p^2}}}`$ (142)
$`\times [{\displaystyle \frac{1}{4}}(1+4{\displaystyle \frac{m^2}{p^2}})^2(p^2)^2(3P^{\mu (\alpha }P^{\beta )\nu }P^{\mu \nu }P^{\alpha \beta })`$
$`+\mathrm{\hspace{0.17em}10}(3\mathrm{\Delta }\xi +{\displaystyle \frac{m^2}{p^2}})^2(p^2)^2P^{\mu \nu }P^{\alpha \beta }|13p^2{\displaystyle \frac{F_2(p)}{F_3(p)}}|^2].`$
To derive the correlation functions (126), we have to take the inverse Fourier transform of the above result. We finally obtain
$$𝒢^{\mu \nu \alpha \beta }(x,x^{})=\frac{\pi }{45}G^2_x^{\mu \nu \alpha \beta }𝒢_\mathrm{A}(xx^{})+\frac{8\pi }{9}G^2_x^{\mu \nu }_x^{\alpha \beta }𝒢_\mathrm{B}(xx^{}),$$
(143)
with
$`\stackrel{~}{𝒢}_\mathrm{A}(p)\theta (p^24m^2)\sqrt{1+4{\displaystyle \frac{m^2}{p^2}}}\left(1+4{\displaystyle \frac{m^2}{p^2}}\right)^2{\displaystyle \frac{1}{\left|F_1(p)\right|^2}},`$ (144)
$`\stackrel{~}{𝒢}_\mathrm{B}(p)\theta (p^24m^2)\sqrt{1+4{\displaystyle \frac{m^2}{p^2}}}\left(3\mathrm{\Delta }\xi +{\displaystyle \frac{m^2}{p^2}}\right)^2{\displaystyle \frac{1}{\left|F_1(p)\right|^2}}\left|13p^2{\displaystyle \frac{F_2(p)}{F_3(p)}}\right|^2,`$ (145)
and $`_x^{\mu \nu \alpha \beta }3_x^{\mu (\alpha }_x^{\beta )\nu }_x^{\mu \nu }_x^{\alpha \beta }`$, and where $`F_l(p)`$, $`l=1,2,3`$, are given in (121) and (131). Notice that, for a massless field ($`m=0`$), we have
$`F_1(p)=1+16\pi Gp^2\stackrel{~}{H}(p;\overline{\mu }^2),`$ (146)
$`F_2(p)={\displaystyle \frac{16}{3}}\pi G\left[(1+180\mathrm{\Delta }\xi ^2)\stackrel{~}{H}(p;\overline{\mu }^2)6\mathrm{{\rm Y}}\right],`$ (147)
$`F_3(p)=148\pi Gp^2\left[60\mathrm{\Delta }\xi ^2\stackrel{~}{H}(p;\overline{\mu }^2)2\mathrm{{\rm Y}}\right],`$ (148)
with $`\overline{\mu }\mu \mathrm{exp}(1920\pi ^2\overline{\alpha })`$ and $`\mathrm{{\rm Y}}\overline{\beta }60\mathrm{\Delta }\xi ^2\overline{\alpha }`$, and where $`\stackrel{~}{H}(p;\mu ^2)`$ is the Fourier transform of $`H(x;\mu ^2)`$ given in (92).
### B Conformal field case
The above correlation functions become simpler when the scalar field is massless and conformally coupled, i.e., when $`m=0`$ and $`\mathrm{\Delta }\xi =0`$, since in this case $`𝒢_\mathrm{B}(x)=0`$ and $`\stackrel{~}{𝒢}_\mathrm{A}(p)`$ reduces to $`\stackrel{~}{𝒢}_\mathrm{A}(p)=\theta (p^2)\left|F_1(p)\right|^2`$. Introducing the function $`\phi (\chi ;\lambda )\left[1\chi \mathrm{ln}\left(\lambda \chi /e\right)\right]^2+\pi ^2\chi ^2`$, with $`\chi 0`$ and $`\lambda >0`$, $`𝒢_\mathrm{A}(x)`$ can be written as
$$𝒢_\mathrm{A}(x)=\frac{(120\pi )^{3/2}}{2\pi ^3L_P^3}\frac{1}{|𝐱|}_0^{\mathrm{}}d|𝐪||𝐪|\mathrm{sin}\left[\frac{\sqrt{120\pi }}{L_P}|𝐱||𝐪|\right]_0^{\mathrm{}}𝑑q^0\mathrm{cos}\left[\frac{\sqrt{120\pi }}{L_P}x^0q^0\right]\frac{\theta (q^2)}{\phi (q^2;\lambda )},$$
(149)
where $`L_P\sqrt{G}`$ is the Planck length, $`\lambda 120\pi e/(L_P^2\overline{\mu }^2)`$, and we use the notation $`x^\mu =(x^0,𝐱)`$ and $`q^\mu =(q^0,𝐪)`$. Notice that, if we assume that $`\overline{\mu }L_P^1`$, then $`\lambda 10^3`$. For those values of the parameter $`\lambda `$ (and also for smaller values), the function $`\phi (\chi ;\lambda )`$ has a minimum at some value of $`\chi `$ that we denote as $`\chi _0(\lambda )`$. This can be found by solving the equation $`\pi ^2\chi _0=\left[1\chi _0\mathrm{ln}(\lambda \chi _0/e)\right]\left[1+\mathrm{ln}(\lambda \chi _0/e)\right]`$ numerically (discarding a solution $`\chi _M(\lambda )<\chi _0(\lambda )`$, at which the function $`\phi (\chi ;\lambda )`$ has a maximum). Since the main contribution to the integral (149) come from the values of $`q^2`$ around $`q^2=\chi _0(\lambda )`$, $`\phi (\chi ;\lambda )`$ can be approximately replaced in this integral by $`\phi _{\mathrm{ap}}(\chi ;\lambda )[1\kappa (\lambda )\chi ]^2+\pi ^2\chi ^2=[\kappa ^2(\lambda )+\pi ^2]\chi ^22\kappa (\lambda )\chi +1`$, with $`\kappa (\lambda )\mathrm{ln}\left(\lambda \chi _0(\lambda )/e\right)`$. For $`(\lambda /5)10^310^7`$, we have $`\kappa 10`$.
Let the spacetime points $`x`$ and $`x^{}`$ be different and spacelike separated. In this case, we can choose an inertial coordinate system for which $`(xx^{})^\mu =(0,𝐱𝐱^{})`$ and $`𝒢^{\mu \nu \alpha \beta }(x,x^{})`$ will be a function of $`𝐱𝐱^{}`$ only that can be written as
$$𝒢^{\mu \nu \alpha \beta }(𝐱𝐱^{})=𝒢_1^{\mu \nu \alpha \beta }(𝐱𝐱^{})+𝒢_2^{\mu \nu \alpha \beta }(𝐱𝐱^{})+𝒢_3^{\mu \nu \alpha \beta }(𝐱𝐱^{}),$$
(150)
with
$$𝒢_a^{\mu \nu \alpha \beta }(𝐱)\frac{\pi }{45}G^2_{a_𝐱}^{\mu \nu \alpha \beta }I_a(𝐱),$$
(151)
$`a=1,2,3`$, where $`I_1(𝐱)𝒢_\mathrm{A}(x)|_{x^\mu =(0,𝐱)}`$, $`I_2(𝐱)(_x^0)^2𝒢_\mathrm{A}(x)|_{x^\mu =(0,𝐱)}`$, $`I_3(𝐱)(_x^0)^4𝒢_\mathrm{A}(x)|_{x^\mu =(0,𝐱)}`$, and $`_{a_𝐱}^{\mu \nu \alpha \beta }`$ are some differential operators. Note that the terms containing an odd number of $`_x^0`$ derivatives are zero. The differential operators $`_{1_𝐱}^{\mu \nu \alpha \beta }`$ are given by $`_1^{\mu \nu \alpha \beta }=3𝒟^{\mu (\alpha }𝒟^{\beta )\nu }𝒟^{\mu \nu }𝒟^{\alpha \beta }`$, with $`𝒟^{\mu \nu }(\eta ^{\mu \nu }\delta ^{ij}\delta ^{\mu i}\delta ^{\nu j})_i_j`$. The non-null components of the remaining operators are $`_2^{00ij}=3^i^j\delta ^{ij}`$, $`_2^{0i0j}=\frac{1}{2}(^i^j+3\delta ^{ij})`$, $`_3^{ijkl}=\delta ^{ij}\delta ^{kl}+3\delta ^{i(k}\delta ^{l)j}`$, $`_2^{ijkl}=2(\delta ^{ij}\delta ^{kl}3\delta ^{i(k}\delta ^{l)j})\delta ^{ij}^k^l\delta ^{kl}^i^j+3(\delta ^{i(k}^{l)}^j+\delta ^{j(k}^{l)}^i)`$, where $`\delta ^{ij}_i_j`$ is the usual (Euclidean space) Laplace operator. From the above expressions, we can see that $`𝒢^{000i}(𝐱𝐱^{})=𝒢^{0ijk}(𝐱𝐱^{})=0`$, but the remaining correlation functions $`𝒢^{\mu \nu \alpha \beta }(𝐱𝐱^{})`$ are in principle non-null.
With the approximation described above, the integrals $`I_a(𝐱)`$ can be written as
$$I_a(𝐱)\frac{(1)^{a+1}}{2\pi ^3}\left(\frac{120\pi }{L_P^2}\right)^{a+1/2}\frac{1}{|𝐱|}_0^{\mathrm{}}d|𝐪|\mathrm{sin}\left[\frac{\sqrt{120\pi }}{L_P}|𝐱||𝐪|\right]|𝐪|J_a(|𝐪|),$$
(152)
where
$`J_1(|𝐪|){\displaystyle _{|𝐪|}^{\mathrm{}}}𝑑q^0{\displaystyle \frac{1}{\phi _{\mathrm{ap}}(q^2;\lambda )}},J_2(|𝐪|){\displaystyle _{|𝐪|}^{\mathrm{}}}𝑑q^0{\displaystyle \frac{(q^0)^2}{\phi _{\mathrm{ap}}(q^2;\lambda )}},`$ (153)
$`J_3(|𝐪|){\displaystyle \frac{|𝐪|}{\kappa ^2(\lambda )+\pi ^2}}+{\displaystyle _{|𝐪|}^{\mathrm{}}}𝑑q^0\left[{\displaystyle \frac{(q^0)^4}{\phi _{\mathrm{ap}}(q^2;\lambda )}}{\displaystyle \frac{1}{[\kappa ^2(\lambda )+\pi ^2]}}\right].`$ (154)
Noting that $`\phi _{\mathrm{ap}}(q^2;\lambda )`$ has four zeros in the complex $`q^0`$ plane at $`\pm p(|𝐪|)`$, $`\pm p^{}(|𝐪|)`$, where $`p(s)`$ (we make $`s|𝐪|`$) is the complex function with
$$\begin{array}{c}\mathrm{Re}p(s)\\ \mathrm{Im}p(s)\end{array}\}=\left[\frac{\sqrt{\left[(\kappa ^2+\pi ^2)s^2+\kappa \right]^2+\pi ^2}\pm (\kappa ^2+\pi ^2)s^2\pm \kappa }{2(\kappa ^2+\pi ^2)}\right]^{1/2},$$
(155)
we can decompose
$$\frac{1}{\phi _{\mathrm{ap}}(q^2;\lambda )}=\frac{1}{4(\kappa ^2+\pi ^2)}\frac{1}{|p|^2\mathrm{Re}p}\left[\frac{q^0+2\mathrm{R}\mathrm{e}p}{(q^0)^2+2\mathrm{Re}pq^0+|p|^2}\frac{(q^02\mathrm{R}\mathrm{e}p)}{(q^0)^22\mathrm{Re}pq^0+|p|^2}\right],$$
(156)
and then we can perform the integrals $`J_a(s)`$, $`a=1,2,3`$. The results for these integrals can be found in Appendix F.
Next, to carry on with the calculation, we need to introduce some suitable approximations for the functions $`J_a(s)`$ in the integrals (152). In order to do so, we study the behavior of these functions for small and large values of $`s`$. For $`sJ_1(s)`$, we find that it can be well approximated by an $`\mathrm{arctan}`$ function. In fact, on the one hand, $`sJ_1(s)`$ tends very quickly to a constant limiting value $`lim_s\mathrm{}sJ_1(s)=a/4`$, where $`a1+(2/\pi )\mathrm{arctan}(\kappa /\pi )`$. On the other hand, for small values of $`s`$, we have $`sJ_1(s)\left[\sqrt{120\pi }a/(2\pi b)\right]s+O(s^2)`$, with $`b(4a/\pi ^2)\left[15\pi \left(\sqrt{\kappa ^2+\pi ^2}\kappa \right)\right]^{1/2}`$. Hence, we can approximate
$$sJ_1(s)\frac{a}{2\pi }\mathrm{arctan}\left(\frac{\sqrt{120\pi }}{b}s\right).$$
(157)
Performing the integral $`I_1(𝐱)`$ \[see Eq. (152)\] with this approximation, we get, for $`|𝐱|0`$,
$$I_1(𝐱)\frac{15}{\pi ^2}\frac{a}{L_P^2}\frac{1}{|𝐱|^2}e^{b|𝐱|/L_P}.$$
(158)
The function $`J_2(s)`$ behaves as $`J_2(s)(a/4)s+O(s^1\mathrm{ln}s)`$ for large values of $`s`$, and as $`J_2(s)(a/4)(120\pi )^{1/2}\gamma +O(s^2)`$, with $`\gamma 240(\kappa ^2+\pi ^2)^{1/2}b^1`$, for small values of $`s`$. This function can be well approximated by
$$J_2(s)\frac{a}{4}\left[s^2+\frac{\gamma ^2}{120\pi }\right]^{1/2},$$
(159)
and, substituting the last expression in the integral $`I_2(𝐱)`$ \[see (152)\], we obtain, for $`|𝐱|0`$,
$$I_2(𝐱)\frac{15}{\pi ^2}\frac{a}{L_P^4}\frac{\gamma ^2}{|𝐱|^2}K_2\left(\gamma |𝐱|/L_P\right),$$
(160)
where $`K_\nu (z)`$ denote the modified Bessel functions of the second kind. For $`J_3(s)`$, we find that $`J_3(s)(a/4)s^3+O(s\mathrm{ln}s)`$ for large values of $`s`$, and that $`J_3(s)(a/4)(120\pi )^{3/2}\delta ^3+O(s)`$, with $`\delta 4(\kappa ^2+\pi ^2)^{1/2}\left[450\pi b^1\left(2\kappa \sqrt{\kappa ^2+\pi ^2}\right)\right]^{1/3}`$, for $`s`$ small. With the approximation
$$J_3(s)\frac{a}{4}\left[s^2+\frac{\delta ^2}{120\pi }\right]^{3/2},$$
(161)
we can compute the integral $`I_3(𝐱)`$ \[see (152)\] for $`|𝐱|0`$, and we find
$$I_3(𝐱)\frac{45}{\pi ^2}\frac{a}{L_P^5}\frac{\delta ^3}{|𝐱|^3}K_3\left(\delta |𝐱|/L_P\right).$$
(162)
Numerical calculations confirm that the above approximations are reasonable. For $`\kappa 10`$, we have $`a,b,\delta 1`$ and $`\gamma 10`$.
The results (158), (160) and (162) are now ready to be substituted into (151), from where we can compute the different contributions to the correlation functions (150). Using the relation $`(d/dz)\left[z^\nu K_\nu (z)\right]=z^\nu K_{\nu +1}(z)`$, and defining $`\sigma _bb|𝐱|/L_P`$, $`\sigma _\gamma \gamma |𝐱|/L_P`$, $`\sigma _\delta \delta |𝐱|/L_P`$, we get, after a rather long but straightforward calculation, the following results for the non-zero components of $`𝒢_a^{\mu \nu \alpha \beta }(𝐱)`$ \[with $`|𝐱|0`$\]:
$`𝒢_1^{0000}(𝐱){\displaystyle \frac{2}{3\pi }}{\displaystyle \frac{ab^6}{L_P^4}}{\displaystyle \frac{e^{\sigma _b}}{\sigma _b^2}}\left[1+{\displaystyle \frac{4}{\sigma _b}}+{\displaystyle \frac{12}{\sigma _b^2}}+{\displaystyle \frac{24}{\sigma _b^3}}+{\displaystyle \frac{24}{\sigma _b^4}}\right],`$ (163)
$`𝒢_1^{00ij}(𝐱){\displaystyle \frac{1}{3\pi }}{\displaystyle \frac{ab^6}{L_P^4}}{\displaystyle \frac{e^{\sigma _b}}{\sigma _b^2}}\left[\delta ^{ij}\left(1+{\displaystyle \frac{5}{\sigma _b}}+{\displaystyle \frac{16}{\sigma _b^2}}+{\displaystyle \frac{32}{\sigma _b^3}}+{\displaystyle \frac{32}{\sigma _b^4}}\right){\displaystyle \frac{x^ix^j}{|𝐱|^2}}\left(1+{\displaystyle \frac{7}{\sigma _b}}+{\displaystyle \frac{24}{\sigma _b^2}}+{\displaystyle \frac{48}{\sigma _b^3}}+{\displaystyle \frac{48}{\sigma _b^4}}\right)\right],`$ (164)
$`𝒢_1^{0i0j}(𝐱)={\displaystyle \frac{3}{2}}𝒢_1^{00ij}(𝐱),`$ (165)
$`𝒢_1^{ijkl}(𝐱){\displaystyle \frac{1}{3\pi }}{\displaystyle \frac{ab^6}{L_P^4}}{\displaystyle \frac{e^{\sigma _b}}{\sigma _b^2}}[(\delta ^{ij}\delta ^{kl}3\delta ^{i(k}\delta ^{l)j})(1+{\displaystyle \frac{6}{\sigma _b}}+{\displaystyle \frac{18}{\sigma _b^2}}+{\displaystyle \frac{30}{\sigma _b^3}}+{\displaystyle \frac{24}{\sigma _b^4}})`$ (166)
$`+\mathrm{\hspace{0.17em}10}\delta ^{i(k}\delta ^{l)j}\left({\displaystyle \frac{1}{\sigma _b^2}}+{\displaystyle \frac{5}{\sigma _b^3}}+{\displaystyle \frac{8}{\sigma _b^4}}\right)`$ (167)
$`+{\displaystyle \frac{1}{|𝐱|^2}}(\delta ^{ij}x^kx^l+\delta ^{kl}x^ix^j3\delta ^{i(k}x^{l)}x^j3\delta ^{j(k}x^{l)}x^i)\left(1+{\displaystyle \frac{5}{\sigma _b}}+{\displaystyle \frac{6}{\sigma _b^2}}{\displaystyle \frac{18}{\sigma _b^3}}{\displaystyle \frac{48}{\sigma _b^4}}\right)`$ (168)
$`{\displaystyle \frac{10}{|𝐱|^2}}(\delta ^{i(k}x^{l)}x^j+\delta ^{j(k}x^{l)}x^i)\left({\displaystyle \frac{1}{\sigma _b}}+{\displaystyle \frac{9}{\sigma _b^2}}+{\displaystyle \frac{33}{\sigma _b^3}}+{\displaystyle \frac{48}{\sigma _b^4}}\right)`$ (169)
$`+{\displaystyle \frac{2}{|𝐱|^4}}x^ix^jx^kx^l(1+{\displaystyle \frac{14}{\sigma _b}}+{\displaystyle \frac{87}{\sigma _b^2}}+{\displaystyle \frac{279}{\sigma _b^3}}+{\displaystyle \frac{384}{\sigma _b^4}})],`$ (170)
$`𝒢_2^{00ij}(𝐱){\displaystyle \frac{1}{3\pi }}{\displaystyle \frac{a\gamma ^6}{L_P^4}}{\displaystyle \frac{K_4(\sigma _\gamma )}{\sigma _\gamma ^2}}\left(3{\displaystyle \frac{x^ix^j}{|𝐱|^2}}\delta ^{ij}\right),`$ (171)
$`𝒢_2^{0i0j}(𝐱){\displaystyle \frac{1}{6\pi }}{\displaystyle \frac{a\gamma ^6}{L_P^4}}{\displaystyle \frac{K_4(\sigma _\gamma )}{\sigma _\gamma ^2}}\left({\displaystyle \frac{x^ix^j}{|𝐱|^2}}+3\delta ^{ij}\right){\displaystyle \frac{5}{3\pi }}{\displaystyle \frac{a\gamma ^6}{L_P^4}}{\displaystyle \frac{K_3(\sigma _\gamma )}{\sigma _\gamma ^3}}\delta ^{ij},`$ (172)
$`𝒢_2^{ijkl}(𝐱){\displaystyle \frac{1}{3\pi }}{\displaystyle \frac{a\gamma ^6}{L_P^4}}{\displaystyle \frac{K_4(\sigma _\gamma )}{\sigma _\gamma ^2}}[2(\delta ^{ij}\delta ^{kl}3\delta ^{i(k}\delta ^{l)j}){\displaystyle \frac{1}{|𝐱|^2}}(\delta ^{ij}x^kx^l+\delta ^{kl}x^ix^j3\delta ^{i(k}x^{l)}x^j`$ (173)
$`3\delta ^{j(k}x^{l)}x^i)]{\displaystyle \frac{4}{3\pi }}{\displaystyle \frac{a\gamma ^6}{L_P^4}}{\displaystyle \frac{K_3(\sigma _\gamma )}{\sigma _\gamma ^3}}(\delta ^{ij}\delta ^{kl}3\delta ^{i(k}\delta ^{l)j}),`$ (174)
$`𝒢_3^{ijkl}(𝐱){\displaystyle \frac{1}{\pi }}{\displaystyle \frac{a\delta ^6}{L_P^4}}{\displaystyle \frac{K_3(\sigma _\delta )}{\sigma _\delta ^3}}\left(\delta ^{ij}\delta ^{kl}3\delta ^{i(k}\delta ^{l)j}\right).`$ (175)
Note that, for $`\sigma 1`$, we have the following asymptotic expansions for the modified Bessel functions in the above expressions:
$$K_4(\sigma )\sqrt{\frac{\pi }{2\sigma }}e^\sigma \left[1+\frac{63}{8}\frac{1}{\sigma }+O\left(\frac{1}{\sigma ^2}\right)\right],K_3(\sigma )\sqrt{\frac{\pi }{2\sigma }}e^\sigma \left[1+\frac{35}{8}\frac{1}{\sigma }+O\left(\frac{1}{\sigma ^2}\right)\right].$$
(176)
### C Correlation functions for the metric perturbations
Starting from the solutions found for the linearized Einstein tensor, which are characterized by the two-point correlation functions (143) \[or, in terms of Fourier transforms, (142)\], we can now solve the equations for the metric perturbations. Working in the harmonic gauge, $`_\nu \overline{h}^{\mu \nu }=0`$ (this zero must be understood in the same statistical sense as above), where $`\overline{h}_{\mu \nu }h_{\mu \nu }(1/2)\eta _{\mu \nu }h`$, and using Eqs. (22) and (E1), these equations reduce to $`\mathrm{}\overline{h}^{\mu \nu }(x)=2G^{\left(1\right)\mu \nu }(x)`$, or, in terms of Fourier transforms, $`p^2\stackrel{~}{\overline{h}}^{_{_{_{_{_{\mu \nu }}}}}}(p)=2\stackrel{~}{G}^{\left(1\right)\mu \nu }(p)`$. As above, we can write $`\overline{h}^{\mu \nu }=\overline{h}^{\mu \nu }_c+\overline{h}_\mathrm{f}^{\mu \nu }`$, where $`\overline{h}_\mathrm{f}^{\mu \nu }`$ is a solution to these equations with zero average, and the two-point correlation functions are given by
$$^{\mu \nu \alpha \beta }(x,x^{})\overline{h}^{\mu \nu }(x)\overline{h}^{\alpha \beta }(x^{})_c\overline{h}^{\mu \nu }(x)_c\overline{h}^{\alpha \beta }(x^{})_c=\overline{h}_\mathrm{f}^{\mu \nu }(x)\overline{h}_\mathrm{f}^{\alpha \beta }(x^{})_c.$$
(177)
We can now seek solutions of the form $`\stackrel{~}{\overline{h}}_\mathrm{f}^{_{_{_{_{_{\mu \nu }}}}}}(p)=2D(p)\stackrel{~}{G}_\mathrm{f}^{\left(1\right)\mu \nu }(p)`$, where $`D(p)`$ is a Lorentz invariant scalar distribution in Minkowski spacetime, which is the most general solution of $`p^2D(p)=1`$. Note that, since the linearized Einstein tensor is conserved, solutions of this form automatically satisfy the harmonic gauge condition. As above, we can write $`D(p)=[1/p^2]_r+D_h(p)`$, where $`D_h(p)`$ is the most general solution to the associated homogeneous equation and, correspondingly, we have $`\stackrel{~}{\overline{h}}_\mathrm{f}^{_{_{_{_{_{\mu \nu }}}}}}(p)=\stackrel{~}{\overline{h}}_p^{_{_{_{_{_{\mu \nu }}}}}}(p)+\stackrel{~}{\overline{h}}_h^{_{_{_{_{_{\mu \nu }}}}}}(p)`$. However, since $`D_h(p)`$ has support on the set of points for which $`p^2=0`$, it is easy to see from Eq. (142) \[from the factor $`\theta (p^24m^2)`$\] that $`\stackrel{~}{\overline{h}}_h^{_{_{_{_{_{\mu \nu }}}}}}(p)\stackrel{~}{G}_\mathrm{f}^{\left(1\right)\alpha \beta }(p^{})_c=0`$ and, thus, the two-point correlation functions (177) can be computed from $`\stackrel{~}{\overline{h}}_\mathrm{f}^{_{_{_{_{_{\mu \nu }}}}}}(p)\stackrel{~}{\overline{h}}_\mathrm{f}^{_{_{_{_{_{\alpha \beta }}}}}}(p^{})_c=\stackrel{~}{\overline{h}}_p^{_{_{_{_{_{\mu \nu }}}}}}(p)\stackrel{~}{\overline{h}}_p^{_{_{_{_{_{\alpha \beta }}}}}}(p^{})_c`$. From Eq. (142) and due to the factor $`\theta (p^24m^2)`$, it is also easy to see that the prescription $`[]_r`$ is irrelevant in this correlation function and we obtain
$$\stackrel{~}{\overline{h}}_p^{_{_{_{_{_{\mu \nu }}}}}}(p)\stackrel{~}{\overline{h}}_p^{_{_{_{_{_{\alpha \beta }}}}}}(p^{})_c=\frac{4}{(p^2)^2}\stackrel{~}{G}_p^{\left(1\right)\mu \nu }(p)\stackrel{~}{G}_p^{\left(1\right)\alpha \beta }(p^{})_c,$$
(178)
where $`\stackrel{~}{G}_p^{\left(1\right)\mu \nu }(p)\stackrel{~}{G}_p^{\left(1\right)\alpha \beta }(p^{})_c`$ is given in (142). The right hand side of this equation is a well defined bi-distribution, at least for $`m0`$ (the $`\theta `$ function provides the suitable cutoff). In the massless field case, since the noise kernel is obtained as the limit $`m0`$ of the noise kernel for a massive field, it seems that the natural prescription to avoid the divergencies on the lightcone $`p^2=0`$ is a Hadamard finite part (see Refs. for its definition). Taking this prescription, we also get a well defined bi-distribution for the massless limit of the last expression. Finally, we find the result
$$^{\mu \nu \alpha \beta }(x,x^{})=\frac{4\pi }{45}G^2_x^{\mu \nu \alpha \beta }_\mathrm{A}(xx^{})+\frac{32\pi }{9}G^2_x^{\mu \nu }_x^{\alpha \beta }_\mathrm{B}(xx^{}),$$
(179)
where $`\stackrel{~}{}_\mathrm{A}(p)[1/(p^2)^2]\stackrel{~}{𝒢}_\mathrm{A}(p)`$ and $`\stackrel{~}{}_\mathrm{B}(p)[1/(p^2)^2]\stackrel{~}{𝒢}_\mathrm{B}(p)`$, with $`\stackrel{~}{𝒢}_\mathrm{A}(p)`$ and $`\stackrel{~}{𝒢}_\mathrm{B}(p)`$ given by (145). The two-point correlation functions for the metric perturbations can be easily obtained using $`h_{\mu \nu }=\overline{h}_{\mu \nu }(1/2)\eta _{\mu \nu }\overline{h}_\alpha ^\alpha `$.
## VI Discussion
Our main results for the correlation functions are (143) and (179). In the case of a conformal field, the correlation functions of the linearized Einstein tensor have been explicitly evaluated and the results are given in (175). From the exponential factors $`e^\sigma `$ in these results, we see that the correlation functions of the linearized Einstein tensor are in this case characterized by correlation lengths of the order of the Planck length. A similar behavior is expected for the correlation functions of the metric perturbations. Hence, as expected in this case, the correlation functions are negligibly small for points separated by distances large compared to the Planck length. At such scales, the dynamics of gravitational perturbations of Minkowski spacetime can be simply described by semiclassical gravity . Deviations from semiclassical gravity are only important for points separated by Planckian or sub-Planckian scales. However, for such scales, our results (175) are not reliable, since we expect that gravitational fluctuations of genuine quantum nature to be relevant and, thus, the classical description breaks down. It is interesting to note, however, that these results for correlation functions are non-analytic in their characteristic correlation lengths. This kind of non-analytic behavior is actually quite typical of the solutions of Langevin-type equations with dissipative terms. An example in the context of a reduced version of the semiclassical Einstein-Langevin equation is given in Ref. ).
For background solutions of semiclassical gravity with other scales present apart from the Planck scales (for instance, for matter fields in a thermal state), stress-energy fluctuations may be important at larger scales. For such backgrounds, stochastic semiclassical gravity might predict correlation functions with characteristic correlation lengths much larger than the Planck scales, so as to be relevant and reliable on a certain range of scales. It seems quite plausible, nevertheless, that these correlation functions would remain non-analytic in their characteristic correlation lengths. This would imply that these correlation functions could not be obtained from a calculation involving a perturbative expansion in the characteristic correlation lengths. In particular, if these correlation lengths are proportional to the Planck constant $`\mathrm{}`$, the gravitational correlation functions could not be obtained from an expansion in $`\mathrm{}`$. Hence, stochastic semiclassical gravity might predict a behavior for gravitational correlation functions different from that of the analogous functions in perturbative quantum gravity . This is not necessarily inconsistent with having neglected action terms of higher order in $`\mathrm{}`$ when considering semiclassical gravity as an effective theory .
We conclude this section with some comments about a technical point on the obtained solutions of stochastic semiclassical gravity. It concerns the issue that the Einstein-Langevin equations, as well as the semiclassical Einstein equations, contain derivatives of order higher than two. Due to this fact, these equations can have some “pathological” solutions (e.g., “runaway” solutions) which are presumably unphysical . Thus, one needs to apply some criterion to discern the “physical” from the unphysical solutions. However, as it is discussed in Ref. (see also Refs. ), even in the context of “pure” (non-stochastic) semiclassical gravity, this is still an open problem. Two main proposals, both based in the works by Simon , have been made concerning this issue: the “perturbative expandability” (in $`\mathrm{}`$) criterion and the “reduction of order” procedure .
The first proposal consists in identifying a subclass of “physical” solutions which are analytic in the Planck constant $`\mathrm{}`$. This proposal has been successful in eliminating the instability of Minkowski spacetime found by Horowitz . However, on the one hand, this proposal seems to be too restrictive since, as it has been pointed out in Ref. , one could not describe effects such as the continuous mass loss of a black hole due to Hawking radiation. On the other hand, there can be situations in which the formal series obtained when seeking approximate perturbative solutions (to a finite order in $`\mathrm{}`$) does not converge to a solution to the semiclassical equations . In our case, if we had tried to find solutions to Eq. (117) as a Taylor expansion in $`\mathrm{}`$, we would have obtained a series for $`\stackrel{~}{G}_{\mu \nu }^{\left(1\right)}(p)`$ which, as the above solutions, would be linear and local in $`\stackrel{~}{\xi }_{\alpha \beta }(p)`$, but whose corresponding two-point correlation functions for the conformal field case would not converge to (143).
The “reduction of order” procedure provides in some cases a reasonable way to modify the semiclassical equations in order to eliminate spurious solutions. But, as it has been emphasized in Ref. , it is not clear at all whether a reduction of order procedure can always be applied to the semiclassical Einstein equation (and how this procedure should be applied). For the Einstein-Langevin equation, this issue has not been, to our knowledge, properly addressed. A naive application of the prescription to Eq. (117) seems to downplay the role of the dissipative terms with respect to the noise source. In fact, to lowest order, we obtain $`G^{\left(1\right)\mu \nu }=16\pi G\xi ^{\mu \nu }`$, where there is no contribution of the dissipation kernel. From this equation, we get the well-known result $`G^{\left(1\right)\mu \nu }_c=0`$ , and also $`𝒢^{\mu \nu \alpha \beta }(x,x^{})=(16\pi )^2L_P^4N^{\mu \nu \alpha \beta }(x,x^{})`$. For a massless field, using Eqs. (40), (89) and (93), this gives $`𝒢^{\mu \nu \alpha \beta }(x,x^{})=(2/15)(L_P^4/\pi ^2)\left[(1/6)_x^{\mu \nu \alpha \beta }+60\mathrm{\Delta }\xi ^2_x^{\mu \nu }_x^{\alpha \beta }\right]\left[𝒫f\left[1/\left((xx^{})^2\right)^2\right]+\pi ^3\delta ^4(xx^{})\right]`$. For the two-point correlation functions (177), we get, in the harmonic gauge, $`^{\mu \nu \alpha \beta }(x,x^{})=(4\pi /45)L_P^4_x^{\mu \nu \alpha \beta }_\mathrm{A}(xx^{})+(32\pi /9)L_P^4_x^{\mu \nu }_x^{\alpha \beta }_\mathrm{B}(xx^{})`$, with $`\stackrel{~}{}_\mathrm{A}(p)\theta (p^24m^2)(p^2)^2\sqrt{1+4m^2/p^2}\left(1+4m^2/p^2\right)^2`$ and $`\stackrel{~}{}_\mathrm{B}(p)\theta (p^24m^2)(p^2)^2\sqrt{1+4m^2/p^2}\left(3\mathrm{\Delta }\xi +m^2/p^2\right)^2`$. Comparing the last results for the massless case with the ones obtained in Sec. V, we note that the main qualitative feature is the absence of the exponential factors $`e^\sigma `$, which make the two-point correlation functions to decay much more slowly with the distance, i.e., like a power instead of an exponential law. This fact is due to the lack of dissipative terms in the reduced order equations. The conclusion is that one should probably implement a more sophisticated version of the reduction of order procedure so as to keep some contribution of the dissipation kernel in the reduced order equations.
For these reasons, in our work we have not attempted any of these procedures and we have simply sought some solutions to the full equations (117). Our solutions for the conformal field case have the physically reasonable feature of having negligible two-point functions for points separated by scales larger than the Planck length.
###### Acknowledgements.
We are grateful to Esteban Calzetta, Jaume Garriga, Bei-Lok Hu, Ted Jacobson and Albert Roura for very helpful suggestions and discussions. This work has been partially supported by the CICYT Research Project number AEN98-0431, and the European Project number CI1-CT94-0004.
## A The kernels for a vacuum state
The kernels for a vacuum state can be computed in terms of the Wightman and Feynman functions defined in Eq. (6) using
$`0|\widehat{t}_n^{ab}(x)\widehat{t}_n^{cd}(y)|0=4(N_n^{abcd}(x,y)+iH_{\mathrm{A}_n}^{abcd}(x,y))=__x^a__y^cG_n^+(x,y)__x^b__y^dG_n^+(x,y)`$ (A1)
$`+__x^a__y^dG_n^+(x,y)__x^b__y^cG_n^+(x,y)+2𝒟_x^{ab}\left(__y^cG_n^+(x,y)__y^dG_n^+(x,y)\right)`$ (A2)
$`+\mathrm{\hspace{0.17em}2}𝒟_y^{cd}\left(__x^aG_n^+(x,y)__x^bG_n^+(x,y)\right)+2𝒟_x^{ab}𝒟_y^{cd}(G_n^{+2}(x,y)),`$ (A3)
where $`𝒟^{ab}`$ is the differential operator
$$𝒟_x^{ab}(\xi \frac{1}{4})g^{ab}(x)\mathrm{}_x+\xi (R^{ab}(x)_x^a_x^b),$$
(A4)
and
$`H_{\mathrm{S}_n}^{abcd}(x,y)={\displaystyle \frac{1}{4}}\mathrm{Im}[__x^a__y^cG_{F_n}(x,y)__x^b__y^dG_{F_n}(x,y)+__x^a__y^dG_{F_n}(x,y)__x^b__y^cG_{F_n}(x,y)`$ (A5)
$`+\mathrm{\hspace{0.17em}2}𝒟_x^{ab}\left(__y^cG_{F_n}(x,y)__y^dG_{F_n}(x,y)\right)+2𝒟_y^{cd}\left(__x^aG_{F_n}(x,y)__x^bG_{F_n}(x,y)\right)`$ (A6)
$`+2𝒟_x^{ab}𝒟_y^{cd}\left(G_{F_n}^{\mathrm{\hspace{0.33em}\hspace{0.25em}2}}(x,y)\right)+{\displaystyle \frac{1}{2}}[g^{ab}(x)(__y^cG_{F_n}(x,y)__y^d+__y^dG_{F_n}(x,y)__y^c)`$ (A7)
$`+g^{cd}(y)(__x^aG_{F_n}(x,y)__x^b+__x^bG_{F_n}(x,y)__x^a)]{\displaystyle \frac{\delta ^n(xy)}{\sqrt{g(x)}}}+(g^{ab}(x)𝒟_y^{cd}`$ (A8)
$`+g^{cd}(y)𝒟_x^{ab})\left({\displaystyle \frac{\delta ^n(xy)}{\sqrt{g(x)}}}G_{F_n}(x,y)\right)+{\displaystyle \frac{1}{4}}g^{ab}(x)g^{cd}(y)G_{F_n}(x,y)(\mathrm{}_xm^2\xi R(x)){\displaystyle \frac{\delta ^n(xy)}{\sqrt{g(x)}}}].`$ (A9)
## B Momentum integrals
Some useful expressions for the momentum integrals in dimensional regularization defined in (47) and (49) are:
$`I_{0_n}={\displaystyle \frac{i}{(4\pi )^2}}m^2\left({\displaystyle \frac{m^2}{4\pi \mu ^2}}\right)^{_{\frac{n4}{2}}}\mathrm{\Gamma }\left(1{\displaystyle \frac{n}{2}}\right)={\displaystyle \frac{i}{(4\pi )^2}}{\displaystyle \frac{4m^2}{(n2)}}\kappa _n+O(n4),`$ (B1)
$`I_{0_n}^\mu =0,`$ (B2)
$`I_{0_n}^{\mu \nu }=m^2\eta ^{\mu \nu }{\displaystyle \frac{I_{0_n}}{n}},`$ (B3)
$`J_n(p)={\displaystyle \frac{i}{(4\pi )^2}}\left[2\kappa _n+\varphi (p^2)+O(n4)\right],`$ (B4)
$`J_n^\mu (p)={\displaystyle \frac{J_n(p)}{2}}p^\mu ,`$ (B5)
$`J_n^{\mu \nu }(p)={\displaystyle \frac{J_n(p)}{4}}\left[p^\mu p^\nu \left(1+4{\displaystyle \frac{m^2}{p^2}}\right){\displaystyle \frac{p^2P^{\mu \nu }}{(n1)}}\right]+{\displaystyle \frac{I_{0_n}}{2}}{\displaystyle \frac{1}{p^2}}[p^\mu p^\nu +{\displaystyle \frac{p^2P^{\mu \nu }}{n1}}],`$ (B6)
$`J_n^{\mu \nu \alpha }(p)={\displaystyle \frac{J_n(p)}{8}}\left[p^\mu p^\nu p^\alpha \left(1+4{\displaystyle \frac{m^2}{p^2}}\right){\displaystyle \frac{p^2}{(n1)}}\left(P^{\mu \nu }p^\alpha +P^{\mu \alpha }p^\nu +P^{\alpha \nu }p^\mu \right)\right]`$ (B7)
$`+{\displaystyle \frac{I_{0_n}}{4}}{\displaystyle \frac{1}{p^2}}[3p^\mu p^\nu p^\alpha +{\displaystyle \frac{p^2}{(n1)}}\left(P^{\mu \nu }p^\alpha +P^{\mu \alpha }p^\nu +P^{\alpha \nu }p^\mu \right)],`$ (B8)
$`J_n^{\mu \nu \alpha \beta }(p)={\displaystyle \frac{J_n(p)}{16}}[p^\mu p^\nu p^\alpha p^\beta (1+4{\displaystyle \frac{m^2}{p^2}}){\displaystyle \frac{p^2}{(n1)}}(P^{\mu \nu }p^\alpha p^\beta +P^{\nu \alpha }p^\mu p^\beta +P^{\nu \beta }p^\mu p^\alpha `$ (B9)
$`+P^{\mu \alpha }p^\nu p^\beta +P^{\mu \beta }p^\nu p^\alpha +P^{\alpha \beta }p^\mu p^\nu )+(1+4{\displaystyle \frac{m^2}{p^2}}\left)^2{\displaystyle \frac{(p^2)^2}{(n^21)}}\right(P^{\mu \nu }P^{\alpha \beta }`$ (B10)
$`+P^{\mu \alpha }P^{\nu \beta }+P^{\mu \beta }P^{\nu \alpha })]`$ (B11)
$`+{\displaystyle \frac{I_{0_n}}{8}}{\displaystyle \frac{1}{p^2}}[(7{\displaystyle \frac{12}{n}}{\displaystyle \frac{m^2}{p^2}})p^\mu p^\nu p^\alpha p^\beta +({\displaystyle \frac{1}{n1}}{\displaystyle \frac{4}{n}}{\displaystyle \frac{m^2}{p^2}})p^2(P^{\mu \nu }p^\alpha p^\beta +P^{\nu \alpha }p^\mu p^\beta `$ (B12)
$`+P^{\nu \beta }p^\mu p^\alpha +P^{\mu \alpha }p^\nu p^\beta +P^{\mu \beta }p^\nu p^\alpha +P^{\alpha \beta }p^\mu p^\nu )`$ (B13)
$`({\displaystyle \frac{1}{n^21}}{\displaystyle \frac{4(2n1)}{n(n^21)}}{\displaystyle \frac{m^2}{p^2}})(p^2)^2(P^{\mu \nu }P^{\alpha \beta }+P^{\mu \alpha }P^{\nu \beta }+P^{\mu \beta }P^{\nu \alpha })],`$ (B14)
where $`p^2P^{\mu \nu }\eta ^{\mu \nu }p^2p^\mu p^\nu `$, $`\kappa _n`$ is defined in (18),
$$\varphi (p^2)_0^1𝑑\alpha \mathrm{ln}\left(1+\frac{p^2}{m^2}\alpha (1\alpha )iϵ\right)=i\pi \theta (p^24m^2)\sqrt{1+4\frac{m^2}{p^2}}+\phi (p^2),$$
(B16)
with $`ϵ0^+`$, and
$`\phi (p^2)`$ $`{\displaystyle _0^1}𝑑\alpha \mathrm{ln}\left|1+{\displaystyle \frac{p^2}{m^2}}\alpha (1\alpha )\right|=2+\sqrt{1+4{\displaystyle \frac{m^2}{p^2}}}\mathrm{ln}\left|{\displaystyle \frac{\sqrt{1+4\frac{m^2}{p^2}}+1}{\sqrt{1+4\frac{m^2}{p^2}}1}}\right|\theta \left(1+4{\displaystyle \frac{m^2}{p^2}}\right)`$ (B18)
$`+\mathrm{\hspace{0.17em}2}\sqrt{14{\displaystyle \frac{m^2}{p^2}}}\mathrm{arccotan}\left(\sqrt{14{\displaystyle \frac{m^2}{p^2}}}\right)\theta \left(14{\displaystyle \frac{m^2}{p^2}}\right).`$
We can also write $`\varphi (p^2)`$ in a more compact way as
$$\varphi (p^2)=2+\sqrt{1+4\frac{m^2}{p^2}}\mathrm{ln}\left(\frac{\sqrt{1+4(m^2iϵ)/p^2}+1}{\sqrt{1+4(m^2iϵ)/p^2}1}\right).$$
(B19)
Other useful integrals in momentum space defined in (29) are
$`I(p)={\displaystyle \frac{1}{4(2\pi )^3}}\theta (p^0)\theta (p^24m^2)\sqrt{1+4{\displaystyle \frac{m^2}{p^2}}},`$ (B20)
$`I^\mu (p)={\displaystyle \frac{I(p)}{2}}p^\mu ,`$ (B21)
$`I^{\mu \nu }(p)={\displaystyle \frac{I(p)}{4}}[p^\mu p^\nu \left(1+4{\displaystyle \frac{m^2}{p^2}}\right){\displaystyle \frac{p^2P^{\mu \nu }}{3}}],`$ (B22)
$`I^{\mu \nu \alpha }(p)={\displaystyle \frac{I(p)}{8}}[p^\mu p^\nu p^\alpha \left(1+4{\displaystyle \frac{m^2}{p^2}}\right){\displaystyle \frac{p^2}{3}}\left(P^{\mu \nu }p^\alpha +P^{\mu \alpha }p^\nu +P^{\alpha \nu }p^\mu \right)],`$ (B23)
$`I^{\mu \nu \alpha \beta }(p)={\displaystyle \frac{I(p)}{16}}[p^\mu p^\nu p^\alpha p^\beta (1+4{\displaystyle \frac{m^2}{p^2}}){\displaystyle \frac{p^2}{3}}(P^{\mu \nu }p^\alpha p^\beta +P^{\nu \alpha }p^\mu p^\beta +P^{\nu \beta }p^\mu p^\alpha +P^{\mu \alpha }p^\nu p^\beta `$ (B24)
$`+P^{\mu \beta }p^\nu p^\alpha +P^{\alpha \beta }p^\mu p^\nu )+(1+4{\displaystyle \frac{m^2}{p^2}})^2{\displaystyle \frac{(p^2)^2}{15}}(P^{\mu \nu }P^{\alpha \beta }+P^{\mu \alpha }P^{\nu \beta }+P^{\mu \beta }P^{\nu \alpha })],`$ (B25)
## C Products of Wightman functions
For the products of derivatives of Wightman functions involved in the calculations of Sec. III A, we obtain the following expressions:
$`\mathrm{\Delta }^{+2}(x)=(2\pi )^2{\displaystyle \frac{d^4p}{(2\pi )^4}e^{ipx}I(p)},`$ (C1)
$`^\mu \mathrm{\Delta }^+(x)^\nu \mathrm{\Delta }^+(x)=(2\pi )^2{\displaystyle \frac{d^4p}{(2\pi )^4}e^{ipx}\left[I^\mu (p)p^\nu I^{\mu \nu }(p)\right]},`$ (C2)
$`^\mu ^\nu \mathrm{\Delta }^+(x)^\alpha ^\beta \mathrm{\Delta }^+(x)=(2\pi )^2{\displaystyle \frac{d^4p}{(2\pi )^4}e^{ipx}\left[I^{\mu \nu }(p)p^\alpha p^\beta 2I^{\mu \nu (\alpha }(p)p^{\beta )}+I^{\mu \nu \alpha \beta }(p)\right]},`$ (C3)
with $`I(p)`$, $`I^\mu (p)`$, $`I^{\mu \nu }(p)`$, $`I^{\mu \nu \alpha }(p)`$ and $`I^{\mu \nu \alpha \beta }(p)`$ given by Eqs. (B20)-(LABEL:I\_mu-nu-alpha-beta(P)). From these expressions, using the results of Appendix B, we obtain
$`^\mu \mathrm{\Delta }^+(x)^\nu \mathrm{\Delta }^+(x)=\pi ^2_x^\mu _x^\nu {\displaystyle \frac{d^4p}{(2\pi )^4}e^{ipx}I(p)}{\displaystyle \frac{\pi ^2}{3}}_x^{\mu \nu }{\displaystyle \frac{d^4p}{(2\pi )^4}e^{ipx}\left(1+4\frac{m^2}{p^2}\right)I(p)},`$ (C5)
(C6)
$`^\mu ^{(\alpha }\mathrm{\Delta }^+(x)^{\beta )}^\nu \mathrm{\Delta }^+(x)={\displaystyle \frac{\pi ^2}{4}}_x^\mu _x^\nu _x^\alpha _x^\beta {\displaystyle \frac{d^4p}{(2\pi )^4}e^{ipx}I(p)}`$ (C7)
$`{\displaystyle \frac{\pi ^2}{12}}(_x^{\mu \nu }_x^\alpha _x^\beta +_x^{\alpha \beta }_x^\mu _x^\nu ){\displaystyle \frac{d^4p}{(2\pi )^4}e^{ipx}\left(1+4\frac{m^2}{p^2}\right)I(p)}`$ (C8)
$`{\displaystyle \frac{\pi ^2}{60}}(_x^{\mu \nu }_x^{\alpha \beta }+2_x^{\mu (\alpha }_x^{\beta )\nu }){\displaystyle \frac{d^4p}{(2\pi )^4}e^{ipx}\left(1+4\frac{m^2}{p^2}\right)^2I(p)}.`$ (C9)
## D Products of Feynman functions
For the products of derivatives of Feynman functions that we need for the calculations of Sec. III B, we obtain the following results:
$`\mu ^{(n4)}\mathrm{\Delta }_{F_n}^2(x)={\displaystyle \frac{d^np}{(2\pi )^n}e^{ipx}J_n(p)},`$ (D1)
$`\mu ^{(n4)}^\mu \mathrm{\Delta }_{F_n}(x)^\nu \mathrm{\Delta }_{F_n}(x)={\displaystyle \frac{d^np}{(2\pi )^n}e^{ipx}\left[J_n^\mu (p)p^\nu J_n^{\mu \nu }(p)\right]},`$ (D2)
$`\mu ^{(n4)}^\mu ^\nu \mathrm{\Delta }_{F_n}(x)^\alpha ^\beta \mathrm{\Delta }_{F_n}(x)={\displaystyle \frac{d^np}{(2\pi )^n}e^{ipx}\left[J_n^{\mu \nu }(p)p^\alpha p^\beta 2J_n^{\mu \nu (\alpha }(p)p^{\beta )}+J_n^{\mu \nu \alpha \beta }(p)\right]},`$ (D3)
(D4)
$`\mu ^{(n4)}\mathrm{\Delta }_{F_n}(0)=I_{0_n},`$ (D5)
$`\mu ^{(n4)}^\mu \mathrm{\Delta }_{F_n}(x)^\nu \delta ^n(x)={\displaystyle \frac{d^np}{(2\pi )^n}e^{ipx}\left(I_{0_n}^\mu p^\nu I_{0_n}^{\mu \nu }\right)},`$ (D6)
$`\mu ^{(n4)}\mathrm{\Delta }_{F_n}(x)\mathrm{}\delta ^n(x)={\displaystyle \frac{d^np}{(2\pi )^n}e^{ipx}\left(p^2I_{0_n}+2p_\mu I_{0_n}^\mu +I_{0_n}^\mu _\mu \right)}.`$ (D7)
Using the results of Appendix B, we find from the above expressions
$`\mu ^{(n4)}^\mu \mathrm{\Delta }_{F_n}(x)^\nu \mathrm{\Delta }_{F_n}(x)={\displaystyle \frac{1}{4}}_x^\mu _x^\nu {\displaystyle \frac{d^np}{(2\pi )^n}e^{ipx}J_n(p)}+{\displaystyle \frac{1}{12}}_x^{\mu \nu }{\displaystyle \frac{d^np}{(2\pi )^n}e^{ipx}(1+4\frac{m^2}{p^2})J_n(p)}`$ (D8)
$`+{\displaystyle \frac{1}{2}}{\displaystyle \frac{d^np}{(2\pi )^n}e^{ipx}\left[I_{0_n}\left(\frac{p^\mu p^\nu }{p^2}+\frac{1}{3}P^{\mu \nu }\right)\frac{i}{(4\pi )^2}\frac{1}{9}(p^2+6m^2)P^{\mu \nu }\right]}+O(n4),`$ (D9)
(D10)
$`\mu ^{(n4)}^\mu ^{(\alpha }\mathrm{\Delta }_{F_n}(x)^{\beta )}^\nu \mathrm{\Delta }_{F_n}(x)={\displaystyle \frac{1}{16}}_x^\mu _x^\nu _x^\alpha _x^\beta {\displaystyle \frac{d^np}{(2\pi )^n}e^{ipx}J_n(p)}`$ (D11)
$`+{\displaystyle \frac{1}{48}}(_x^{\mu \nu }_x^\alpha _x^\beta +_x^{\alpha \beta }_x^\mu _x^\nu ){\displaystyle \frac{d^np}{(2\pi )^n}e^{ipx}(1+4\frac{m^2}{p^2})J_n(p)}`$ (D12)
$`+{\displaystyle \frac{1}{240}}(_x^{\mu \nu }_x^{\alpha \beta }+2_x^{\mu (\alpha }_x^{\beta )\nu }){\displaystyle \frac{d^np}{(2\pi )^n}e^{ipx}\left(1+4\frac{m^2}{p^2}\right)^2J_n(p)}`$ (D13)
$`{\displaystyle \frac{1}{8}}{\displaystyle }{\displaystyle \frac{d^np}{(2\pi )^n}}e^{ipx}\{I_{0_n}[{\displaystyle \frac{1}{p^2}}(1+{\displaystyle \frac{12}{n}}{\displaystyle \frac{m^2}{p^2}})p^\mu p^\nu p^\alpha p^\beta +{\displaystyle \frac{1}{3}}(P^{\mu \nu }p^\alpha p^\beta +P^{\alpha \beta }p^\mu p^\nu )`$ (D14)
$`+{\displaystyle \frac{4}{n}}{\displaystyle \frac{m^2}{p^2}}\left(P^{\mu \nu }p^\alpha p^\beta +P^{\alpha \beta }p^\mu p^\nu +2P^{\mu (\alpha }p^{\beta )}p^\nu +2P^{\nu (\alpha }p^{\beta )}p^\mu \right)`$ (D15)
$`+{\displaystyle \frac{1}{15}}(p^2+{\displaystyle \frac{28}{n}}m^2)(P^{\mu \nu }P^{\alpha \beta }+2P^{\mu (\alpha }P^{\beta )\nu })]{\displaystyle \frac{i}{(4\pi )^2}}{\displaystyle \frac{1}{9}}(p^2+6m^2)(P^{\mu \nu }p^\alpha p^\beta +P^{\alpha \beta }p^\mu p^\nu )`$ (D16)
$`{\displaystyle \frac{i}{4\pi ^2}}{\displaystyle \frac{1}{225}}(2(p^2)^2+20m^2p^2+45m^4)(P^{\mu \nu }P^{\alpha \beta }+2P^{\mu (\alpha }P^{\beta )\nu })\}+O(n4),`$ (D17)
where $`P^{\mu \nu }`$ is the projector orthogonal to $`p^\mu `$ defined above.
## E Linearized tensors around flat spacetime
Some curvature tensors linearized around flat spacetime are given by the following expressions:
$`G^{\left(1\right)\mu \nu }`$ $`=`$ $`R^{\left(1\right)\mu \nu }{\displaystyle \frac{1}{2}}\eta ^{\mu \nu }R^{\left(1\right)},`$ (E1)
$`D^{\left(1\right)\mu \nu }`$ $`=`$ $`^\mu ^\nu R^{\left(1\right)}+{\displaystyle \frac{1}{2}}\eta ^{\mu \nu }\mathrm{}R^{\left(1\right)}3\mathrm{}R^{\left(1\right)\mu \nu },`$ (E2)
$`B^{\left(1\right)\mu \nu }`$ $`=`$ $`2(^\mu ^\nu R^{\left(1\right)}\eta ^{\mu \nu }\mathrm{}R^{\left(1\right)}),`$ (E3)
with
$`R^{\left(1\right)\mu \nu }={\displaystyle \frac{1}{2}}(_\alpha ^\mu h^{\nu \alpha }+_\alpha ^\nu h^{\mu \alpha }\mathrm{}h^{\mu \nu }^\mu ^\nu h),`$ (E4)
$`R^{\left(1\right)}=\eta _{\alpha \beta }R^{\left(1\right)\alpha \beta }=^\alpha ^\beta h_{\alpha \beta }\mathrm{}h,`$ (E5)
and
$$R^{\left(1\right)\mu \nu \alpha \beta }=\frac{1}{2}\left(^\mu ^\beta h^{\nu \alpha }+^\nu ^\alpha h^{\mu \beta }^\mu ^\alpha h^{\nu \beta }^\nu ^\beta h^{\mu \alpha }\right).$$
(E6)
In four spacetime dimensions, the linearized Weyl tensor is given by
$`C^{\left(1\right)\mu \nu \alpha \beta }={\displaystyle \frac{1}{12}}[6(\eta ^{\nu \rho }\eta ^{\alpha \sigma }^\mu ^\beta +\eta ^{\mu \rho }\eta ^{\beta \sigma }^\nu ^\alpha \eta ^{\nu \rho }\eta ^{\beta \sigma }^\mu ^\alpha \eta ^{\mu \rho }\eta ^{\alpha \sigma }^\nu ^\beta )+3(\eta ^{\mu \alpha }\eta ^{\rho \sigma }^\nu ^\beta `$ (E7)
$`+\eta ^{\mu \alpha }\eta ^{\nu \rho }\eta ^{\beta \sigma }\mathrm{}\eta ^{\mu \alpha }\eta ^{\nu \rho }^\beta ^\sigma \eta ^{\mu \alpha }\eta ^{\beta \sigma }^\nu ^\rho +\eta ^{\nu \beta }\eta ^{\rho \sigma }^\mu ^\alpha +\eta ^{\nu \beta }\eta ^{\mu \rho }\eta ^{\alpha \sigma }\mathrm{}`$ (E8)
$`\eta ^{\nu \beta }\eta ^{\mu \rho }^\alpha ^\sigma \eta ^{\nu \beta }\eta ^{\alpha \sigma }^\mu ^\rho \eta ^{\nu \alpha }\eta ^{\rho \sigma }^\mu ^\beta \eta ^{\nu \alpha }\eta ^{\mu \rho }\eta ^{\beta \sigma }\mathrm{}+\eta ^{\nu \alpha }\eta ^{\mu \rho }^\beta ^\sigma `$ (E9)
$`+\eta ^{\nu \alpha }\eta ^{\beta \sigma }^\mu ^\rho \eta ^{\mu \beta }\eta ^{\rho \sigma }^\nu ^\alpha \eta ^{\mu \beta }\eta ^{\nu \rho }\eta ^{\alpha \sigma }\mathrm{}+\eta ^{\mu \beta }\eta ^{\nu \rho }^\alpha ^\sigma +\eta ^{\mu \beta }\eta ^{\alpha \sigma }^\nu ^\rho )`$ (E10)
$`+\mathrm{\hspace{0.17em}2}(\eta ^{\mu \alpha }\eta ^{\nu \beta }\eta ^{\nu \alpha }\eta ^{\mu \beta })(^\rho ^\sigma \eta ^{\rho \sigma }\mathrm{})]h_{\rho \sigma }.`$ (E11)
## F The integrals $`J_a(s)`$
For the integrals $`J_a(s)`$, $`a=1,2,3`$, defined in (154), we find the following results
$`J_1(s)={\displaystyle \frac{1}{4(\kappa ^2+\pi ^2)|p|^2}}\{{\displaystyle \frac{1}{2\mathrm{Re}p}}\mathrm{ln}\left[{\displaystyle \frac{s^22\mathrm{Re}ps+|p|^2}{s^2+2\mathrm{Re}ps+|p|^2}}\right]`$ (F1)
$`+{\displaystyle \frac{1}{\mathrm{Im}p}}[\pi \mathrm{arctan}\left({\displaystyle \frac{s+\mathrm{Re}p}{\mathrm{Im}p}}\right)\mathrm{arctan}\left({\displaystyle \frac{s\mathrm{Re}p}{\mathrm{Im}p}}\right)]\},`$ (F2)
$`J_2(s)={\displaystyle \frac{1}{4(\kappa ^2+\pi ^2)}}\{{\displaystyle \frac{1}{2\mathrm{Re}p}}\mathrm{ln}\left[{\displaystyle \frac{s^2+2\mathrm{Re}ps+|p|^2}{s^22\mathrm{Re}ps+|p|^2}}\right]`$ (F3)
$`+{\displaystyle \frac{1}{\mathrm{Im}p}}[\pi \mathrm{arctan}\left({\displaystyle \frac{s+\mathrm{Re}p}{\mathrm{Im}p}}\right)\mathrm{arctan}\left({\displaystyle \frac{s\mathrm{Re}p}{\mathrm{Im}p}}\right)]\},`$ (F4)
$`J_3(s)={\displaystyle \frac{1}{4(\kappa ^2+\pi ^2)}}\{4s+{\displaystyle \frac{1}{2\mathrm{Re}p}}[3(\mathrm{Re}p)^2(\mathrm{Im}p)^2]\mathrm{ln}\left[{\displaystyle \frac{s^2+2\mathrm{Re}ps+|p|^2}{s^22\mathrm{Re}ps+|p|^2}}\right]`$ (F5)
$`+{\displaystyle \frac{1}{\mathrm{Im}p}}[(\mathrm{Re}p)^23(\mathrm{Im}p)^2][\pi \mathrm{arctan}\left({\displaystyle \frac{s+\mathrm{Re}p}{\mathrm{Im}p}}\right)\mathrm{arctan}\left({\displaystyle \frac{s\mathrm{Re}p}{\mathrm{Im}p}}\right)]\},`$ (F6)
where $`p`$ is a function of $`s`$ given by expressions (155), which give $`|p|^2=\left[\left[(\kappa ^2+\pi ^2)s^2+\kappa \right]^2+\pi ^2\right]^{1/2}/(\kappa ^2+\pi ^2)`$. |
warning/0001/hep-ph0001082.html | ar5iv | text | # Propagation of ultrahigh-energy neutrinos through the Earth
## Abstract
The dispersion relation in matter of ultrahigh-energy neutrinos above the pole of the $`W`$ resonance ($`E_\nu \text{}\text{10}^7\text{GeV}`$), is studied. We perform our calculation using the real-time formulation of Thermal Field Theory in which the massless limit for the $`W`$ boson is taken. The range of active-to-sterile neutrino oscillation parameters for which there is significant mixing enhancement during propagation through the interior of the Earth, and therefore significant attenuation of neutrino beams in the Earth at high energies, is estimated. Finally, this range is considered in view of the cosmological and astrophysical constraints.
It is now well established that the Earth’s diameter exceeds the attenuation length of neutrinos with energies greater than 25TeV. Such an estimate was based on the calculation of the cross sections for $`\nu N`$ collisions at ultrahigh energies (UHE), ($`E_\nu \text{}1\text{TeV}`$). Because of the smallness of the electron mass, $`\nu e`$ interactions are generally considered as negligible with respect to $`\nu N`$ interactions; $`\nu N`$ interactions therefore provide the dominant signal and account for most of the attenuation of neutrino beams in the interior of the Earth at ultrahigh energies . There is one exception though, the resonant formation of the intermediate $`W^{}`$ boson in $`\overline{\nu }e`$ interactions in the neighborhood of $`E_\nu ^{res}=M_W^2/2m_e\text{6.3}\times \text{10}^{15}\text{eV}`$.
The promising tool for detection of UHE cosmic neutrinos by means of neutrino telescopes consists of recording the long-range muons produced in charged-current $`\nu N`$ interactions that occurs in matter surrounding the detector. Apart from efficient shielding from the flux of atmospheric muons, such upward-going muon events have the advantage of enhancing the effective volume in proportion to the range of the produced muons (typically a few kilometers for $`E_\mu \text{10}\text{TeV}`$). For our purpose it is important to note that the rate for upward-going muons does not depend only on the probability for neutrino conversion to a muon with energy above the threshold energy, but also through the interaction length which is responsible for the attenuation of the neutrino flux due to interactions in the Earth’s interior. The typical situation that occurs for $`E_\nu \text{}\text{10}^5\text{GeV}`$ is that the upward rates depend little on the calculated $`\nu N`$ cross sections, since the enhanced (weakened) interaction rate is nearly compensated by the enhanced (weakened) attenuation of UHE neutrinos propagating through the Earth. On the other hand, the detection of cosmic neutrinos at energies of $`\text{10}^{16}\text{eV}`$ or larger is beset by the problem of the increased importance of attenuation of neutrino beams . However, even so, the upward rates produced by neutrinos from powerful radiation sources, like Active Galactic Nuclei (AGN), should be observable in a detector whose effective area is $`A0.1\text{km}^2`$.
As for event rates involving electron neutrinos, they are generally smaller than the muon event rate by the flux ratio (the initial fluxes of UHE neutrinos originating from AGNs are expected to have a ratio $`\nu _e/\nu _\mu 1/2`$) times the detector length divided by the mean muon range, because of the rapid energy loss of electrons (or annihilation for positrons). Still, it was shown recently that the Landau-Pomeranchuk-Migdal effect may effectively enhance the electron range by detecting upward-going air showers initiated by the $`\nu _e`$ interaction near the Earth’s surface. On the other hand, resonant $`\overline{\nu }e`$ scattering contributes significantly to the attenuation of $`\overline{\nu }_e`$’s, meaning that the flux of electron antineutrinos in the range $`2\times \text{10}^{15}\text{eV}E_\nu 2\times \text{10}^{16}\text{eV}`$ is extinguished for neutrinos traversing the Earth .
In the present paper we are going to consider another mechanism for attenuation of UHE neutrinos propagating through the Earth, namely, matter enhanced neutrino oscillations $`\nu _e\nu _s`$, where $`s`$ is a sterile neutrino (e.g. a singlet under the gauge symmetry of the Standard model). We shall be concerned exclusively with the case of $`\nu _e\nu _s`$ oscillations where the energy of UHE neutrinos is above the resonant energy, $`E_\nu \text{}E_\nu ^{res}`$. Owing to the new form of the effective matter potential in this regime, we are in position to study MSW resonance effect previously ignored in the literature. On the other hand, the matter effect of the Earth through the standard effective potentials ($`E_\nu <<E_\nu ^{res}`$), in the region of oscillation parameters relevant for the solar and atmospheric neutrinos as well as those suggested by the LSND result, is well established now . Even more, a new effect of matter-enhanced neutrino mixing, based on a maximal constructive interference among transition amplitudes, has been discovered recently . In the following, we shall first derive the induced mass squared of the electron neutrino with $`E_\nu \text{}E_\nu ^{res}`$, then we find the range of neutrino parameters for which there is matter-enhanced $`\nu _e\nu _s`$ oscillation during propagation of $`\nu _e`$ through the Earth and finally we discuss if the established range could survive constraints from type II supernovae as well as big bang nucleosynthesis on $`\nu _e\nu _s`$ mixing.
Effects of a medium on neutrino propagation is determined by the difference of potentials, whose standard-model contribution in the context of Thermal Field Theory (TFT) may readily be obtained from the relevant thermal self-energies of a neutrino: charged-current, neutral-current and tadpole. Of these, only the charged-current and tadpole contribution are relevant for our consideration. Let us consider the charged-current diagram in more detail. Since the “target” electrons are massive, the $`W`$ boson may be considered “massless” always when $`s2E_\nu m_eM_W^2`$. This corresponds to energies in the lab frame $`E_\nu \text{}6\times \text{10}^{15}\text{eV}`$. In the opposite limit, $`sM_W^2`$, the $`W`$ boson should be considered “massive” and the usual contact approximation for the $`W`$-propagator is adequate. Using the real-time formulation of TFT, we discover by explicit calculation that the induced mass squared of $`\nu _e`$ is equivalent to the fermion thermal mass squared ,
$$A_\nu ^{cc}=\frac{g^2}{2\pi ^2}_0^{\mathrm{}}k𝑑kn_e(k_0)(E_\nu E_\nu ^{res}),$$
(1)
where $`g0.63`$ is the gauge coupling constant. Note that in contrast to the mass squared induced by the standard MSW potential,
$$A_\nu ^{cc}=2\sqrt{2}G_FN_eE_\nu (E_\nu E_\nu ^{res}),$$
(2)
(1) is independent of neutrino energy and also there is no explicit dependence on the number density of electrons. It should be clearly stated here that actually we are not dealing with field theory in equilibrium since all the electrons in the medium are bound electrons. Still, one is allowed to retain the usual real-time formalism by taking the (11)-component of the electron propagator to be
$$S_{11}(k)=(\overline{)}k+m_e)\left(\frac{1}{k^2m_e^2+iϵ}+2\pi in_e(k_0)\delta (k^2m_e^2)\right),$$
(3)
where now a bound electron is assigned a distribution $`n_e(k_0)`$. Thus (1) describes the plane-wave impulse approximation, which is, for instance, the basic approximation of electron momentum spectroscopy of atoms and molecules .
Although a connection with equilibrium TFT is now established, one may still wonder why (1) is equivalent to the fermion thermal mass, that is, to the characteristic mass scale which appears naturally only in the high-temperature limit of the fermion self-energy. Before going into details, let us stress that to first order in perturbation theory at high temperature, the poles of the full fermion propagator are determined just by the thermal mass . In hot theories, the high-temperature limit means that the temperature is much larger than mass of the particles under consideration ($`m_e,M_W`$ in our case) and the external momenta. Consequently, all particle masses can be ignored for practical purposes. In the standard electro-weak theory we now show that the opposite case, when the external momenta are large and the characteristic scale of the medium is much less than the particle masses, can faithfully mimic the high-temperature limit. Indeed, when $`E_\nu E_\nu ^{res}`$, the four-momentum squared of the $`W`$ boson is always much larger than $`M_W^2`$, and hence the $`W`$ boson can be considered “massless”. To get rid of the electron mass, note that it appears in the numerator of the thermal part of the electron propagator, and both in the numerator and denominator of the vacuum part of $`S_{11}`$ \[see Eq. (3)\]. It is however trivially to see that $`m_e`$ from the numerator disappears when sandwiched between the two electroweak vertices $`\gamma _\mu L`$, where $`L=\frac{1}{2}(1\gamma _5)`$ is the projection operator for the left-chiral fermions. On the other hand, the real part of the neutrino self-energy is given only by a contribution where a cut is through the electron line, since $`W`$ bosons are absent from the medium. This means that there is no contribution from the vacuum part of $`S_{11}`$, where $`m_e`$ appears in the denominator . This completes the proof that the high-energy limit is equivalent to the high-temperature limit of hot gauge theories. Strictly speaking, since the thermal mass always involves the thermal effects from both the $`W`$ boson and the electron propagators, the only difference is that the former is absent in (1).
Using the above interpretation for bound states and keeping the same normalization as for quasifree states, one can rewrite (1) in the following form,
$$A_\nu ^{cc}0.2<k^1>N_e,$$
(4)
where $`<k>`$ is the average momentum of bound electrons. For the rough estimates presented here, it is sufficient to assume that $`<k^1><k>^1`$. Let us choose the average momentum per atom (with the atomic number $`Z`$) as a quantity of interest here. Going back to atomic physics, one can determine this quantity by applying the Thomas-Fermi method to the calculation of the total ionization energy of a neutral atom. The result is
$$<k>^Z4.6Z^{2/3}\text{keV}.$$
(5)
For our purpose, let us recall that the interior of the Earth consists of two regions of slowly varying density - the core and the mantle, with particularly strong density change between the lower mantle and outer core. The density profile of the Earth can be found in . The density of the mantle increases from $`3`$ to $`5.5\text{gcm}^3`$ (the average value is $`4.7\text{gcm}^3`$ and the average electron fraction is $`0.49`$), while the density of the core varies from $`10`$ to $`13\text{gcm}^3`$ (the average value is $`11.8\text{gcm}^3`$ and the average electron fraction is $`0.47`$). The core comprises heavier elements, presumably nickel ($`Z=28`$). Hence from (5) we have
$$<k>_{core}^Z42\text{keV}.$$
(6)
The mantle consists of lighter elements ($`Z=816`$), and our estimate in this case is
$$<k>_{mantle}^Z25\text{keV}.$$
(7)
Before determination from the resonance condition of a range of neutrino masses where maximum mixing enhancement may occur, one should consider the tadpole graph as well. Being a constant independent of the external neutrino momentum, it is the same as in the standard MSW case ($`E_\nu E_\nu ^{res}`$), giving rise to the induced mass squared which grows linearly with energy, i.e.,
$$A_\nu ^{tadpole}=\sqrt{2}G_FN_nE_\nu ,$$
(8)
where $`N_n`$ is the neutron number density. Apart from a negative sign in (8), let us compare the magnitude of (8) with the charged-current contribution (1). It turns out that for $`E_\nu \text{}\text{10}^9\text{GeV}`$ (8) is beginning to dominate over (1). For our purpose, it is however enough to consider $`E_\nu \text{10}^7\text{10}^8\text{GeV}`$ since the short $`\nu _e`$ interaction length for energies $`E_\nu >\text{10}^9\text{GeV}`$ means that the flux of electron neutrinos is extinguished for neutrinos traversing the Earth. Hence for $`E_\nu \text{}\text{10}^9\text{GeV}`$, the total induced mass squared for $`\nu _e`$ is essentially given by (1). It is interesting to note that at a particular energy around $`\text{10}^9\text{GeV}`$ there is a nearly complete cancellation of matter effects in the neutrino propagator due to the sign of (8).
With the above simplifications and the range of densities in the Earth’s interior as discussed before, one finds from the resonance condition, $`A_\nu ^{cc}\mathrm{\Delta }m_{es}^2`$, a range of neutrino masses where maximal mixing enhancement may occur,
$$0.07<\mathrm{\Delta }m_{es}^2/keV^2<0.12(10^9GeV\text{}E_\nu \text{}10^7GeV).$$
(9)
Here we have taken a small mixing angle, $`\mathrm{cos}2\theta 1`$, in order to study oscillation enhancement. $`\mathrm{\Delta }m_{es}^2>0`$ in (9) means that $`\nu _s`$ is heavier than $`\nu _e`$. This is also true for $`\overline{\nu }_e\overline{\nu }_s`$ oscillations as $`A_{\overline{\nu }}^{cc}>0`$ (in contrast to the standard MSW potential where the sign is reversed for antineutrinos).
Notice that because of the resonance condition which is energy independent and the range of densities in the Earth, the range (9) is actually very small. Even so, it lies in the region which might be very interesting to astrophysics as well as cosmology. We recall that the prediction for the sterile neutrino of $``$ keV mass is not in contradiction with any of the present bound. Indeed, the $``$ keV mass is needed if active-to-sterile neutrino oscillations are to solve the pulsar velocity puzzle . In contrast to active-to-active oscillations, this solution is not in conflict with the cosmological bound on stable neutrino masses since the $``$ keV mass sterile neutrino has been proposed as a viable dark-matter candidate .
It is well known that if we add the reported results from the LSND collaboration to the list of neutrino anomalies, then the explanation of all of them requires a four-neutrino scheme with three active neutrinos $`\nu _e`$, $`\nu _\mu `$, $`\nu _\tau `$ and one electroweak-singlet neutrino . This introduces the ’LSND gap’ of a few eV, a few orders of magnitude below our preferred range (9). However, the LSND result is not confirmed (although not completely ruled out) by a similar KARMEN experiment .
It is easy to estimate the range of active-to-sterile neutrino parameters for which there is significant enhancement mixing during propagation through the Earth. For significant transitions to developed, it is necessary that the propagation distance be greater than about a quarter of a wavelength at resonance . This constraint gives us a lower limit on the mixing angle. Taking the longest distance through the Earth ($`2R_{earth}`$) we have,
$$\frac{\pi E_\nu }{\mathrm{\Delta }m_{es}^2\mathrm{sin}2\theta }<2R_{earth},$$
(10)
which for $`E_\nu =\text{10}^7\text{GeV}`$ and $`2R_{earth}=\text{1.27}\times \text{10}^9\text{cm}`$ gives $`(\mathrm{sin}2\theta )_{core}>\text{2}\times \text{10}^3`$ and $`(\mathrm{sin}2\theta )_{mantle}>\text{3.4}\times \text{10}^3`$.
There is however a stronger limit on $`\mathrm{sin}2\theta `$ coming from the condition for unsuppressed oscillations. The oscillation frequency must be real, otherwise the system is critically overdamped, the oscillations would be fully incoherent, and hence in fact there will be no oscillations. Since for neutrino energy in the range $`\text{10}^{15}\text{eV}E_\nu \text{10}^{21}\text{eV}`$, the cross section scales with $`E_\nu `$ as $`\sigma E_\nu ^{0.4}`$ , one finds the mean free path for neutrinos with $`E_\nu =\text{10}^7\text{GeV}`$ to be $`l=\text{0.1}\times 2R_{earth}`$. The condition for unsuppressed oscillations,
$$l_m<2l,$$
(11)
then gives $`(\mathrm{sin}2\theta )_{core}>\text{4}\times \text{10}^2`$ and $`(\mathrm{sin}2\theta )_{mantle}>\text{7}\times \text{10}^2`$. Notice that although the Earth is opaque to UHE neutrinos, the oscillations may proceed unsuppressed whenever the above requirement is satisfied.
Let us finally check up the range ($`\mathrm{\Delta }m_{es}^2/keV^2\text{0.12},\mathrm{sin}2\theta >\text{4}\times \text{10}^2`$) from a viewpoint of astrophysics and cosmology. The effect of a resonant $`\nu _e\nu _s`$ mixing on a type II supernovae was considered in . The bounds on ($`\mathrm{\Delta }m_{es}^2,\mathrm{sin}2\theta `$) derived in are valid only if the sterile neutrinos have a mean free path larger than the radius of the supernova core after passing the resonance; this is the case if $`\mathrm{sin}2\theta <\text{3}\times \text{10}^2`$. Our range is therefore unaffected by the type II supernovae constraint. On the other hand, a naive bound on the $`\nu _e\nu _s`$ mixing from big bang nucleosynthesis was derived, $`\mathrm{\Delta }m_{es}^2\mathrm{sin}^42\theta \text{}\text{5}\times \text{10}^6\text{eV}^2`$ . Notice a disagreement of our preferred range with the above naive bound. However, the naive calculations ignored the creation of $`\nu \overline{\nu }`$ asymmetries by active-sterile oscillations in the early universe; these may efficiently suppress $`\nu _e\nu _s`$ oscillations, and therefore invalidate the conclusions drawn from naive calculations (even maximal $`\nu _\mu \nu _s`$ oscillations as a solution of the atmospheric neutrino anomaly cannot be excluded ). The measurement of the flux of UHE neutrinos could thus provide us with a new test for this cosmological scenario.
To summarize, the dispersion relation for electron (anti)neutrinos in the Earth’s interior for energies above the pole of the $`W`$ resonance, is derived. Then we have considered MSW oscillations for cosmic neutrinos traversing the Earth by including the charged-current self energy diagram for $`\nu _e`$. We have shown that the range of neutrino masses where maximal enhancement may occur could be interesting from a viewpoint of astrophysics and cosmology. Let us finally stress that in order to study a nadir angle dependence beyond $`\text{34}^o`$, where neutrinos always propagate outside the core, a weaker attenuation of a $`\nu _e`$ beam would require the inclusion of the tadpole self-energy for energies beyond $`\text{10}^9\text{GeV}`$. This interesting case in now under study.
Acknowledgments. The author acknowledges the support of the Croatian Ministry of Science and Technology under the contract 1 – 03 – 068. |
warning/0001/cond-mat0001048.html | ar5iv | text | # Phase ordering at the lambda transition in liquid 4He
## 1 Introduction
As suggested by London and substantiated by Feynman there is an intimate relation between the lambda transition in <sup>4</sup>He and the Bose-Einstein condensation of the ideal Bose gas (IBG). The IBG explains some basic properties of liquid <sup>4</sup>He like for example the irrotational superfluid flow. We propose a modification of the IBG that leads to a realistic expression for the specific heat.
The single particle functions of the IBG are of the form
$$\phi _𝐤\mathrm{exp}(\mathrm{i}𝐤𝐫)=\mathrm{exp}[\mathrm{i}(k_1x+k_2y+k_3z)].$$
(1)
Without changing the kinetic energy we may alternatively consider the single particle functions
$$\phi _𝐤^{\mathrm{p}.\mathrm{o}.}\mathrm{exp}\left[\mathrm{i}(k_2y+k_3z)\right]\mathrm{sin}(qx+\varphi ),$$
(2)
where $`q=k_1`$. For an arbitrary direction (chosen here as the $`x`$ direction) the functions $`\phi _𝐤^{\mathrm{p}.\mathrm{o}.}`$ have a specific phase $`\varphi `$; they are phase ordered (p.o.). This phase ordering implies the following correlation: For $`k_1=k_1^{}=q`$, the probability densities $`|\phi _𝐤^{\mathrm{p}.\mathrm{o}.}|^2`$ and $`|\phi _𝐤^{}^{\mathrm{p}.\mathrm{o}.}|^2`$ have common minima and maxima. That means that two particles with the same momentum in one extinguished direction have an additional positive spatial correlation. No such correlations are present for the single particle functions (1).
The phase ordering leads to a correlation energy of the form
$$E_{\mathrm{corr}}\underset{q}{}\nu _q^{\mathrm{\hspace{0.17em}2}}\text{with}\nu _q=\underset{k_2,k_3}{}n_𝐤.$$
(3)
Here $`n_𝐤`$ is the number of atoms with momentum $`𝐤=(q,k_2,k_3)`$ and $`\nu _q`$ is the number of atoms with common momentum $`q`$ in $`x`$ direction. In the group of $`\nu _q`$ atoms each one has an extra correlation with the other atoms; this results in a $`\nu _q^{\mathrm{\hspace{0.17em}2}}`$ contribution to the energy. This correlation energy has been investigated in detail in Ref. .
We will use the IGB many-body wave functions together with Jastrow factors as proposed by Chester . The effective interaction on the model level is then attractive (the hard core is cut out by the Jastrow factors). This means that the energy $`E_{\mathrm{corr}}`$ is negative and that the correlations are favoured by the condition of minimal free energy.
For a finite correlation energy one must assume that the phase ordering is local. The phase correlations of the single particle functions (2) could be enforced by physical boundary conditions for the considered macroscopic volume $`V`$ (instead of periodic boundary conditions for the single particle functions (1)). In this case the correlations are a surface effect implying $`E_{\mathrm{corr}}V^{2/3}`$, and $`E_{\mathrm{corr}}/V0`$ for $`V\mathrm{}`$. Assuming local phase ordering introduces a parameter $`V_0`$ that defines the finite range over which the single particle functions are correlated. The directions ($`x`$ direction in Eq. (2)) of phase ordering are in general different in various subvolumes. The statistical average over these subvolumes ensures the isotropy of the bulk liquid. Evaluating the energy (3) with the IBG expectation values $`n_𝐤`$ yields a logarithmic singularity , i.e., $`E_{\mathrm{corr}}t\mathrm{ln}|t|`$, where $`t`$ is the relative temperature. Adjusting the amplitude of this singularity to the experimental value fixes the model parameter $`V_0`$ (leading to $`V_0/v10^5`$, where $`v`$ is the volume per atom).
It is the aim of this paper to calculate the entropy change due to the assumed local phase ordering. This calculation is independent of the previous determination of the correlation energy (3). In order to calculate the phase ordering entropy we establish a relation between the functions (2) and coherent states; this relation is analogous to Anderson’s connection between the condensate wave function and locally coherent states. For coherent states the phase variances are given by $`\mathrm{\Delta }\varphi _𝐤=1/(2\sqrt{n_𝐤})`$. Restricting the phases by $`\mathrm{\Delta }\varphi _𝐤`$ yields the phase ordering entropy $`S_{\mathrm{p}.\mathrm{o}}(n_𝐤)`$. Using the occupation numbers $`n_𝐤`$ of the IBG form we find logarithmic singularity for the thermodynamic entropy, $`S_{\mathrm{p}.\mathrm{o}}t\mathrm{ln}|t|`$. For the specific heat per particle, $`c=A\mathrm{ln}|t|`$, we obtain the result $`A=3k_\mathrm{B}/2`$, where $`k_\mathrm{B}`$ is Boltzmann’s constant.
As a theory of the lambda transition, the presented approach is a modest attempt only. Instead of determining exact many-body states we make a guess of the relevant correlations. Instead of evaluating the partition sum we use the IBG expectation values (arguing that the correlation effects are of minor influence on the occupation of the individual levels). The determination of the exact many-body states and the exact evaluation of the partition sum is, of course, not feasible. One might, however, try to find a suitable variational ansatz for the many-body states for which the free energy can be calculated and minimized (like it may be done for BCS states). The present evaluation of the entropy together with the previous evaluation of the correlation energy (3) are a first step towards such an approach. As it stands, we present a phenomenological many-body model that yields a promising result, namely a logarithmic singularity with a sensible strength.
As a many-body approach for the lambda transition the presented model lies outside the main stream of realistic models for liquid helium. There are several well-known and thoroughly studied microscopic many-body approaches (for an overview see Ref. ), like diagrammatic (Green’s function) approaches , variational methods (in particular the correlated basis function method and its generalizations ), and Monte Carlo calculations. As far as the approaches are based on elementary excitations they are restricted to low temperatures. For temperatures around the lambda point, there are successful applications of path-integral Monte Carlo simulations and of the variational approach . These approaches do, however, not yield analytic and realistic expressions for the asymptotic behavior. What is nowadays widely considered as the theory of the lambda transition and its asymptotic properties is the renormalization group theory . In contrast to the model presented here, this theory is not a many-body description.
Several authors have considered coherent states in the description of liquid helium, too. Their investigations concern mostly the connection between the condensate and coherent states as in Anderson’s paper . The local phase ordering that we consider is not specifically related to any of these other investigations.
## 2 Phase ordering entropy
### 2.1 Many-body states
Chester studied the density matrix for many-body wave functions of the form $`\mathrm{\Psi }=𝒮F\mathrm{\Psi }_{\mathrm{IBG}}`$. Here $`S`$ denotes the symmetrization operator, $`F=_{i,j}f(r_{ij})`$ is the Jastrow factor, $`\mathrm{\Psi }_{\mathrm{IBG}}(𝐫_j,n_𝐤)`$ describes the IBG state, and $`𝐫_j`$ defines the positions of the helium atoms. Using hard sphere factors for $`f(r)`$, Chester obtained the IBG values $`n_𝐤_{\mathrm{IBG}}`$ for the average occupation numbers $`n_𝐤`$. McMillan used the ansatz $`\mathrm{\Psi }=𝒮F`$ for the ground state; for realistic atom-atom interactions he determined the function $`f(r)`$ by minimizing the energy. This yields a realistic pair correlation function and a realistic binding energy.
We follow basically Chester’s approach replacing, however, the single particle functions (1) by the phase ordered ones (2). The many-body wave function is then of the form
$$\mathrm{\Psi }=𝒮F\underset{j=1}{\overset{N}{}}\phi _{𝐤_j}^{\mathrm{p}.\mathrm{o}.}(𝐫_j)=𝒮F\underset{𝐤}{}\left[\phi _𝐤^{\mathrm{p}.\mathrm{o}.}\right]^{n_𝐤}.$$
(4)
The factor $`[\phi _𝐤]^{n_k}`$ in the last expression stands for the product $`\phi _𝐤(𝐫_{\nu +1})\phi _𝐤(𝐫_{\nu +2})\mathrm{}\phi _𝐤(𝐫_{\nu +n_k})`$. The total number of atoms is denoted by $`N`$.
As explained in the Introduction, a finite correlation energy (3) requires finite extensions of the phase ordered single particle functions. For this purpose the volume $`V`$ is thought to be divided into $`V/V_0`$ subvolumes of size $`V_0`$, and the single particle functions are restricted to such subvolumes. The phase ordering can then be obtained by requiring physical boundary conditions for the subvolumes, or by considering coherent states in these subvolumes (Sec. 2.3). The product over the subvolumes might be explicitly displayed in the last expression of Eq. (4). We will adjust the prefactor of the momentum integrals such that the result corresponds to the total volume $`V`$.
The Jastrow factors are of no influence on the statistical counting. Without phase ordering (i.e., using the single particle functions (1) in Eq. (4)) one obtains, therefore, the well-known IBG expression for the entropy:
$$S_{\mathrm{IGB}}=k_\mathrm{B}\underset{𝐤}{}\left[(1+n_𝐤)\mathrm{ln}(1+n_𝐤)n_𝐤\mathrm{ln}(n_𝐤)\right].$$
(5)
The entropy change due to the phase ordering may be divided into two parts:
1. All $`\nu _q`$ atoms in states with $`k_1=q`$ must adopt the same (mean) phase. This leads to a replacement of the entropy (5) by a modified expression $`S_{\mathrm{IBG}}^{}`$.
2. A well-defined phase requires a small phase variance. The entropy change due to this restriction is denoted by $`S_{\mathrm{p}.\mathrm{o}.}`$. This decisive contribution is called phase ordering entropy.
### 2.2 Mean phase
Admitting the values $`\varphi =0`$ and $`\varphi =\pi /2`$ in Eq. (2) means that we use the functions $`\mathrm{sin}(qx)`$ and $`\mathrm{cos}(qx)`$. This is equivalent to the use of $`\mathrm{exp}(\mathrm{i}k_1x)`$ and $`\mathrm{exp}(\mathrm{i}k_1x)`$, i.e., to the case (1) of no phase ordering. Using the single particle functions (2) with one definite phase (lets say $`\varphi =0`$) effectively implies that only every second single particle state is occupied. This can be accounted for by the substitutions $`_𝐤(1/2)_𝐤`$ and $`n_𝐤2n_𝐤`$ in the expression (5), i.e.,
$$S_{\mathrm{IBG}}^{}=\frac{k_\mathrm{B}}{2}\underset{𝐤}{}\left[(1+2n_𝐤)\mathrm{ln}(1+2n_𝐤)2n_𝐤\mathrm{ln}(2n_𝐤)\right].$$
(6)
The step from the entropy (5) to the entropy (6) takes into account the choice $`\varphi =0`$ for all atoms. Deviating from this, we could admit different $`\varphi _q`$ values in Eq. (2) for different $`q`$’s without destroying the considered correlations. This is, however, a negligible effect because the increase in entropy would be $`\mathrm{\Delta }S/N=_q\mathrm{}/N𝒪(N^{2/3})0`$.
The expectation value of $`S_{\mathrm{IBG}}^{}`$ has similar temperature dependence as that of $`S_{\mathrm{IBG}}`$ (Sec. 3.2).
### 2.3 Phase variances
Anderson has established a connection between the condensate wave function and locally coherent states. The coherent states are constructed in suitably chosen subvolumes $`\mathrm{\Delta }V`$. The volumes must be large enough in order to define a mean phase $`\overline{\varphi }`$ with a small variance $`\mathrm{\Delta }\varphi 1`$; this requires that the number of atoms in $`\mathrm{\Delta }V`$ is large compared to 1. On the other hand, the volumes $`\mathrm{\Delta }V`$ must be small enough in order to ensure $`\overline{\varphi }\text{const.}`$ in spite of macroscopic perturbations or flows.
The connection between the condensate wave function and locally coherent states may be carried over to noncondensed phase ordered states (2) provided that $`n_𝐤1`$. This means that we may construct coherent states that correspond to the single particle functions (2). As in Anderson’s work these coherent states are restricted to finite volumes. For the following calculation it is not necessary to quantify the size of these subvolumes.
Appendix A presents an explicit construction of coherent states that correspond to phase ordered single particle functions (2). Here we restrict ourselves to the presented argument based on the analogy to known work .
If the atoms with momentum $`𝐤`$ form a coherent state then their phase variance reads
$$\mathrm{\Delta }\varphi _𝐤=\frac{1}{2\sqrt{\overline{n}_𝐤}}.$$
(7)
The bar denotes the quantum mechanical average.
In Eq. (1) we may admit arbitrary phases $`\varphi _j`$ for each atom, $`\phi _𝐤(𝐫_j)\mathrm{exp}(\mathrm{i}𝐤𝐫_j+\varphi _j)`$. These phases $`\varphi _j`$ are of no influence on physical quantities (like the single particle kinetic energy or an interaction matrix element). In particular, they are of no influence on the statistical counting. Such phases are usually ignored.
Contrary to this, arbitrary phases $`\varphi _j`$ in Eq. (2), $`\phi _𝐤^{\mathrm{p}.\mathrm{o}.}\mathrm{sin}(qx_j+\varphi _j)`$, would destroy the considered correlations. The phase ordering requires, therefore, that these phases $`\varphi _j`$ are adequately restricted, i.e., to an interval of the order (7) around a common mean value. A random phase corresponds to a phase variance $`\mathrm{\Delta }\varphi _{\mathrm{ra}}=\pi /\sqrt{3}`$. (We use the phase definition by Barnett et al. ; the numerical factors are, however, without influence on the critical part of the entropy.) Restricting the phase $`\varphi _j`$ of the $`j`$th atom to the interval (7) corresponds to a reduction factor $`\mathrm{\Delta }\varphi _𝐤/\mathrm{\Delta }\varphi _{\mathrm{ra}}`$. For all atoms this leads to the entropy contribution
$$S_{\mathrm{p}.\mathrm{o}.}=k_\mathrm{B}\mathrm{ln}\underset{𝐤}{}\left(\frac{\mathrm{\Delta }\varphi _𝐤}{\mathrm{\Delta }\varphi _{\mathrm{ra}}}\right)^{\overline{n}_𝐤}\frac{k_\mathrm{B}}{2}\underset{𝐤}{}\overline{n}_𝐤\mathrm{ln}\left(\overline{n}_𝐤\right)(\overline{n}_𝐤1).$$
(8)
This expression appears to be the most simple and plausible way of accounting for the phase restrictions.
Phase ordering requires that the considered phases are reasonably well defined, i.e., $`\mathrm{\Delta }\varphi _𝐤1`$ or $`\overline{n}_𝐤1`$. The condition $`\overline{n}_𝐤1`$ has been used in the last step in Eq. (8). We note that the critical terms in the expectation value of the entropy are solely due to the contributions from $`\overline{n}_𝐤1`$. The expectation value of $`S_{\mathrm{p}.\mathrm{o}.}`$ will exhibit a logarithmic singularity (Sec. 3.3), i.e., the phase ordering entropy is the decisive contribution in our model.
For $`n_𝐤1`$, expression (5) yields $`S_{\mathrm{IBG}}k_\mathrm{B}\mathrm{ln}n_𝐤`$. From this and Eq. (8) it follows that the low $`𝐤`$ contributions to the sum $`S_{\mathrm{IBG}}+S_{\mathrm{p}.\mathrm{o}.}`$ are negative which might seem disturbing. The entropy must, however, not be attributed separately to each $`𝐤`$ level; it rather results from the distribution of all particles onto the available levels. Consider, for example, an isolated level with the single particle energy $`\epsilon _1`$ that is occupied by $`n_1`$ atoms. Then we may attribute the energy $`\epsilon _1n_1`$ to this level but we cannot say that $`k_\mathrm{B}\mathrm{ln}n_1`$ is the entropy of this (isolated) level. Alternatively one may argue as follows: The low $`𝐤`$ contributions to $`S_{\mathrm{IBG}}`$ (or $`S_{\mathrm{IBG}}^{}`$) are small and noncritical. Consequently, any entropy expression that leads to the experimentally observed critical behavior (i.e., to $`S_{\mathrm{exp}}t\mathrm{ln}|t|`$, where $`t`$ is the relative temperature) must override the low $`𝐤`$ contribution of $`S_{\mathrm{IBG}}`$ (similarly as it is done by $`S_{\mathrm{p}.\mathrm{o}.}`$).
## 3 Thermodynamic entropy
### 3.1 Statistical assumptions
In order to determine the thermodynamic entropy we need the expectation values $`n_𝐤`$ of the occupation numbers. For the wave function (4) with the single particle functions (1) and with hard sphere Jastrow factors, Chester obtained the IBG values, i.e., $`n_𝐤=n_𝐤_{\mathrm{IBG}}`$. This means that the IBG occupation numbers remain valid in spite of hard core interactions. The considered phase ordering correlations are small in the sense that $`|E_{\mathrm{corr}}|/Nk_\mathrm{B}T_\lambda `$. Therefore, it appears to be a sensible approximation to use the IBG occupation numbers for the phase ordered states, too. The phenomenological assumption
$$n_𝐤=n_𝐤_{\mathrm{IBG}}$$
(9)
is the basis of our description of the phase transition. It implies that we consider an almost ideal Bose gas model. This classification refers to the statistical assumption (9) but not to an assumption about weak interactions.
The thermodynamic entropy is obtained by
$$S(T,V,N)=S(n_𝐤)=S(n_𝐤),$$
(10)
where
$$S(n_𝐤)=S_{\mathrm{IBG}}^{}+S_{\mathrm{p}.\mathrm{o}.}.$$
(11)
The contributions $`S_{\mathrm{IBG}}^{}`$ and $`S_{\mathrm{p}.\mathrm{o}.}`$ are defined in Eqs. (6) and (8), respectively. For the relevant states ($`n_𝐤1`$) the quantum mechanical variances $`\mathrm{\Delta }n_𝐤/\overline{n}_𝐤1/\sqrt{\overline{n}_𝐤}1`$ are small compared to the statistical variances. Therefore, we omit a distinction between $`\overline{n}_𝐤`$ and $`n_𝐤`$.
The average occupation numbers of the IBG are
$$n_𝐤=\frac{1}{\mathrm{exp}[(\epsilon _𝐤\mu )/k_\mathrm{B}T]1}=\frac{1}{\mathrm{exp}(x^2+\tau ^2)1}.$$
(12)
Here $`\epsilon _𝐤=\mathrm{}^2k^2/2m`$ is the single particle energy and $`\mu `$ is the chemical potential. We use the dimensionless quantities $`x`$ and $`\tau `$,
$$x^2=\frac{\epsilon _𝐤}{k_\mathrm{B}T}=\frac{\lambda ^2k^2}{4\pi }\text{and}\tau ^2=\frac{\mu }{k_\mathrm{B}T},$$
(13)
where $`\lambda =2\pi \mathrm{}/\sqrt{2\pi mk_\mathrm{B}T}`$ denotes the thermal wave length. The transition temperature $`T_\lambda `$ is given by the condition
$$\lambda (T_\lambda )=\left[v^{}\zeta (3/2)\right]^{1/3},$$
(14)
where $`\zeta (3/2)2.6124`$ denotes Riemann’s zeta function and $`v=V/N`$ is the volume per particle. In the following we use the relative temperature
$$t=t(T,v)=\frac{TT_\lambda }{T_\lambda }.$$
(15)
The chemical potential $`\mu `$ vanishes at the transition point. For $`|t|1`$ the particle number condition $`N=n_𝐤`$ yields
$$\tau (t)=\sqrt{\frac{\mu }{k_\mathrm{B}T}}=\{\begin{array}{ccc}at+bt^2+\mathrm{}& & (t>0)\hfill \\ a^{}|t|+b^{}t^2+\mathrm{}& & (t<0).\hfill \end{array}$$
(16)
The coefficients are known for the IBG, in particular $`a=(3/4)\zeta (3/2)/\sqrt{\pi }`$ and $`a^{}=b^{}=\mathrm{}=0`$. As a generalization of the IBG we admit $`\tau =a^{}|t|+\mathrm{}`$ with $`a^{}0`$ for $`t<0`$. This generalization does not change the character of the transition (neither the point of transition nor the critical exponents). Using the second line in Eq. (16) and the particle number condition we obtain the condensate fraction
$$\frac{n_0(t)}{N}=\left(\frac{3}{2}+\frac{2\sqrt{\pi }a^{}}{\zeta (3/2)}\right)|t|+gt^2+\mathrm{}.$$
(17)
Here $`n_0(t)`$ is the expectation value of the occupation number of the lowest energy level.
There are a number of reasons for generalizing the IBG by admitting $`\tau =a^{}|t|+\mathrm{}`$ with a coefficient $`a^{}0`$ in the expansion (16): Formally, this generalization leads to a more symmetric form of this expansion. Physically, a positive value of $`a^{}`$ implies a temperature dependent energy gap $`\mathrm{\Delta }k_\mathrm{B}Ta^2t^2`$ between the condensate and the noncondensed particles and implies that the lowest level is more rapidly occupied than in the IBG. Such a behavior leads to a better agreement between the calculated and the experimental temperature dependence of various quantities, in particular of the superfluid density . As a last argument we mention that $`a^{}>0`$ removes the unphysical divergence ($`1/k^2`$ for $`k0`$) of the static structure function of the IBG for $`t<0`$. Our main result will not depend on the parameter $`a^{}`$.
The well-known expectation value of the IBG entropy (5) is given by
$$\frac{S_{\mathrm{IBG}}(T,V,N)}{Nk_\mathrm{B}}=\frac{5}{2}\frac{v}{\lambda ^3}g_{5/2}(\tau )+\tau ^2,$$
(18)
where
$$g_\nu (\tau )=\underset{n=1}{\overset{\mathrm{}}{}}\frac{\mathrm{exp}(n\tau ^2)}{n^\nu }=\underset{n=1}{\overset{\mathrm{}}{}}\frac{z^n}{n^\nu }$$
(19)
defines Riemann’s generalized zeta function. Usually $`z=\mathrm{exp}(\beta \mu )=\mathrm{exp}(\tau ^2)`$ is taken as the argument of this function. We prefer the argument $`\tau `$ because of its close relation ($`\tau |t|`$) to the relative temperature.
For the IBG entropy one has to use $`\tau 0`$ in Eq. (18) for $`t<0`$. The modification $`\tau =a^{}|t|`$ with $`a^{}>0`$ leads to a specific heat that is continuous at $`T_\lambda `$ (as in the pure IBG) but that falls off more rapidly for $`t<0`$.
### 3.2 Mean phase contribution
Because of the common mean phases the IBG expression (5) is replaced by $`S_{\mathrm{IBG}}^{}`$, Eq. (6). The expectation value of the entropy $`S_{\mathrm{IBG}}^{}`$ is
$$\frac{S_{\mathrm{IBG}}^{}(T,V,N)}{Nk_\mathrm{B}}=\frac{1}{2}\underset{𝐤}{}\left[(1+2n_𝐤)\mathrm{ln}(1+2n_𝐤)2n_𝐤\mathrm{ln}(2n_𝐤)\right].$$
(20)
For a first, crude estimate we take into account only those contributions that come from $`n_𝐤1`$. This yields $`S_{\mathrm{IBG}}k_\mathrm{B}_𝐤\mathrm{ln}n_𝐤`$ and $`S_{\mathrm{IBG}}^{}(k_\mathrm{B}/2)_𝐤\mathrm{ln}n_𝐤`$, i.e., $`S_{\mathrm{IBG}}^{}(T,V,N)S_{\mathrm{IBG}}(T,V,N)/2`$.
The actual results for $`S_{\mathrm{IBG}}`$ and $`S_{\mathrm{IBG}}^{}`$ near the transition point are
$`S_{\mathrm{IBG}}(T,V,N)`$ $``$ $`1.28Nk_\mathrm{B}\left(1+{\displaystyle \frac{3}{2}}t+𝒪(t^2)\right),`$ (21)
$`S_{\mathrm{IBG}}^{}(T,V,N)`$ $``$ $`0.96Nk_\mathrm{B}\left(1+{\displaystyle \frac{3}{2}}t+𝒪(t^2)\right).`$ (22)
Eq. (21) follows from the expression (18). The value of $`S_{\mathrm{IBG}}^{}`$ at the lambda point has been determined by a numerical integration of the r.h.s. of Eq. (20). The next term in the expansion is due to the prefactor $`v/\lambda ^3T^{3/2}`$ in Eq. (18). This prefactor originates from the replacement of the momentum sum by an integral over dimensionless variables; it occurs in both entropy expressions in the same way. The higher terms $`𝒪(t^2)`$ in Eqs. (21) and (22) do not agree; moreover, the coefficients of these $`t^2`$ terms are different for $`t>0`$ and $`t<0`$.
### 3.3 Phase variance contribution
The expectation value of the phase ordering entropy $`S_{\mathrm{p}.\mathrm{o}.}`$, Eq. (8), reads
$$\frac{S_{\mathrm{p}.\mathrm{o}.}(T,V,N)}{Nk_\mathrm{B}}=\frac{1}{2N}\underset{𝐤}{}n_𝐤\mathrm{ln}n_𝐤(|t|<0.1).$$
(23)
For the phase ordering we required $`\mathrm{\Delta }\varphi _𝐤1`$ in Eq. (8), lets say $`\mathrm{\Delta }\varphi _𝐤<0.1`$ or $`n_𝐤>10^2`$. Because of $`n_𝐤1/(x^2+\tau ^2)`$ this implies $`x0.1`$ and $`\tau 0.1`$ or $`|t|<0.1`$. The condition $`|t|<0.1`$ indicates the temperature range in which phase ordering and the corresponding entropy contribution are expected to be relevant.
For evaluating the entropy (23) the momentum sum is replaced by an integral,
$$\underset{𝐤}{}\mathrm{}\frac{V}{(2\pi )^3}_0^{\mathrm{}}𝑑k\mathrm{\hspace{0.17em}4}\pi k^2\mathrm{}.$$
(24)
This step has to be accompanied by a separate consideration of a potential condensate fraction; this will be done below in Eq. (28).
For a subvolume of size $`V_0`$ the r.h.s. of Eq. (24) would be $`V_0/(2\pi )^3_0^{\mathrm{}}𝑑k\mathrm{\hspace{0.17em}4}\pi k^2`$. The subsequent summation over all subvolumes yields the factor $`V/V_0`$; this summation has been taken into account in the prefactor on the r.h.s. of Eq. (24). In Appendix B we discuss the validity of the step (24) with respect to finite spacing $`\mathrm{\Delta }k=\pi /V_0^{1/3}`$ of the momentum values.
We insert the replacement (24), the dimensionless quantities (13) and the occupation numbers (12) into Eq. (23):
$$\frac{S_{\mathrm{p}.\mathrm{o}.}}{Nk_\mathrm{B}}=\frac{2}{\sqrt{\pi }}\frac{v}{\lambda ^3}_0^{\mathrm{}}𝑑xx^2\frac{\mathrm{ln}[\mathrm{exp}(x^2+\tau ^2)1]}{\mathrm{exp}(x^2+\tau ^2)1}=\frac{2}{\sqrt{\pi }}\frac{v}{\lambda ^3}J(\tau ).$$
(25)
The critical $`\tau `$ dependence of the integral $`J(\tau )`$ can be determined analytically (App. C),
$$J(\tau )=J(0)\pi \tau \mathrm{ln}\tau +\pi \tau [1\mathrm{ln}(2)]+𝒪(\tau ^2).$$
(26)
Using this, $`\tau =\left(3\zeta (3/2)/4\sqrt{\pi }\right)|t|+𝒪(t^2)`$, and $`v/\lambda ^3=1/\zeta (3/2)+𝒪(t)`$ we obtain
$$S_{\mathrm{p}.\mathrm{o}.}(T,V,N)=\frac{3Nk_\mathrm{B}}{2}t\mathrm{ln}|t|\pm \mathrm{}(|t|<0.1)$$
(27)
for the leading singularity of the phase ordering entropy for $`t>0`$. As we will see, this result holds for $`t<0`$, too.
For $`t<0`$ the prescription (24) has to be accompanied by the replacement
$$n_𝐤n_𝐤+n_0(t)\delta (𝐤).$$
(28)
A phase variance $`\mathrm{\Delta }\varphi _01/\sqrt{n_0}0`$ stands for a macroscopic phase coherence. This is appropriate for a superfluid flow but not for the considered local phase ordering. For the phase variances (7) we use, therefore, the continuous part of the occupation numbers $`n_𝐤`$ only. The contribution $`S_{\mathrm{p}.\mathrm{o}.}^{\mathrm{cond}}`$ in Eq. (23) due to the condensed particles reads then
$$\frac{S_{\mathrm{p}.\mathrm{o}.}^{\mathrm{cond}}}{Nk_\mathrm{B}}=\frac{n_0(t)}{2N}\mathrm{ln}n_{𝐤0}=\frac{n_0(t)}{2N}\mathrm{ln}[\mathrm{exp}(\tau ^2)1].$$
(29)
This can be evaluated by using Eq. (17) for $`n_0(t)/N`$ and Eq. (16) for $`\tau (t)`$.
For $`t<0`$ the contributions (25) and (29) have to be added. The $`a^{}`$ term of $`n_0(t)/N=[3/2+2\sqrt{\pi }a^{}/\zeta (3/3)]|t|+\mathrm{}`$ in Eq. (29) cancels exactly the logarithmic term in Eq. (25). The surviving leading term comes from $`n_0(t)/N=(3/2)|t|+\mathrm{}`$ and $`\mathrm{ln}[\mathrm{exp}(\tau ^2)1]=2\mathrm{ln}|t|+\mathrm{}`$ yielding again Eq. (27).
### 3.4 Specific heat
The specific heat per particle may be written as
$$c_V=\frac{1}{N}\left(\frac{S}{t}\right)_{V,N}=\{\begin{array}{ccc}A\mathrm{ln}|t|+B+\mathrm{}\hfill & & (t>0)\hfill \\ A^{}\mathrm{ln}|t|+B^{}+\mathrm{}\hfill & & (t<0).\hfill \end{array}$$
(30)
From Eq. (27) we obtain the theoretical amplitudes
$$A_{\mathrm{theor}}=A_{\mathrm{theor}}^{}=\frac{3}{2}k_\mathrm{B}.$$
(31)
If the experimental data for $`c_V`$ are fitted by the form (30) one finds $`A_{\mathrm{exp}}0.63k_\mathrm{B}`$ and $`A_{\mathrm{exp}}^{}0.59k_\mathrm{B}`$; these values are taken from Eq. (4.2) (for low pressure) of Ref. . It is remarkable that we obtain a parameter free result and a sensible size for these amplitudes. The experimental ratio $`A_{\mathrm{exp}}/A_{\mathrm{exp}}^{}`$ deviates from 1 by a few percent. In our model, the contribution (25) of the noncondensed particles to $`S_{\mathrm{p}.\mathrm{o}.}`$ has different signs above and below the lambda point. A correction in this contribution could, therefore, lead to a deviation from $`A_{\mathrm{theor}}/A_{\mathrm{theor}}^{}=1`$.
Experimentally the specific heat $`c_P`$ at constant pressure is more readily accessible. It is this quantity that has been measured over several decades of the relative temperature, eventually down to $`|t|=2\times 10^9`$ in a microgravity experiment . In the context of these measurements and their analyses we note the following points:
1. Theoretically, a logarithmic behavior of one of the two quantities, $`c_P`$ or $`c_V`$, implies an analogous behavior of the other quantity in the experimentally accessible range. The statement that only one of these quantities may diverge for $`|t|0`$ applies to very low (experimentally inaccessible) $`|t|`$ values only .
The experimental values for $`c_P`$ may be fitted by the form (30), too. The corresponding coefficients (as given by Eq. (47) of Ref. ) differ from that for $`c_V`$ by about 5 to 10%. For comparing our result (30) with the experiment we may, therefore, also consider the experimental $`c_P`$.
2. If the critical exponents $`\alpha `$ and $`\alpha ^{}`$ are used as fit parameters one finds small negative values. The following points indicate that this does not rule out a logarithmic singularity (corresponding to $`\alpha 0`$ and $`\alpha ^{}0`$) for the leading term:
1. Arp fitted the data of many experiments in order to obtain an equation of state for helium. He considered specifically the question of the critical exponent $`\alpha `$ (Sec. 13.1.1 of Ref. ) and found no statistical significance for a deviation from a leading logarithmic form. Consequently, he adopted the value $`\alpha =\alpha ^{}=0`$ (standing for a logarithmic function) in his expression for the specific heat $`c_V`$.
2. By the scaling law $`\alpha ^{}=23\nu `$ the critical exponent for the specific heat for $`t<0`$ is related to that of the superfluid density $`\varrho _\mathrm{s}|t|^\nu `$. The deviations from $`\alpha ^{}=0`$ and $`\nu =2/3`$ found in standard fits are consistent with this scaling law.
In Ref. we presented an alternative fit formula for the superfluid density in which the leading exponent $`\nu `$ equals exactly 2/3. This fit formula reproduced the data better than standard fit formulas. This means that the value $`\nu =2/3`$ (corresponding to $`\alpha ^{}=0`$) is not ruled out by the experiment, and that the value of the leading exponent (found in a fit) might depend on the higher order terms (used in the fit).
In view of these points, we consider it to be an open question whether the true behavior deviates from a logarithmic one in a way that cannot be accounted for by higher order terms.
We turn now to the coefficients $`B`$ and $`B^{}`$ of the specific heat (30). From $`S_{\mathrm{IBG}}^{}`$, Eq. (22), we get a contribution of $`1.44k_\mathrm{B}`$ for both, $`B_{\mathrm{theor}}`$ and $`B_{\mathrm{theor}}^{}`$. In Eq. (27) we add the terms that are linear in $`t`$; these terms follow from Eqs. (25) and (29). Including all contributions we obtain
$$B_{\mathrm{theor}}0.52k_\mathrm{B},B_{\mathrm{theor}}^{}=3.7k_\mathrm{B}.$$
(32)
For $`B_{\mathrm{theor}}^{}`$ we used the parameter value $`a^{}3`$ that is (rather uniquely) determined by adjusting the model expression for the superfluid density to the experimental temperature dependence.
The value for $`B_{\mathrm{theor}}`$ compares reasonably well with the experimental value $`B_{\mathrm{exp}}0.84k_\mathrm{B}`$ (from Eq. (4.5) of Ref. for low pressure).
There is a large discrepancy between $`B_{\mathrm{theor}}^{}`$ and the experimental value $`B_{\mathrm{exp}}^{}2.0k_\mathrm{B}`$. In the IBG, the coefficient of the $`t^2`$ term in the free energy $`F_{\mathrm{IBG}}`$ is the same above and below the transition. This means that the IBG does not reproduce the jump of the specific heat that normally accompanies the occurrence of a macroscopic order parameter field. For a crude estimate of this missing jump we consider the Landau free energy $`F_{\mathrm{Landau}}/(Nk_\mathrm{B}T_\lambda )=rt|\psi |^2+u|\psi |^4`$. This Landau model yields a jump of the specific heat, $`\mathrm{\Delta }c_V/k_\mathrm{B}=r^2/2u`$. The identification $`|\psi |^2=n_0(t)/N`$ relates the coefficient in $`|\psi |^2=(r/2u)|t|`$ to that in Eq. (17). This connection leads to $`\mathrm{\Delta }c_V=r[3/2+2\sqrt{\pi }a^{}/\zeta (3/2)]`$. Using again $`a^{}3`$ yields $`\mathrm{\Delta }c_V5.6r`$. For the parameter $`r`$ in the expression for $`F_{\mathrm{Landau}}/(Nk_\mathrm{B}T_\lambda )`$ we may expect $`r𝒪(1)`$. The resulting $`\mathrm{\Delta }c_V5.6k_\mathrm{B}`$ could close the gap between $`B_{\mathrm{theor}}^{}`$ of Eq. (32) and $`B_{\mathrm{exp}}^{}`$.
### 3.5 Entropy at the lambda point
As a last point we consider the value of the entropy at the lambda point. Evaluating the phase ordering entropy (25) as it stands yields
$$S_{\mathrm{p}.\mathrm{o}.}(T_\lambda )=\frac{2J(0)}{\sqrt{\pi }\zeta (3/2)}Nk_\mathrm{B}0.51Nk_\mathrm{B}.$$
(33)
This result contains, however, contributions from momenta for which the condition $`n_𝐤1`$ is not satisfied.
We estimate the value $`S_{\mathrm{p}.\mathrm{o}.}^{\mathrm{est}}(T_\lambda )`$ that comes from the contributions with $`n_𝐤1`$ alone. For small and positive $`t`$ values we may assume the form $`S_{\mathrm{p}.\mathrm{o}.}/NS_{\mathrm{p}.\mathrm{o}.}^{\mathrm{est}}(T_\lambda )/NA_{\mathrm{theor}}t\mathrm{ln}(t)+(B_{\mathrm{p}.\mathrm{o}.}+A_{\mathrm{theor}})t`$. Following Eq. (23) we argued that the phase ordering entropy should be relevant in the range $`|t|<0.1`$ only. This implies $`S_{\mathrm{p}.\mathrm{o}.}(t=0.1)0`$ leading to $`S_{\mathrm{p}.\mathrm{o}.}^{\mathrm{est}}(T_\lambda )0.3Nk_\mathrm{B}`$.
We conclude that the result (33) might be too large by roughly a factor of 2. One may, therefore, expect that the contribution $`S_{\mathrm{p}.\mathrm{o}.}(T_\lambda )`$ accounts for the difference between $`S_{\mathrm{IBG}}^{}0.96Nk_\mathrm{B}`$ and the experimental value $`S_{\mathrm{exp}}(T_\lambda )0.76Nk_\mathrm{B}`$.
## 4 Concluding remarks
We summarize the main features of the presented phenomenological many-body model:
1. Following Chester we use the IBG many-body states multiplied by Jastrow factors.
2. We assume that local phase correlations are the relevant correlations near the lambda point. These correlations are specified by the single particle functions (2).
3. The local phase ordering can be described by coherent states. The phase variances of these coherent states are the basis for the statistical counting of the phase restrictions. This leads to the expression (8) for the phase ordering entropy.
4. For evaluating the phase ordering entropy we use the IBG expression for the average occupation numbers.
It is clear that global phase ordering plays a decisive role in liquid helium below the transition point. This makes our assumption that local phase ordering are relevant correlations near the lambda point to some extent plausible. We have proposed a specific kind of the local phase ordering in the framework of an almost ideal Bose gas model. The use of the IBG occupation numbers is a phenomenological assumption.
Our model yields a logarithmic singularity of the specific heat in a straightforward and rather simple way: The expression (8) for the phase ordering entropy follows from the phase variances (7). Evaluating this expression with the IBG occupation number yields the logarithmic singularity (27) (immediately for $`t>0`$, and after a slight generalization of the occupation numbers for $`t<0`$).
We summarize the main results of our almost ideal Bose gas model:
1. The model yields a realistic expression for the specific heat of liquid helium at the $`\lambda `$ transition. For the strength of the logarithmic singularity the remarkable result $`A=A^{}=3k_\mathrm{B}/2`$ is obtained.
2. The model retains the essential features of the IBG. Therefore, it strengthens the suggested relation between the Bose-Einstein condensation and the $`\lambda `$ transition.
3. The model offers an intriguing picture for the relevant correlations near the $`\lambda `$ point. This picture provides also perspectives beyond the specific heat (some of which are indicated in App. B).
Various other quantities may be calculated in the framework of our almost ideal Bose gas model: In an earlier paper we calculated the energy as a function of the temperature; the present derivation of the specific heat appears to be more direct. The critical exponent $`\beta =1/2`$ of the condensate fraction may be reconciled with the actual behavior of the superfluid density by assuming that noncondensed particles move coherently with the condensate . This idea leads to observable consequences that have been discussed in Ref. .
## Appendix A Coherent states
In order to justify the use of relation (7) we establish a connection between the phase ordered single particle functions (2) and coherent states. Our procedure is analogous to Anderson’s construction of coherent states for the condensate.
The many-body wave function (4) can be written in the form
$$\mathrm{\Psi }=𝒮F\underset{j=1}{\overset{n_0}{}}\phi _0(𝐫_j)\underset{j=n_0+1}{\overset{n_0+n_1}{}}\phi _1(𝐫_j)\mathrm{}.$$
(34)
The single particle functions (phase ordered or not) are denoted by $`\phi _i`$, where $`i=0,1,2,\mathrm{}`$ follows the energy sequence. Without the Jastrow factor the off-diagonal single particle density reads
$$\varrho (𝐫,𝐫^{})=\mathrm{\Psi }\left|\widehat{\psi }^+(𝐫^{})\widehat{\psi }(𝐫)\right|\mathrm{\Psi }=n_0\phi _0(𝐫)^{}\phi _0(𝐫^{})+n_1\phi _1(𝐫)^{}\phi _1(𝐫)+\mathrm{},$$
(35)
where $`\widehat{\psi }^+`$ and $`\widehat{\psi }`$ are the particle creation and annihilation operator, respectively. Including the Jastrow factor the condensate contribution in Eq. (35) is depleted from $`n_0`$ to a lower value $`n_\mathrm{c}`$. Similar effects are to be expected for next low-momentum contributions.
All terms in Eq. (35) besides the first term vanish for $`|𝐫𝐫^{}|\mathrm{}`$. This off-diagonal long range order is considered as the decisive criterium for superfluidity.
Anderson related the off-diagonal long range order to the mean field aspect by constructing localized coherent states. The coherent states $`|\mathrm{coh}`$ are defined such that $`\mathrm{coh}|\widehat{\psi }|\mathrm{coh}\mathrm{exp}(\mathrm{i}\varphi )`$, where $`\varphi `$ is the phase of the condensate wave function $`\phi _0\mathrm{exp}(\mathrm{i}\varphi )`$. Accordingly, we will construct coherent states $`|\mathrm{coh}`$ for which $`\mathrm{coh}|\widehat{\psi }|\mathrm{coh}`$ is proportional to the phase ordered single particle function (2). Anderson’s aim was to show that the phase of the condensate and the corresponding particle number may be treated as macroscopic variables. Our aim is to justify Eq. (7) by relating the wave functions (2) to coherent states.
Anderson introduced finite volumes $`\mathrm{\Delta }V`$ in which the phase $`\varphi `$ of the condensate wave function $`\phi _0`$ is approximately constant. We consider phase ordered single particle functions that are restricted to finite subvolumes of size $`V_0`$. In the following we identify $`\mathrm{\Delta }V`$ with $`V_0`$ and treat both cases simultaneously.
The single particle functions $`\phi _k`$ shall form an orthonormal set in one subvolume of size $`V_0`$. Then we may write down the following relations between the field operator $`\widehat{\psi }(𝐫)`$ and the annihilation operator $`\widehat{c}_𝐤`$:
$`\widehat{\psi }(𝐫)`$ $`=`$ $`{\displaystyle \underset{𝐤}{}}\phi _𝐤(𝐫)\widehat{c}_𝐤,`$ (36)
$`\widehat{c}_𝐤`$ $`=`$ $`{\displaystyle _{V_0}}d^3r\phi _𝐤^{}(𝐫)\widehat{\psi }(𝐫).`$ (37)
A state with $`n_𝐤`$ particles of momentum $`𝐤`$ is given by
$$|n=|n_𝐤=\frac{1}{\sqrt{n_𝐤!}}\widehat{c}_𝐤^+|\mathrm{vac},$$
(38)
where $`|\mathrm{vac}`$ is the vacuum state. In the following we restrict ourselves to one specific momentum and omit the index $`𝐤`$. The coherent state is constructed as
$$|\mathrm{coh}=\mathrm{exp}(|z|^2/2)\underset{n_=0}{\overset{\mathrm{}}{}}\frac{z^n}{\sqrt{n!}}|n.$$
(39)
This state depends on a complex number $`z`$. Because of $`\overline{n}=\mathrm{coh}|\widehat{c}_𝐤^+\widehat{c}_𝐤|\mathrm{coh}=|z|^2`$ we may set
$$z=\sqrt{\overline{n}}\mathrm{exp}(\mathrm{i}\varphi ).$$
(40)
For the state (39) with Eqs. (38) and (40) we find
$$\mathrm{coh}|\widehat{\psi }(𝐫)|\mathrm{coh}=\sqrt{\overline{n}}\mathrm{exp}(\mathrm{i}\varphi )\phi _𝐤(𝐫)$$
(41)
for the expectation value of the field operator.
Using
$$\phi _0(𝐫)=\{\begin{array}{ccc}1/\sqrt{V_0}\hfill & & (𝐫V_0)\hfill \\ 0\hfill & & (𝐫V_0),\hfill \end{array}$$
(42)
for the lowest state, Eq. (41) becomes $`\mathrm{coh}|\widehat{\psi }|\mathrm{coh}=\sqrt{\overline{n}/V_0}\mathrm{exp}(\mathrm{i}\varphi )`$. So far we have reproduced Anderson’s construction.
We are now going to connect the phase ordered single particle functions (2) with coherent states. For this purpose we introduce the single particle functions
$`\phi _{𝐤,+}(𝐫)`$ $``$ $`\mathrm{exp}\left[\mathrm{i}(k_2y+k_3z)\right]\mathrm{exp}(+\mathrm{i}qx)`$ (43)
$`\phi _{𝐤,}(𝐫)`$ $``$ $`\mathrm{exp}\left[\mathrm{i}(k_2y+k_3z)\right]\mathrm{exp}(\mathrm{i}qx)`$ (44)
These functions shall be orthonormalized in the considered subvolume. Analogously to Eqs. (38) with (37) we introduce the $`n`$-particle states $`|n_+`$ and $`|n_{}`$ that correspond to $`\phi _{𝐤,+}`$ and $`\phi _{𝐤,}`$, respectively. The coherent state
$$|\mathrm{coh}=\frac{\mathrm{exp}(|z|^2/2)}{\sqrt{2}}\underset{n=0}{\overset{\mathrm{}}{}}\frac{z^n|n_+z^n|n_{}}{\sqrt{n!}}$$
(45)
with $`z=\sqrt{\overline{n}}\mathrm{exp}(\mathrm{i}\varphi )`$ yields then
$$\mathrm{coh}|\widehat{\psi }(𝐫)|\mathrm{coh}=\sqrt{\overline{n}}\phi _𝐤^{\mathrm{p}.\mathrm{o}.}(𝐫).$$
(46)
The phase variance of the coherent state (45) is given by
$$\mathrm{\Delta }\varphi _𝐤=\frac{1}{2\sqrt{\overline{n}_𝐤}}.$$
(47)
Here we have attached the index $`𝐤`$ again; the above discussion referred to a single $`𝐤`$ value. Eq. (47) is a well-known result for a coherent state of the standard form (39). For the state (45) we may write $`|\mathrm{coh}=(|\mathrm{coh}_+|\mathrm{coh}_{})/\sqrt{2}`$ in an obvious notation. The phase operator acts separately within each of the two sets of states, $`\{|n_+\}`$ or $`\{|n_{}\}`$. Moreover, the states of one of these sets are orthogonal to that of the other set (because of the functions (43) and (44) are mutually orthogonal). Therefore, the phase variance of the state (45) is one half of the sum of the phase variances of the states $`|\mathrm{coh}_+`$ and $`|\mathrm{coh}_{}`$ that are of the standard form (39). This means that we obtain the standard result (47) for our somewhat special coherent state (45), too.
The reason for establishing the relation between the phase ordered single particle functions (2) and coherent states (45) is the necessity to justify the relation (47) or, equivalently, Eq. (7). The relation (7) is the basis for the phase ordering entropy (8).
Instead of $`|\mathrm{coh}=(|\mathrm{coh}_+|\mathrm{coh}_{})/\sqrt{2}`$ we may also consider the coherent state $`(|\mathrm{coh}_++|\mathrm{coh}_{})/\sqrt{2}`$. This orthogonal state corresponds to a cosinus in Eq. (2) instead of a sinus. The considered correlation effect requires that only one of these sets of states is occupied. This $`2:1`$ reduction has been discussed in Section 2.2.
With respect to the condition $`n_𝐤1`$ we note that the average occupation numbers $`n_𝐤`$, Eq. (12), do not depend on the size $`V_0`$ of the subvolumes. The finite size implies, however, a finite spacing of the momentum values (App. B) that takes care of the finite number $`N_0=V_0/v`$ of atoms in one subvolume.
## Appendix B Phase coherence volume
We start by explaining why finite phase correlations require that the single particle functions (2) have to be localized within finite volumes. Then we show that the lowest single particle functions may extend over several of these subvolumes. This leads to the notion of a phase coherence volume for which we obtain $`V_{\mathrm{coh}}V_0/|t|^2`$, where $`V_0`$ is the size of one subvolume. We discuss this result in a number of points.
Two phase ordered states (2) with momentum $`𝐤`$ and $`𝐤^{}`$ are correlated if $`k_1=k_1^{}`$. The restriction $`k_1=k_1^{}`$ cancels one momentum sum in the expression for the energy. A momentum sum $`_k𝑑k/\mathrm{\Delta }k`$ is proportional to $`1/\mathrm{\Delta }k=V^{1/3}/2\pi `$ or to $`N^{1/3}`$. The correlation energy $`E_{\mathrm{corr}}`$ is, therefore, proportional to $`N^{2/3}`$ only. This is consistent with the observation that a phase ordering of the kind (2) can be obtained by physical boundary conditions at the wall of the considered volume. For an infinite volume such a surface effect vanishes ($`E_{\mathrm{corr}}/N1/N^{1/3}0`$). A finite correlation effect may, however, be obtained by dividing the total volume $`V`$ into $`V/V_0`$ subvolumes of size $`V_0`$ and by requiring physical boundary conditions at the walls of these subvolumes. For the whole system (including a factor $`V/V_0`$ for summing over all subvolumes) we obtain then a result of the form
$$\frac{E_{\mathrm{corr}}}{N}=\frac{w_0}{N_0^{1/3}},$$
(48)
where $`w_0<0`$ is some average strength of the attractive part of the interaction (for a more detailed discussion see Ref. ). This result shows that a finite correlation effect requires a finite size $`V_0=N_0v`$ of the subvolumes.
Finite volumes $`V_0`$ imply a finite spacing $`\mathrm{\Delta }k=\pi /V_0^{1/3}`$ of the momentum values. Alternatively one may start from a finite spacing and admit only the momentum values
$$q_n=q_0+n\mathrm{\Delta }k,(n=0,1,2,\mathrm{}),$$
(49)
in the single particle functions (2). The spacing (49) leads to a finite correlation energy even if the single particle functions are not localized. In a macroscopic system the possible $`q`$ values are, however, dense, and the entropy drives the particles into the occupation of all available states. This is the reason why the finite spacing (49) can be realized only for single particle functions localized within subvolumes of the size $`V_0=(\pi /\mathrm{\Delta }k)^3`$. In this sense, a finite spacing implies finite volumes.
This discussion shows that we may either start from finite volumes or from a finite $`\mathrm{\Delta }k`$ value. There is, however, the following difference: Starting from the finite spacing (49) we may admit a $`q_0`$ value below $`\mathrm{\Delta }k`$ without damaging the correlation energy. A single particle function with $`q_0<\mathrm{\Delta }k`$ can, however, not be localized within $`V_0`$.
In view of this observation we consider the following modified picture. Only the single particle functions (2) with $`q_{n1}`$ are localized within $`V_0`$, the single particle function $`\phi _0`$ with $`q_0`$ may have a larger extension. We present an estimate for the volume $`V_{\mathrm{coh}}`$ of the lowest single particle function $`\phi _0`$: Let $`V_{\mathrm{coh}}`$ be some multiple of $`V_0`$, that means $`V_{\mathrm{coh}}=WV_0`$. In this case, there are $`W`$ single particle states with $`q<\mathrm{\Delta }k`$ in the volume $`V_{\mathrm{coh}}`$ out of which only one (the one with $`q_0`$) is occupied. A redistribution of $`n_0`$ atoms over these $`W`$ states would increase the entropy by $`\mathrm{\Delta }sk_\mathrm{B}\mathrm{ln}(n_0)^W`$. At the same time these atoms would loose their correlation energy, $`\mathrm{\Delta }e_{\mathrm{corr}}n_0w_0/N_0^{1/3}`$. The stability condition $`T_\lambda \mathrm{\Delta }s<\mathrm{\Delta }e_{\mathrm{corr}}`$ yields the upper bound $`W<[w_0/(N_0^{1/3}k_\mathrm{B}T_\lambda )]n_0/\mathrm{ln}n_0n_0`$. Using $`n_0=n_{𝐤0}=\tau ^2t^2`$ leads to
$$V_{\mathrm{coh}}(t)=WV_0\frac{V_0}{t^2}.$$
(50)
We discuss this result in a number of points:
1. The lowest single particle functions $`\phi _0`$ extend over volumes of the size $`V_{\mathrm{coh}}`$. Within a volume $`V_{\mathrm{coh}}`$ the function $`\phi _0`$ defines the direction of the phase ordering ($`x`$ direction in Eq. (2)) and the phase $`\varphi `$. Therefore, $`V_{\mathrm{coh}}(t)`$ of Eq. (50) constitutes a phase coherence volume.
2. Strictly finite volumes for all particles would mean that the lower bound of the integral $`J(\tau )`$ in Eq. (25) would be $`\mathrm{\Delta }x`$ instead of zero. This would imply a cut of the logarithmic singularity. The replacement (24) may, however, be justified by a lower bound $`q_01/V_{\mathrm{coh}}`$ and an average over somewhat different $`q_0`$ and $`\mathrm{\Delta }k`$ values in different parts of the macroscopic system. A closer examination of this point might lead to a modification of the logarithmic singularity at very low $`|t|`$ values.
3. In the estimate leading to Eq. (50) we used the continuous part of the occupation number $`n_{𝐤0}`$ also for $`t<0`$. Using instead the condensate occupation number $`n_0(t)`$ leads to an infinite volume, $`V_{\mathrm{coh}}^{}=\mathrm{}`$. This is the adequate phase coherence volume for the potentially macroscopic range of a superfluid flow.
4. Different phase coherence volumes of size $`V_{\mathrm{coh}}`$ within the macroscopic system will, in general, correspond to different phase directions ($`x`$ direction in Eq. (2)). A specific phase direction implies a small anisotropy of the static structure function. In principle, the volumes $`V_{\mathrm{coh}}`$ and their increasing size for $`t0`$ should, therefore, be observable.
5. When approaching the lambda point from above the phase coherence volumes grow according to Eq. (50). At the lambda point the coherence volume becomes infinite. This picture has some similarity with a ferromagnetic system.
Below the lambda point the finite volume (50) refers to the phase ordering of real single particle functions of the form (2). At the same time there is a potentially infinite range of phase coherence of a superfluid current. The corresponding phase coherence volume $`V_{\mathrm{coh}}^{}=\mathrm{}`$ refers to the complex condensate wave $`\phi _0\mathrm{exp}[\mathrm{i}\varphi (𝐫)]`$.
6. Adjusting the correlation energy (48) to the experimental strength of the logarithmic singularity yields $`N_0^{1/3}50`$. The phase $`\varphi `$ should be approximately constant within $`V_0`$ (see App. A). It appears, therefore, tempting (but also speculative) to set $`|\varphi |_{\mathrm{max}}V_0^{1/3}`$ and to use the relation $`𝐮_\mathrm{s}=\mathrm{}\varphi /m`$ for the velocity of a potential superfluid flow. The maximum (or critical) velocity for such a flow would then be
$$v_{\mathrm{crit}}=\frac{\mathrm{}|\varphi |_{\mathrm{max}}}{m}\frac{\mathrm{}}{mV_0^{1/3}}=\frac{\mathrm{}}{mv^{1/3}}\frac{1}{N_0^{1/3}}0.8\frac{\mathrm{m}}{\mathrm{s}}.$$
(51)
In this way the finite size $`V_0`$ renders a possible connection between the “natural” velocity scale $`\mathrm{}/(mv^{1/3})`$ and realistic values for the critical velocity. Using Eq. (50) we obtain $`v_{\mathrm{crit}}|t|^{2/3}\mathrm{m}/\mathrm{s}`$ for the critical velocity near the transition point.
## Appendix C Evaluation of $`J(\tau )`$
We determine the critical ($`\tau 1`$) behavior of the integral
$$J(\tau )=_0^{\mathrm{}}𝑑xx^2\frac{\mathrm{ln}[\mathrm{exp}(x^2+\tau ^2)1]}{\mathrm{exp}(x^2+\tau ^2)1}.$$
(52)
Writing $`\mathrm{ln}[\mathrm{exp}(x^2+\tau ^2)1]=\mathrm{ln}[1\mathrm{exp}(x^2\tau ^2)]+x^2+\tau ^2`$ we obtain
$$J(\tau )=_0^{\mathrm{}}𝑑x\frac{x^2\mathrm{ln}[1\mathrm{exp}(x^2\tau ^2)]}{\mathrm{exp}(x^2+\tau ^2)1}+_0^{\mathrm{}}𝑑x\frac{x^2(x^2+\tau ^2)}{\mathrm{exp}(x^2+\tau ^2)1}.$$
(53)
The second integral yields a result that is of the form
$$R_0=\text{const.}+𝒪(\tau ^2).$$
(54)
In the following $`R_0`$ stands for any expression of this kind. By $`𝒪(\tau ^2)`$ we mean terms that are proportional to $`\tau ^2`$ or to higher powers of $`\tau `$.
In the first integral in Eq. (53) we use the following expansions into powers of $`y=x^2+\tau ^2`$ or $`\mathrm{exp}(y)`$:
$$\frac{\mathrm{ln}[1\mathrm{exp}(y)]}{\mathrm{exp}(y)1}=\{\begin{array}{cc}\frac{\mathrm{ln}(y)}{y}\frac{\mathrm{ln}(y)+1}{2}+\frac{y(2\mathrm{ln}(y)+7)}{24}\pm \mathrm{}\hfill & (y1)\\ \mathrm{exp}(2y)\frac{3\mathrm{exp}(3y)}{2}\frac{11\mathrm{exp}(4y)}{6}\pm \mathrm{}\hfill & (y1).\end{array}$$
(55)
We divide the integration into one part from 0 to 1 and another part from 1 to $`\mathrm{}`$, and insert the appropriate expansion (55). All terms except the one with $`\mathrm{ln}(y)/y`$ yield contributions of the form (54). We evaluate the remaining integral over $`x^2\mathrm{ln}(y)/y=\mathrm{ln}(y)\tau ^2\mathrm{ln}(y)/y`$,
$`J(\tau )`$ $`=`$ $`{\displaystyle _0^1}𝑑x\mathrm{ln}(x^2+\tau ^2)\tau ^2{\displaystyle _0^1}𝑑x{\displaystyle \frac{\mathrm{ln}(x^2+\tau ^2)}{x^2+\tau ^2}}+R_0`$ (56)
$`=`$ $`\pi \tau \tau ^2{\displaystyle _0^{\mathrm{}}}𝑑x{\displaystyle \frac{\mathrm{ln}(x^2+\tau ^2)}{x^2+\tau ^2}}+R_0.`$
The first integral yielded $`\pi \tau +R_0`$ (see number 2733 in Ref. ). Because of $`_1^{\mathrm{}}𝑑x\mathrm{ln}(x^2+\tau ^2)/(x^2+\tau ^2)=R_0`$ the upper limit of the remaining integral could be set equal to infinity. In this integral we substitute $`x=\tau z`$:
$`J(\tau )`$ $`=`$ $`\pi \tau \tau {\displaystyle _0^{\mathrm{}}}𝑑z{\displaystyle \frac{\mathrm{ln}(\tau ^2)}{1+z^2}}\tau {\displaystyle _0^{\mathrm{}}}𝑑z{\displaystyle \frac{\mathrm{ln}(1+z^2)}{1+z^2}}+R_0`$ (57)
$`=`$ $`J(0)\pi \tau \mathrm{ln}\tau +\pi \tau \left[1\mathrm{ln}(2)\right]+𝒪(\tau ^2).`$
The first integral is elementary, the second one may be found under number 4.295 in Ref. . A numerical integration yields $`J(0)1.183`$. |
warning/0001/astro-ph0001458.html | ar5iv | text | # The X-ray spectrum of a disk illuminated by ions
## 1 Introduction
The X-ray spectra $`F(E)`$ of AGN as well as galactic black hole candidates (BHC) in their hard states are characterized by a power law of index $`1`$ and a high energy cutoff $`E_\mathrm{c}`$ around 200 keV (e.g. Mushotzky et al. 1993, Ulrich et al. 1997, Cappi et al. 1997, Zdziarski et al. 1996). Such spectra are well known to be describable by Comptonization in an electron scattering layer of optical depth $`\tau 0.5`$ and temperature $`TE_\mathrm{c}/2`$. One of the classsical problems in X-ray astronomy is to explain why $`\tau `$ and $`T`$ should have just these values, with little variation between sources. Theoretical arguments can be given that Comptonization is in fact the most important interaction between matter and radiation at temperatures of 10–100 keV (e.g. Shakura and Sunyaev 1973), for the inferred radiation energy densities near accreting black holes, but this does not tell us what the the thickness and temperature of the interaction region are.
An optically thick accretion disk would produce spectra peaking at 1 keV and 10–100 eV for BHC and AGN, repectively. The conditions in BHC and AGN allow (at accretion rates well below Eddington) for a second form of accretion, an ion supported advection torus (Shapiro et al. 1976, Liang 1979, Ichimaru 19xx, Rees et al. 1982, Abaramowicz et al. 1988, Narayan and Yi 1994, Fabian and Rees 1995, Narayan et al. 1995, 1996). The ions in this flow are near their virial temperature, the electrons much cooler because of their weak interaction with the ions and their strong interaction with the radiation field. Such flows could produce, in principle, the kind of spectrum observed (Narayan et al. 1995, 1996), but more physics must be invoked to restrict the optical depth and temperature of the flow to the observed narrow ranges (cf. Haardt 1997, Maraschi and Haardt 1997).
### 1.1 Evaporating disks inside tori
Various geometries for the accretion flow near a black hole have been developed, for a review see Collin (1997). One of the possibilities is an ion supported advection torus coexisting with an optically thick accretion disk embedded in it (see fig. 1 in Collin). This possibility is attractive because the spectra of BHC in their high states show evidence of the simultaneous presence of an optically thick, thermal, accretion disk and a hotter component which produces a power law tail at higher photon energies (e.g. Mitsuda et al. 1984, Tanaka 1997).
Theoretically, one would expect exchange of both mass an energy to take place between the disk and the advection torus. Heating of the disk surface by the hot ion supported flow above would lead to an ‘evaporation’ of the disk surface, feeding mass into the torus. Such evaporation has been studied in detail for the case of disks in Cataclysmic Variable systems by Meyer and Meyer-Hofmeister (1994). In the inner regions of the disk the mass available in the disk is smallest, and the energy budget potentially available for evaporation largest. If the mass flow from disk into torus increases with the energy dissipation rate, and if a steady state develops, one could therefore envisage a structure consisting of three regions: an outer one in which only a geometrically thin optically thick disk is present, inside this a composite region with an evaporating disk inside a hot ion supported advection torus, and inside this a region in which only an ion supported flow exists because all disk mass has evaporated (Meyer and Meyer-Hofmeister 1994, Meyer-Hofmeister and Meyer 1999).
Depending on details of the processes of mass and energy exchange between disk and torus, the boundaries between these regions may vary. It is not necessary that the structure is steady. The model has, in principle, sufficient ingredients to allow for variability and may perhaps be developed further in the context of the various forms of variability seen in BHC.
### 1.2 Energy exchange between disk and torus
Energy exchange between disk and torus may be mediated by particles or by radiation. In its simplest form of interaction, the torus provides hard photons that are reprocessed by the disk, a model used extensively for interpretations of X-ray spectra (e.g. Wandel and Liang 1991, Reynolds et al. 1994, Petrucci and Henri 1997, Gilfanov et al. 1999 and references therein).
A more internally consistent model for this interaction is that of Haardt and Maraschi (1991). To simplify the discussion, assume that the accretion takes place predominantly through an ion torus or a hot ‘corona’ above the disk (this assumption can easily be relaxed). The radiation produced by the torus illuminates the disk below, which thermalizes it into an approximate blackbody spectrum. These (soft) photons are Comptonized in the hot torus. In this model, approximately half the energy comes out as soft radiation and half as Comptonized photons. It correctly predicts the slope of the spectrum, and fixes a relation between the temperature and optical depth of the Comptonizing layer. An assumption about the temperate of of Comptonizing region is still needed to get the right cutoff energy $`E_\mathrm{c}`$, though pair ceation limits the possible temperatures. For further developments of this model see Haardt (1997).
A second channel of energetic interaction is the hot protons in the torus with temperature near the virial temperature, $`T_\mathrm{p}T_\mathrm{v}160r_\mathrm{g}/r`$ MeV. At the distance dominating the energy release, $`r7r_\mathrm{g}`$, the protons thus have a temperature around 20 MeV. At this energy, they have a significant penetration depth into the cool disk. They are slowed down mainly by Coulomb interactions with the ensemble of electrons inside their Debye sphere. The ‘stopping depth’, expressed in terms of the corresponding Thomson optical depth, is
$$\tau _\mathrm{s}\frac{m_\mathrm{p}}{3m_\mathrm{e}\mathrm{ln}\mathrm{\Lambda }}\frac{\beta ^4}{\psi x\psi ^{}},$$
(1)
(e.g. Ryter et al., 1970) where $`\beta =v_\mathrm{z}/c(kT_\mathrm{p}/m_\mathrm{p}c^2)^{1/2}`$ is the vertical component of the proton velocity, $`\theta =kT/m_\mathrm{e}c^2`$ is a measure of the temperature of the heated layer, $`\mathrm{ln}\mathrm{\Lambda }20`$ is a Coulomb logarithm, $`x^2=\beta ^2/(2\theta )`$, and $`\psi `$ the error function
$$\psi =\frac{2}{\sqrt{\pi }}_0^xe^{x^2}dx.$$
(2)
This formula holds for nonrelativistic temperatures; the relativistic generalization has been given by Stepney (1983) and Stepney and Guilbert (1983). At low temperature $`kT<m_\mathrm{e}/m_\mathrm{p}kT_\mathrm{p}`$, $`x`$ is small and the factor involving the error function can be expanded. This yields
$`\tau _\mathrm{s}`$ $`{\displaystyle \frac{m_\mathrm{p}}{m_\mathrm{e}\mathrm{ln}\mathrm{\Lambda }}}\beta \theta ^{3/2},`$ (3)
$``$ $`\left({\displaystyle \frac{kT_\mathrm{p}}{50\mathrm{M}eV}}\right)^{1/2}\left({\displaystyle \frac{kT}{50\mathrm{k}eV}}\right)^{3/2}.`$ (4)
## 2 Comptonization in a layer heated by protons
### 2.1 Estimating the depth of the Comptonizing layer
Heating by protons yields a Comptonizing layer of thickness equal to the stopping depth $`\tau _\mathrm{s}`$. This depth is a function of the temperature in the layer, by (1). The temperature on the other hand is determined by the heating and cooling processes, so that a consistent calculation of heating and cooling will yield both the temperature and the optical depth of the layer. With a simple estimate, we can now show that this will yield $`\tau _\mathrm{s}`$ and $`T`$ in roughly the right range.
The cooling process in the layer is the inverse Compton process, i.e. the energy loss electrons experience as they scatter the soft photons from the cool disk below. We assume that these soft photons are all (or mostly) produced by thermalization of Comptonized photons from the heated layer, as in the model of Haardt and Maraschi (1991, hereafter HM). Since approximately half the Comptonized photons escape and the other half illuminates the thermalizing layer, the energy flux in the soft photons at the base of the layer must be about the same as that in the escaping Comptonized photons. Such a balance is possible only if the Comptonization is sufficiently strong. In terms of the Compton $`y`$-parameter $`y4\theta \tau _\mathrm{s}`$, it requires that $`y1`$. If the temperature is too low, the energy transfer from the electrons to the soft photons is too low and the layer heats up until $`y1`$, and vice versa. Since the $`y`$-parameter also determines the slope of the X-ray spectrum, the model yields a fixed spectral slope, which is in the range of the observed values. In HM the depth of the layer is a free parameter; in the present model, it is fixed by requiring the stopping depth $`\tau _\mathrm{s}`$ to be consistent with the resulting temperature $`T`$. A simple estimate is obtained by setting $`y=1`$, or $`\theta =1/(4\tau _\mathrm{s})`$, and inserting into (1). This yields, assuming again that the protons are near their virial temperature:
$$\tau _\mathrm{s}^{5/2}=\frac{m_\mathrm{p}}{8\sqrt{6}\mathrm{ln}\mathrm{\Lambda }m_\mathrm{e}}\left(\frac{r_\mathrm{g}}{r}\right)^{1/2},$$
(5)
or
$$\tau _\mathrm{s}1.3\left(\frac{7r_\mathrm{g}}{r}\right)^{1/5},$$
(6)
and
$$kTm_\mathrm{e}c^2/(4\tau _\mathrm{s})60\left(\frac{r}{7r_\mathrm{g}}\right)^{1/5}\mathrm{k}eV.$$
(7)
We conclude that proton illumination yields optical depths and temperatures in the right range, with only a weak dependence on the assumed distance from the black hole. Obviously, the estimate is rather crude, and more detailed calculations of the energy transfer from the protons to the electrons, as well as the Comptonization process are needed to test the model. In the following we make a first step in this direction, by means of a radiative transfer calculation.
### 2.2 Model problem
The aim of the calculation reported below is to compute the temperature as a function of depth, together with the emergent photon spectrum from a layer heated by protons who deposit their energy according to (1). A preliminary account of the calculation has been given in Spruit (1997).
We approximate the heating rate to be distributed uniformly over a layer with depth $`\tau _\mathrm{s}`$, where $`\tau _\mathrm{s}`$ is computed from the average temperature in this layer using eq. (1). I reality, the heating is somewhat non-uniform because of the dependence of the energy loss rate of the proton on both the proton velocity and the temperature. The approximation of a uniform energy input in the layer is deemed sufficient for the present exploratory purpose.
In this layer, we solve the radiative transfer equation iteratively together with the temperature $`T(\tau )`$, and the layer stopping depth $`\tau _\mathrm{s}`$ such that the heating is in balance with cooling by Comptonization of the soft photons.
The assumptions and simplifications that go into the model are as follows. In the heated layer, the only photon process considered is electron scattering (Comptonization). Below this, we assume that electron scattering continues to be the dominant process down to some depth $`\tau _\mathrm{b}`$. At $`\tau _\mathrm{b}`$, the downward photons are assumed to be absorbed and their energy reradiated upward as a black body spectrum. Thus, the gradual thermalization with depth by free-free processes is simplified by a step at depth $`\tau _\mathrm{b}`$. The value of $`\tau _\mathrm{s}`$ is determined from the proton velocity $`\beta =v/c`$ and the mean temperature in the heated layer. Obvious improvements are possible on these simplifications, by explicitly taking into account photon production/destruction process, an by a more accurate treatment of the energy loss of the ions as they penetrate into the disk.
The radiative transport part of the problem is simplified by reducing the angular dependences to a one-stream model: only vertically upward and downward moving photons are considered, and the electron distribution is similarly reduced from a 3-D to a one-dimensional Maxwellian distribution. For the scattering cross section and the electron Maxwellian the relativistic expressions are used. The one-stream simplification is made for programming convenience only: leaving out the full angular dependences leads to very simple expressions. Discretization by a reasonable number of angles would still yield a very modest problem in terms of computing time, and is an obvious next step to improve the calculations.
### 2.3 Numerical method
The transport equation is of the form (e.g. Rybicki and Lightman, 1976):
$`{\displaystyle \frac{\mathrm{d}n}{\mathrm{d}z}}=`$ $`{\displaystyle }\mathrm{d}^3𝐩{\displaystyle }\mathrm{d}^2𝛀{\displaystyle \frac{\mathrm{d}\sigma }{\mathrm{d}𝛀}}[f_\mathrm{e}(𝐩^{})n(\omega ^{})\times `$ (8)
$`(1+n(\omega ))f_\mathrm{e}(𝐩)n(\omega )(1+n(\omega ^{}))],`$ (9)
where $`n(\omega )`$ is the photon occupation number, $`f_\mathrm{e}`$ the electron momentum distribution, $`𝐩`$ (respectively $`𝐩^{}`$) the electron momentum, $`\omega `$ ($`\omega ^{}`$) the photon momentum vectors before (after) scattering, and $`𝛀`$ the scattering angle. The dependences of $`𝐩^{}`$ and $`\omega ^{}`$ on ($`𝐩,\omega ,𝛀`$) follow from the collision kinematics.
This equation is discretized in $`N_\omega `$ logarithmically spaced photon-energy bins. As photon energy scale we use $`\omega =h\nu /m_\mathrm{e}c^2`$. As depth scale we use the Thomson optical depth $`\tau _\mathrm{T}`$. The number of depth points is fixed, but the grid is stretched such that at each iteration the stopping depth $`\tau _\mathrm{s}`$ is located at the same grid point. Half the grid points are used for the heated layer ($`0<\tau <\tau _\mathrm{s}`$), the other half for the scattering layer below.
If $`n_i^+`$ and $`n_i^{}`$ are the occupation numbers of the upward and downward moving photons in bin $`i`$, the result is the set of $`2N_\omega `$ equations
$$\frac{\mathrm{d}}{\mathrm{d}\tau }n_i^\pm =n_i^\pm \underset{j}{}[(B_{ji}B_{ij})n_j^{}B_{ij}]+\underset{j}{}B_{ji}n_j^{},$$
(10)
where $`B_{ij}(T)`$ is the scattering cross section integrated over the appropriate frequency and electron momenta, for scattering from bin $`i`$ into bin $`j`$, in units of the Thomson cross section $`\sigma _\mathrm{T}`$. At the top boundary there are only outgoing photons:
$$n_i^{}=0(\tau =0),$$
(11)
while at the lower boundary the downward photons are thermalized into a blackbody spectrum of upward photons:
$$n_i^+=n_{\mathrm{B}B}(\omega _i)=1/[1\mathrm{exp}(\omega /\theta _\mathrm{b})](\tau =\tau _\mathrm{b}).$$
(12)
Here $`n_{\mathrm{B}B}`$ is the black-body occupation number at temperature $`\theta _\mathrm{b}`$. This temperature follows from the condition that the net energy flux at depth $`\tau _\mathrm{b}`$ vanishes:
$$\omega ^3[n^{}(\tau _\mathrm{b})n_{\mathrm{B}B}(\omega ,\theta _\mathrm{b})]=0.$$
(13)
This is a result of our assumption that the internal energy production rate in the cool disk can be neglected in comparison with the incident proton energy flux (generalizations are easily made by adding a term corresponding to the internal disspation in the disk).
The condition of energy balance between the assumed heating rate $`h(\tau )`$ and the Compton cooling by the soft photons is
$$\frac{\mathrm{d}F}{\mathrm{d}z}=h,$$
(14)
where the energy flux $`F`$ is given by
$$F=\underset{i}{}\omega ^3(n_i^+n_i^{}).$$
(15)
The equations (10) are discretized in optical depth by centered first order differences. The resulting set of nonlinear algebraic equations is solved by an iterative process. It turned out that full simultaneous linearization of the transfer equation, the boundary conditions and the energy equation had very poor convergence properties. Instead, an iteration was done in which only the transfer equation and the upper boundary condition were linearized, while the lower boundary condition and the energy equation were dealt with by a modified succesive-substitution process after each iteration of the transfer equation. Convergence, however, was still problematic for large optical depths and for cases where the assumed photon energy range extended too far beyond the cutoff energy. For the cases reported here, where the optical depth is not too large, on the order of 30–100 iterations were required for an accuracy of $`10^3`$ in luminosity.
## 3 Results
The parameters of the problem are the total energy flux $`F`$ (per unit surface area of the disk), the velocity of the incident protons, and the optical depth $`\tau _\mathrm{b}`$ of the thermalizing lower boundary. Assuming the protons to be thermal and virialized, their mean vertical velocity component $`\beta _z`$ as a function of the distance $`r`$ from the compact object is
$$\beta _z=\left(\frac{r_\mathrm{g}}{6r}\right)^{1/2}$$
(16)
The temperature as a function of depth is shown in figure 1 for a few values of $`F`$ for $`r/r_\mathrm{g}=7`$. These fluxes correspond to effective temperatures $`\theta _{\mathrm{e}ff}=(F/\sigma )^{1/4}/(m_\mathrm{e}c^2)`$ of $`\mathrm{3\hspace{0.17em}10}^4`$, $`10^3`$ and $`\mathrm{3\hspace{0.17em}10}^3`$, respectively. The emergent spectrum for these energy fluxes is shown in figure 4, for $`r/r_\mathrm{g}=7`$. The dependence of the spectrum on $`r/r_\mathrm{g}`$ at a fixed $`F`$ is shown in figure 3. The penetration depths are shown in figure 5. All of these cases were computed for $`\tau _\mathrm{b}=1`$.
The jump in temperature at $`\tau =\tau _\mathrm{s}`$ is a consequence of the fact that the heating rate jumps at $`\tau _\mathrm{s}`$. The photon field is continuous at $`\tau _\mathrm{s}`$ as shown in figure 2, as expected because of their diffusion through the scattering layer. In order to balance the energy input, the electrons have to loose enough energy by scattering soft photons. This requires them to be hot where the energy input rate is large, and causes the temperature to follow the heating rate. The jump will be smoothed if a more realistic model for the proton penetration is used. If the protons are not monoenergetic and unidirectional but taken from a thermal distribution, for example, the transition in heating rate will be smoother.
The similarity of the spectra, apart from shifts in amplitude and photon energy, is remarkable. The temperature of the heated layer increases only very weakly with increasing energy flux. The cutoff energy increases somewhat with distance from the hole, but again this is dependence is rather weak. On account of the modest optical depth of the layer, the spectrum shows a prominent contribution from unscattered soft photons from the reprocessing depth $`\tau _\mathrm{b}`$. This peak is smeared out somewhat when the spectra are convolved over distance from the hole, at a given accretion rate. This is shown in figure 6.
The final parameter of the problem is the depth of the thermalizing lower boundary, $`\tau _\mathrm{b}`$. As this depth is increased, the thickness of the passive scattering layer between the thermalizing depth and the heated surface layer increases. This has two effects on the spectrum. First, the strength of the ‘blue bump’ (consisting of unComptonized thermal photons from the lower boundary) decreases. Secondly the spectrum becomes flatter at the high energy end. This is shown in figure 7. With increasing depth of the thermalizing boundary, the spectrum becomes both smoother and flatter.
This can be understood as a combination of two factors. First, the low energy photons produced at $`z_\mathrm{b}`$ must travel through a scattering layer before reaching the base of the Comptonizing layer, and for this a gradient in photon density is needed. On arrival at the base of the Comptonizing layer, the soft spectrum is therefore not a black body any more, but a ‘diluted black body’. For a given soft flux, there are fewer photons but of higher energy than if the thermalization were to take place directly at the base of the Comptonizing layer. Secondly, fewer Comptonized photons reach the thermalization layer, since the scattering layer in between effectively reflects them to some degree. The combination of these effects causes the output spectrum to be harder than if the Comptonizing layer makes direct contact with the thermalizing surface.
How much of the hardening effect remains if other radiation processes than pure scattering are taken into account requires more detailed calculations taking into account free-free and bound-free transitions. The thermalization of the downward traveling Comptonized photons into soft photons is done by these processes, and represented only crudely by the thermalization at a fixed depth assumed here.
## 4 Discussion
With an admittedly somewhat simplified radiative transfer model I have shown that heating of a cool disk by protons from an ion supported advection torus produces X-ray spectra that are very reminiscent of the hard spectra of accreting black holes. Like the Maraschi and Haardt model (and for the same reason), it yields approximately the right spectral slope, but in addition, it also predicts approximately the right optical depth of the Comptonizing layer and the cutoff energy of the spectrum.
The temperature of the heated layer is insensitive to the energy flux, and stays around 40–60 keV. Instead of getting hotter at high energy flux, the incident energy is spent in upscattering a larger number of soft photons. As in the Maraschi and Haardt model, the reason for this lies in the energy balance condition. In order for the incident energy to be radiated as Comptonized flux, the Compton $`y`$-parameter has to be of of the order unity, and the temperature of the layer adjusts accordingly. What is new in the results presented here is that the optical depth of the Comptonizing layer also comes out naturally in the right range, due to the physics of Coulomb interaction between the virialized protons with the electrons in the cool disk. This process also is fairly insensitive to the proton temperature (within the relevant range), so that the emergent spectrum is only a weak function of distance from the accreting object (figure 4).
The agreement of the result with one of the most puzzling features of the hard X-ray spectra of accreting compact objects, viz. the uniformity of the spectral shape, makes it likely that ion illumination plays a major role in the physics of these objects.
### 4.1 Connection between disk and ADAF
The ion torus (ADAF) is underluminous for its accretion rate, because the ions do not have time to transfer their energy to the radiating electrons before being swallowed by the hole. The presence of a cool disk illuminated by the ions of the torus would increase the luminosity of the system. Since the illumination process produces a hard spectrum, like the torus itself, the presence of the cool disk can still be compatible with the observed hard spectra, but the accreting system would not be as underluminous as an ADAF without a cool disk. If any interaction at all takes place, the ADAF can not be underluminous by several orders of magnitude, as in some proposed ADAF applications.
The ion energy lost by illumination acts as a cooling agent on the hot ion torus, and limits the conditions under which it can exist. The hot ion supported flow has to be fed by a cool disk in some way or other, so that there must be at least a small region where the two coexist and the illumination process takes place. In the so-called intermediate state in black hole accreters (e.g. Rutledge etal), a soft and hard component coexist in the X-ray spectrum. If ion tori actually exist in these systems, this is observational evidence that disk and torus can coexist with both contributing significantly to the luminosity.
### 4.2 Lithium
An independent observational indication for ion illumination is the observation of high Li abundances in the secondary stars of X-ray binaries (Martín et al. 1992, 1994a, 1995). The energy of virialized protons hitting the cool disk at $`r=7r_\mathrm{g}`$ (where the gravitational energy release peaks for accretion onto a hole) is around 50 MeV, just in the range where Li production by spallation of CNO elements becomes efficient. Since the observed secondaries are of spectral types known to destroy Li on a rather short time scale, a significant continual source of Li is needed. As shown by Martín et al. (1994) the energetics of the accretion process is enough to explain the observed amount of Li on the secondary, if a fraction $`10^3`$ of the Li produced in the disk finds its way to the secondary (in the form of a disk wind, for example). Another consequence of ion illumination would be Li and Be production by He nuclei from the virialized flow reacting with He nuclei in the disk. These reactions peak around 50 Mev/nucleon, and are accompanied with emission of $`\gamma `$lines at 431 and 478 keV. It is possible that the $`\gamma `$lines observed sometimes around this energy (Gilfanov et al. 1991, Sunyaev et al. 1992) are another signature of ion illumination (Martín et al. 1994a,b).
## Acknowledgments
This work was done in the context of Human Capital and Mobility network ‘Accretion onto compact objects’, CHRX-CT93-0329. |
warning/0001/cond-mat0001323.html | ar5iv | text | # Condensate growth in trapped Bose gases
## I INTRODUCTION
The discovery of Bose-Einstein condensation (BEC) in trapped atomic vapors of <sup>87</sup>Rb , <sup>7</sup>Li , and <sup>23</sup>Na has initiated a period of intense experimental and theoretical activity. A great deal of information is now available about the equilibrium properties of these novel systems , but much remains to be understood about their nonequilibrium behavior. One of the most basic aspects concerns the nonequilibrium growth of the condensate which occurs in the process of cooling a nondegenerate trapped Bose gas to a final temperature below the BEC transition. This important problem was addressed even before the first observation of BEC in trapped atomic gases , and has interesting implications for the general problem of second-order phase transitions, from superfluidity in liquid <sup>4</sup>He to problems in cosmology .
Up to now, the most detailed study of condensate formation was carried out using a gas of <sup>23</sup>Na atoms confined within a highly anisotropic cigar-shaped trap . In these experiments, the sodium atoms were evaporatively cooled to a temperature just above the critical temperature and subsequently quenched by applying a rapid rf sweep. The latter step removes all, or at least a large fraction, of the atoms above a certain energy, after which the Bose gas relaxes to a new equilibrium state below the critical temperature. The growth of the condensate during the equilibration process was monitored using a nondestructive imaging technique which provided a direct measure of the size of the condensate as a function of time. In this way, the characteristic time scale for the growth of the condensate could be determined, and for the particular system studied, was found to be of the order of 100 ms.
A theoretical description of these experiments requires a theory that can account for the coupled nonequilibrium dynamics of both the noncondensed and condensed components of a trapped Bose gas, and includes in particular the collisional processes which transfer atoms between the two components. Thus far several such theories have been developed, which roughly speaking fall into two categories. One class of theories focuses on describing the dynamics of the average value of the order parameter for BEC, i.e., the condensate wave function, whereas the other incorporates also the fluctuations around this mean value. The latter of course, becomes important when the fluctuations are large compared to the mean value, i.e., close to the critical temperature. This is analogous to the situation in laser theory .
A theory that describes both the average value for the order parameter as well its fluctuations can be obtained in two, essentially equivalent ways. First, one can start from a master equation for the many-body density matrix and derive an equation of motion for the one-particle density matrix by means of a perturbative treatment of the interactions. This was the route followed by Gardiner and Zoller , in a series of papers. Second, one can use field-theoretic methods to obtain a nonperturbative Fokker-Planck equation that describes the nonequilibrium dynamics of the gas. This was the formulation developed by Stoof . These two approaches in principle yield a description of the nonequilibrium dynamics that is capable of obtaining the complete probability distribution for the order parameter.
Alternatively, a theory describing the dynamics of the mean-field value for the macroscopic wave function can also be obtained in a straightforward decoupling approach, which has been implemented by Kirkpatrick and Dorfmann , Proukakis et al. , Walser et al. , and in most detail for the trapped case by Zaremba, Nikuni, and Griffin . In this approach, one assumes the order parameter to be nonzero at all temperatures, and decouples the hierarchy of equations of motion that exists for the correlation functions of the second-quantized field operators. Thus one again obtains a perturbative expansion for these equations of motion.
The first quantitative calculations of condensate growth for trapped Bose gases were carried out by Gardiner et al. , and although good qualitative agreement with experiment was found, a number of quantitative discrepancies remained. For example, the reported experimental growth rates were up to a factor 30 larger than the initial theoretical results, and had a temperature dependence opposite to that predicted . By removing some simplifying approximations in subsequent calculations, the theoretical results were improved, but discrepancies of up to a factor of 3 still remained in some cases. From a purely theoretical point of view, one can attribute some of these discrepancies to the approximations made in the calculations. First, the dynamics of the noncondensate was to a large extent neglected. Although the time evolution of the occupancy of low-lying states was included in the simulations, the high energy states were represented by an equilibrium particle reservoir having a fixed chemical potential. This latter assumption is inconsistent with the nonequilibrium initial state established by the experimental quench procedure. Second, the effect of the mean field of the condensate on the noncondensate was included rather crudely by a linear rescaling of the low-energy density of states of the noncondensed atoms.
Our aim in the present paper is to improve on these calculations by taking fully into account the relaxational dynamics of the thermal, or noncondensed, component that takes place in the presence of the mean field of the condensate. We do this by starting from the above mentioned theories describing the growth process, which provide us with a nonlinear Schrödinger equation for the condensate and a kinetic equation for the thermal component. This coupled set of equations is still difficult to deal with and a number of physically motivated approximations are made to simplify the problem. We assume that the condensate grows adiabatically, having an equilibrium spatial distribution determined by the instantaneous number of atoms in the condensate. This assumption is also made in earlier work . The noncondensate is treated by solving a semiclassical Boltzmann equation in the ergodic approximation, which again has been used previously by numerous authors . These assumptions allow us to obtain numerically a detailed description of the growth of the condensate, including the effects of both the dynamics of the thermal cloud and its mean-field interaction with the condensate.
The paper is organized as follows. In Sec. II we summarize the theory of the nonequilibrium dynamics of a trapped Bose gas as developed previously. In Sec. III we introduce the central assumption, the ergodic approximation, that allows us to numerically solve the Boltzmann equation. In addition, we briefly discuss the adiabatic approximation for the condensate. In Sec. IV we treat in some detail particle number and energy conservation. Sec. V introduces the Thomas-Fermi approximation and gives some analytical results for the density of states and other quantities of interest. The numerical solution of our kinetic equations is discussed in Sec. VI and our results for the growth of a condensate are presented. We end in Sec. VII with a discussion and an outlook.
## II NONEQUILIBRIUM DYNAMICS
As discussed in the previous section, the nonequilibrium dynamics of a trapped Bose gas is governed by a set of equations for the condensate and noncondensate components. These equations have been presented in various forms , but they all describe the coherent dynamics of the gas due to mean-field interactions, as well as the incoherent dynamics associated with atomic collisions. The equations and notation we use are taken from Refs. .
The noncondensate is treated using a semiclassical Boltzmann equation for the phase space distribution function $`f(𝐫,𝐩,t)`$. This semiclassical description is justified when the largest level spacing in the external trapping potential is small compared to the thermal excitation energy. Moreover, mean-field interactions are included at the level of the Hartree-Fock approximation. In this situation, the quantum kinetic equation for the thermal excitations takes the form
$`{\displaystyle \frac{f(𝐫,𝐩,t)}{t}}+{\displaystyle \frac{𝐩}{m}}\mathbf{}f(𝐫,𝐩,t)\mathbf{}U(𝐫,t)\mathbf{}_𝐩f(𝐫,𝐩,t)=C_{12}[f]+C_{22}[f].`$ (1)
Here, the effective potential $`U(𝐫,t)U_{\mathrm{ext}}(𝐫)+2g[n_c(𝐫,t)+\stackrel{~}{n}(𝐫,t)]`$ is the sum of the external trapping potential $`U_{\mathrm{ext}}`$ and the self-consistent Hartree-Fock mean field. The latter is determined by the condensate density $`n_c(𝐫,t)`$, defined below, and the noncondensate density $`\stackrel{~}{n}(𝐫,t)`$ given by
$`\stackrel{~}{n}(𝐫,t)={\displaystyle \frac{d𝐩}{(2\pi \mathrm{})^3}f(𝐫,𝐩,t)}.`$ (2)
As usual, we treat the interactions in the $`s`$-wave approximation which results in the bare interaction being replaced by a contact interaction with an effective coupling constant $`g=4\pi \mathrm{}^2a/m`$ proportional to the $`s`$-wave scattering length $`a`$. The effective coupling constant is in fact equal to the two-body $`T`$-matrix, and to emphasize this connection, it is denoted by $`T^{2B}`$ in some works . The collision terms appearing in Eq. (1) are given by
$`C_{22}[f]`$ $``$ $`{\displaystyle \frac{4\pi }{\mathrm{}}}g^2{\displaystyle \frac{d𝐩_2}{(2\pi \mathrm{})^3}\frac{d𝐩_3}{(2\pi \mathrm{})^3}\frac{d𝐩_4}{(2\pi \mathrm{})^3}(2\pi \mathrm{})^3\delta (𝐩+𝐩_2𝐩_3𝐩_4)}`$ (4)
$`\times \delta (E+E_2E_3E_4)\left[(1+f)(1+f_2)f_3f_4ff_2(1+f_3)(1+f_4)\right],`$
$`C_{12}[f]`$ $``$ $`{\displaystyle \frac{4\pi }{\mathrm{}}}g^2n_c{\displaystyle \frac{d𝐩_2}{(2\pi \mathrm{})^3}\frac{d𝐩_3}{(2\pi \mathrm{})^3}\frac{d𝐩_4}{(2\pi \mathrm{})^3}(2\pi \mathrm{})^3\delta (m𝐯_c+𝐩_2𝐩_3𝐩_4)}`$ (7)
$`\times \delta (E_c+E_2E_3E_4)(2\pi \mathrm{})^3[\delta (𝐩𝐩_2)\delta (𝐩𝐩_3)\delta (𝐩𝐩_4)]`$
$`\times [(1+f_2)f_3f_4f_2(1+f_3)(1+f_4)],`$
with $`ff(𝐫,𝐩,t)`$, and $`f_if(𝐫,𝐩_i,t)`$. We note that Eq. (7) takes into account the fact that a condensate atom locally has an energy
$`E_c(𝐫,t)=\mu _c(𝐫,t)+{\displaystyle \frac{1}{2}}m𝐯_c^2(𝐫,t),`$ (8)
a momentum $`m𝐯_c(𝐫,t)`$, and a chemical potential $`\mu _c(𝐫,t)`$. These quantities are defined explicitly below. In addition, the energy of a noncondensate atom in the Hartee-Fock approximation is
$`E(𝐫,𝐩,t)={\displaystyle \frac{𝐩^2}{2m}}+U(𝐫,t).`$ (9)
The energy variables $`E_i`$ appearing in Eqs. (4) and (7) are defined as $`E_i=E(𝐫,𝐩_i,t)`$.
In contrast to the thermal cloud, the dynamics of the condensate is determined by a time-dependent dissipative nonlinear Schrödinger equation ,
$`i\mathrm{}{\displaystyle \frac{\mathrm{\Psi }(𝐫,t)}{t}}`$ $`=`$ $`\left\{{\displaystyle \frac{\mathrm{}^2\mathbf{}^2}{2m}}+U_{\mathrm{ext}}(𝐫)+g\left[2\stackrel{~}{n}(𝐫,t)+n_c(𝐫,t)\right]{\displaystyle \frac{}{}}iR(𝐫,t)\right\}\mathrm{\Psi }(𝐫,t),`$ (10)
where the dissipative term, i.e., $`R(𝐫,t)`$, is given by
$`R`$ $``$ $`{\displaystyle \frac{\mathrm{}}{2n_c}}{\displaystyle \frac{d𝐩}{(2\pi \mathrm{})^3}C_{12}[f]}`$ (11)
$`=`$ $`{\displaystyle \frac{g^2}{(2\pi )^5\mathrm{}^6}}{\displaystyle \underset{i=1,4}{}d𝐩_i\delta (m𝐯_c+𝐩_2𝐩_3𝐩_4)}`$ (14)
$`\times \delta (E_c+E_2E_3E_4)[\delta (𝐩_1𝐩_2)\delta (𝐩_1𝐩_3)\delta (𝐩_1𝐩_4)]`$
$`\times [(1+f_2)f_3f_4f_2(1+f_3)(1+f_4)].`$
The appearance of this dissipative term in Eq. (10) is a consequence of the collisional processes, described by $`C_{12}`$, which have the effect of transferring particles between the condensate and noncondensate. The dissipative term is needed in order to ensure overall particle number conservation of the entire system. At a more fundamental level, the condensate wave function is determined by taking the expectation value of a Bose field with respect to a probability distribution that satisfies the Fokker-Planck equation mentioned previously . The Langevin equation one derives within this formulation has the form of a dissipative nonlinear Schrödinger equation with a noise term. It only reduces to Eq. (10) in a mean-field approximation. This points to the need for exercising care in interpreting the order parameter occuring in the nonlinear Schrödinger equation as the condensate wave function. Due to the underlying $`U(1)`$ gauge invariance associated with strict particle number conservation, the expectation value is in fact always equal to zero as a result of the diffusion of the global phase of the condensate wave function. In first instance this effect can be neglected and we are then effectively treating the system as if the $`U(1)`$ gauge invariance is explicitly broken.
It is convenient to rewrite Eq. (10) in terms of amplitude and phase variables defined by $`\mathrm{\Psi }(𝐫,t)=\sqrt{n_c(𝐫,t)}\mathrm{exp}[i\theta (𝐫,t)]`$. Substituting this form of the wave function into Eq. (10), we obtain
$`{\displaystyle \frac{n_c(𝐫,t)}{t}}+\mathbf{}\left[𝐯_c(𝐫,t)n_c(𝐫,t)\right]={\displaystyle \frac{2}{\mathrm{}}}R(𝐫,t)n_c(𝐫,t),`$ (15)
and
$`m{\displaystyle \frac{𝐯_c(𝐫,t)}{t}}+\mathbf{}\left[\mu _c(𝐫,t)+{\displaystyle \frac{m𝐯_c(𝐫,t)^2}{2}}\right]=0.`$ (16)
Here, we have defined the local chemical potential and superfluid velocity by
$`\mu _c(𝐫,t)=U_{\mathrm{ext}}(𝐫)+g\left[n_c(𝐫,t)+2\stackrel{~}{n}(𝐫,t)\right]{\displaystyle \frac{\mathrm{}^2}{2m}}{\displaystyle \frac{\mathbf{}^2\sqrt{n_c(𝐫,t)}}{\sqrt{n_c(𝐫,t)}}},`$ (17)
and
$`𝐯_c(𝐫,t)={\displaystyle \frac{\mathrm{}}{m}}\mathbf{}\theta (𝐫,t),`$ (18)
respectively. It can easily be shown that this set of equations for the condensate and thermal cloud is consistent with the conservation of the total number of particles in the system .
## III ERGODIC APPROXIMATION
Our main objective in this paper is to apply the kinetic theory formulated above to the problem of condensate formation. In order to make progress we introduce a number of additional approximations. The first and most essential, is the assumption of ergodicity which has been widely used in the literature on kinetic theory. This assumes that equilibration of atoms within one energy level occurs on a much shorter time scale than equilibration of atoms between different energy levels. With this assumption, all points in phase space having the same energy are equally probable, and the distribution function therefore only depends on the phase space variables through the energy variable $`E(𝐫,𝐩,t)`$, i.e., $`f(𝐫,𝐩,t)g(E(𝐫,𝐩,t),t)`$. In equilibrium this is certainly correct, but the assumption requires justification for any particular nonequilibrium application. Unfortunately, we are not aware of any explicit checks that have been made which might indicate that the assumption is correct for the situations we wish to consider. Nevertheless, it appears to be physically reasonable that for quantities that vary on a time scale of the order of several collision times, the approximation is sufficiently accurate.
The ergodic approximation allows us to derive a simplified kinetic equation for the energy distribution function $`g(ϵ,t)`$. This is accomplished by means of the relation
$$\rho (ϵ,t)g(ϵ,t)\frac{d𝐫d𝐩}{(2\pi \mathrm{})^3}\delta (ϵE(𝐫,𝐩,t))f(𝐫,𝐩,t),$$
(19)
which shows that the phase-space projection defined on the right-hand side yields the product of $`g(ϵ,t)`$ and the density of states
$`\rho (ϵ,t)`$ $`=`$ $`{\displaystyle \frac{d𝐫d𝐩}{(2\pi \mathrm{})^3}\delta (ϵE(𝐫,𝐩,t))}`$ (20)
$`=`$ $`{\displaystyle \frac{m^{3/2}}{\sqrt{2}\pi ^2\mathrm{}^3}}{\displaystyle _{Uϵ}}𝑑𝐫\sqrt{ϵU(𝐫,t)}.`$ (21)
We note that the density of states is defined on the variable energy range $`U_{\mathrm{min}}(t)ϵ<\mathrm{}`$ where $`U_{\mathrm{min}}(t)`$ is the minimum value of $`U(𝐫,t)`$ at time $`t`$. The time dependence of the density of states is one of the aspects distinguishing the present development from previous work .
We now apply the phase-space projection to the kinetic equation in Eq. (1). As a result of this operation, the streaming terms in the Boltzmann equation, i.e., the second and third terms on the left-hand side of Eq. (1), cancel each other. Only the projection of the time-derivative term survives. This results in
$$\frac{d𝐫d𝐩}{(2\pi \mathrm{})^3}\delta (ϵE(𝐫,𝐩,t))\frac{f(𝐫,𝐩,t)}{t}=\rho (ϵ,t)\frac{g(ϵ,t)}{t}+\rho _\mathrm{w}(ϵ,t)\frac{g(ϵ,t)}{ϵ}.$$
(22)
Here, we have introduced a weighted density of states
$`\rho _\mathrm{w}(ϵ,t)`$ $`=`$ $`{\displaystyle \frac{d𝐫d𝐩}{(2\pi \mathrm{})^3}\delta (ϵE(𝐫,𝐩,t))\frac{U(𝐫,t)}{t}}`$ (23)
$`=`$ $`{\displaystyle \frac{m^{3/2}}{\sqrt{2}\pi ^2\mathrm{}^3}}{\displaystyle _{Uϵ}}𝑑𝐫\sqrt{ϵU(𝐫,t)}{\displaystyle \frac{U(𝐫,t)}{t}}.`$ (24)
This quantity depends explicitly on the time derivative of the noncondensate potential which in turn is determined by the time derivatives of both the condensate and noncondensate densities. Some formal details regarding its evaluation are given in the Appendix.
Noting that
$$\frac{\rho (ϵ,t)}{t}=\frac{\rho _\mathrm{w}(ϵ,t)}{ϵ},$$
(25)
Eq. (22) can be written as
$$\frac{d𝐫d𝐩}{(2\pi \mathrm{})^3}\delta (ϵE(𝐫,𝐩,t))\frac{f(𝐫,𝐩,t)}{t}=\frac{}{t}\left(\rho g\right)+\frac{}{ϵ}\left(\rho _\mathrm{w}g\right).$$
(26)
We thus arrive at the projected kinetic equation
$$\frac{}{t}\left(\rho g\right)+\frac{}{ϵ}\left(\rho _\mathrm{w}g\right)=I_{12}+I_{22},$$
(27)
where the phase space projections of the collision integrals are defined as
$`I_{12}(ϵ,t)`$ $``$ $`{\displaystyle \frac{d𝐫d𝐩}{(2\pi \mathrm{})^3}\delta (ϵE(𝐫,𝐩,t))C_{12}[f]}`$ (29)
$`I_{22}(ϵ,t)`$ $``$ $`{\displaystyle \frac{d𝐫d𝐩}{(2\pi \mathrm{})^3}\delta (ϵE(𝐫,𝐩,t))C_{22}[f]}.`$ (30)
The result in Eq. (27) is the kinetic equation that we solve numerically.
We now derive in some detail explicit expressions for the collision integrals in Eq. (III). Although an expression for $`I_{22}`$ was given in earlier work , we present here an alternative derivation which can also be adapted to the case of the $`I_{12}`$ collision integral. For the $`I_{22}`$ collision integral we have
$`I_{22}(ϵ_1,t)`$ $`=`$ $`{\displaystyle \frac{4\pi g^2}{(2\pi )^9\mathrm{}^{10}}}{\displaystyle 𝑑ϵ_2𝑑ϵ_3𝑑ϵ_4\delta (ϵ_1+ϵ_2ϵ_3ϵ_4)}`$ (34)
$`\times \left[(1+g_1)(1+g_2)g_3g_4g_1g_2(1+g_3)(1+g_4)\right]`$
$`\times {\displaystyle }d𝐫\left({\displaystyle \underset{i=1,4}{}}{\displaystyle }d𝐩_i\right)\delta (𝐩_1+𝐩_2𝐩_3𝐩_4)`$
$`\times \delta (ϵ_1E_1)\delta (ϵ_2E_2)\delta (ϵ_3E_3)\delta (ϵ_4E_4),`$
where we have introduced the short-hand notation $`g_i=g(ϵ_i,t)`$. We consider first the momentum integrals in Eq. (34) which, with the replacement $`𝐩_3𝐩_3`$ and $`𝐩_4𝐩_4`$, can be written as
$`J_{22}`$ $``$ $`{\displaystyle 𝑑𝐩_1𝑑𝐩_2𝑑𝐩_3𝑑𝐩_4\delta (𝐩_1+𝐩_2+𝐩_3+𝐩_4)}`$ (36)
$`\times \delta (ϵ_1E_1)\delta (ϵ_2E_2)\delta (ϵ_3E_3)\delta (ϵ_4E_4)`$
$`=`$ $`{\displaystyle \frac{d𝝃}{(2\pi )^3}\underset{i=1,4}{}𝑑𝐩_ie^{i𝝃𝐩_i}\delta (ϵ_iE_i)}.`$ (37)
In obtaining this expression, we have introduced a Fourier representation of the momentum conserving delta function. Performing the integrals in Eq. (36) with respect to the momentum variables, we obtain
$$J_{22}=(4\pi m)^4\left(\underset{i=1}{\overset{4}{}}\theta (ϵ_iU)\right)\frac{d𝝃}{(2\pi )^3}\frac{\mathrm{sin}\xi p_1\mathrm{sin}\xi p_2\mathrm{sin}\xi p_3\mathrm{sin}\xi p_4}{\xi ^4},$$
(38)
where now it is understood that $`p_i=\sqrt{2m(ϵ_iU)}`$. The product of theta functions can be replace by $`\theta (ϵ_{\mathrm{min}}U)`$, with $`ϵ_{\mathrm{min}}`$ the minimum value of the four energy variables. Performing the remaining integral with respect to the $`𝝃`$ variable, we find
$`J_{22}=(2\pi )^3m^4\theta (ϵ_{\mathrm{min}}U)`$ $`[`$ $`|p_1p_2+p_3+p_4||p_1p_2+p_3p_4|`$ (39)
$`+`$ $`|p_1p_2p_3p_4||p_1p_2p_3+p_4|`$ (40)
$`+`$ $`|p_1+p_2+p_3p_4||p_1+p_2p_3p_4|`$ (41)
$`+`$ $`|p_1+p_2p_3+p_4||p_1+p_2+p_3+p_4|].`$ (42)
This expression is valid for arbitrary values of the momenta but simplifies when energy conservation is taken into account. Since the energy conserving delta function $`\delta (ϵ_1+ϵ_2ϵ_3ϵ_4)`$ in Eq. (34) imposes the constraint $`p_1^2+p_2^2=p_3^2+p_4^2`$, Eq. (39) can be reduced to
$$J_{22}=4(2\pi )^3m^4\theta (ϵ_{\mathrm{min}}U)\sqrt{2m(ϵ_{\mathrm{min}}U)}.$$
(43)
Substituting this expression for $`J_{22}`$ into Eq. (34), we finally obtain
$`I_{22}(ϵ_1,t)`$ $`=`$ $`{\displaystyle \frac{m^3g^2}{2\pi ^3\mathrm{}^7}}{\displaystyle 𝑑ϵ_2𝑑ϵ_3𝑑ϵ_4\rho (ϵ_{\mathrm{min}})\delta (ϵ_1+ϵ_2ϵ_3ϵ_4)}`$ (45)
$`\times \left[(1+g_1)(1+g_2)g_3g_4g_1g_2(1+g_3)(1+g_4)\right],`$
where we have used the definition of the density of states in Eq. (20). This is precisely the result obtained by Snoke and Wolfe and Luiten, Reynolds, and Walraven , using a different method. We note that if all energies are expressed in units of $`\mathrm{}\overline{\omega }`$, the $`I_{22}`$ integral has an overall factor of $`(a/l)^2\overline{\omega }`$ where $`l=\sqrt{\mathrm{}/m\overline{\omega }}`$ is the average harmonic oscillator length. This factor defines a characteristic time which can be used as the time unit in the simulations.
The $`I_{12}`$ collision integral can be dealt with in a similar way if the superfluid velocity $`𝐯_c`$ in Eq. (7) is set to zero. The validity of this approximation follows from our assumption that the condensate grows adiabatically. The magnitude of the superfluid velocity $`𝐯_c`$ is then typically of the order of $`\dot{R}(t)`$, where $`R(t)`$ is the radius of the condensate. This velocity is small compared to the characteristic velocities $`𝐩/m\sqrt{2k_BT/m}`$ of the thermal atoms participating in a collision, which justifies the neglect of $`m𝐯_c`$ in Eq. (7). The expression for $`I_{12}`$ then reads
$`I_{12}(ϵ_1,t)`$ $`=`$ $`{\displaystyle \frac{4\pi g^2}{(2\pi )^6\mathrm{}^7}}{\displaystyle 𝑑𝐫n_c(𝐫,t)𝑑ϵ_2𝑑ϵ_3𝑑ϵ_4\delta (E_c(𝐫,t)+ϵ_2ϵ_3ϵ_4)}`$ (48)
$`\times \left[\delta (ϵ_1ϵ_2)\delta (ϵ_1ϵ_3)\delta (ϵ_1ϵ_4)\right]\left[(1+g_2)g_3g_4g_2(1+g_3)(1+g_4)\right]`$
$`\times {\displaystyle }d𝐩_2{\displaystyle }d𝐩_3{\displaystyle }d𝐩_4\delta (𝐩_2𝐩_3𝐩_4)\delta (ϵ_2E_2)\delta (ϵ_3E_3)\delta (ϵ_4E_4).`$
If we now define $`J_{12}`$ analogously to $`J_{22}`$, we have
$`J_{12}`$ $``$ $`{\displaystyle 𝑑𝐩_2𝑑𝐩_3𝑑𝐩_4\delta (𝐩_2+𝐩_3+𝐩_4)\delta (ϵ_2E_2)\delta (ϵ_3E_3)\delta (ϵ_4E_4)}`$ (49)
$`=`$ $`{\displaystyle \frac{d𝝃}{(2\pi )^3}\underset{i=2}{\overset{4}{}}𝑑𝐩_ie^{i𝝃𝐩_i}\delta (ϵ_iE_i)}`$ (50)
$`=`$ $`(4\pi m)^3\theta (ϵ_{\mathrm{min}}U){\displaystyle \frac{d𝝃}{(2\pi )^3}\frac{\mathrm{sin}\xi p_2\mathrm{sin}\xi p_3\mathrm{sin}\xi p_4}{\xi ^3}}`$ (51)
$`=`$ $`8\pi ^2m^3\theta (ϵ_{\mathrm{min}}U)S(p_2,p_3,p_4),`$ (52)
where $`ϵ_{\mathrm{min}}`$ is the minimum value of $`ϵ_2`$, $`ϵ_3`$, and $`ϵ_4`$, and
$`S(p_2,p_3,p_4){\displaystyle \frac{1}{2}}[\mathrm{sgn}(p_2+p_3p_4)`$ $`+`$ $`\mathrm{sgn}(p_2p_3+p_4)`$ (53)
$`\mathrm{sgn}(p_2+p_3+p_4)`$ $``$ $`\mathrm{sgn}(p_2p_3p_4)].`$ (54)
Note that this is a boolean function which takes on values of 0 and 1. Inserting the expression for $`J_{12}`$ into Eq. (48), we finally obtain for the $`I_{12}`$ collision integral the result
$`I_{12}(ϵ_1,t)`$ $`=`$ $`{\displaystyle \frac{m^3g^2}{2\pi ^3\mathrm{}^7}}{\displaystyle 𝑑ϵ_2𝑑ϵ_3𝑑ϵ_4\left[\delta (ϵ_1ϵ_2)\delta (ϵ_1ϵ_3)\delta (ϵ_1ϵ_4)\right]}`$ (57)
$`\times \left[(1+g_2)g_3g_4g_2(1+g_3)(1+g_4)\right]`$
$`\times {\displaystyle _{Uϵ_{\mathrm{min}}}}d𝐫n_c(𝐫,t)S(p_2,p_3,p_4)\delta (E_c(𝐫,t)+ϵ_2ϵ_3ϵ_4).`$
A comparison of this expression with $`I_{22}`$ in Eq. (45) shows that the remaining spatial integral acts as an effective density of states for scattering into the condensate. It can be evaluated analytically in the Thomas-Fermi approximation for the condensate, as shown in Sec. V. The kinetic equation in Eq. (27) and the projected collision integrals in Eqs. (45) and (57) are the main results of this section.
Before closing this section we point out one difficulty encountered when a numerical solution of Eq. (27) is attempted. As discussed following Eq. (20), the time dependence of the mean field potential $`U(𝐫,t)`$ implies that the density of states in Eq. (20), and hence the energy distribution function $`g(ϵ,t)`$, are defined on a variable energy range. To eliminate this variation, it is convenient to introduce the shifted energy variable
$$\overline{ϵ}ϵU_{\mathrm{min}}(t),$$
(58)
which leads to a fixed energy range $`0\overline{ϵ}<\mathrm{}`$. The density of states in terms of this new energy variable is given by $`\rho (ϵ,t)=\rho (\overline{ϵ}+U_{\mathrm{min}},t)\overline{\rho }(\overline{ϵ},t)`$. With $`\overline{ϵ}`$ and $`t`$ as independent variables, the kinetic equation in Eq. (27), can be rewritten as
$$\frac{}{t}\left(\overline{\rho }\overline{g}\right)+\frac{}{\overline{ϵ}}\left(\overline{\rho }_\mathrm{w}\overline{g}\right)=\overline{I}_{12}+\overline{I}_{22}.$$
(59)
Here, both $`\overline{\rho }(\overline{ϵ},t)`$ and $`\overline{\rho }_\mathrm{w}(\overline{ϵ},t)`$ are defined by making the replacement
$$U(𝐫,t)U(𝐫,t)U_{\mathrm{min}}(t)\overline{U}(𝐫,t),$$
(60)
in Eqs. (20) and (23). Similarly, from the definition of the collision integrals in Eqs. (4) and (7), it can also be seen that the change of energy variable leads to the replacement of $`U(𝐫,t)`$ by $`\overline{U}(𝐫,t)`$ in this case as well. Thus, the final kinetic equation in terms of the $`\overline{ϵ}`$ variable is unchanged in form from the original equation. We will henceforth drop the overbar on the functions defined in terms of $`\overline{ϵ}`$, with the understanding that the shifted potential $`\overline{U}(𝐫,t)`$ is to be used wherever the potential appears in the original expressions. The possibility of using a fixed energy range in the solution of the kinetic equation simplifies the numerical calculations considerably.
## IV COLLISIONAL INVARIANTS
In this section we explicitly consider two important quantities that should be conserved as the Bose gas condenses and equilibrates, namely, the total number of particles and the total energy of the trapped Bose gas. Together, they determine the final equilibrium state of the Bose-condensed gas, i.e., the number of particles in the condensate, its chemical potential, and the temperature of the vapor.
### A Particle Number Conservation
The time rate of change of the total number of particles consists of the time rate of change of the number of condensed particles plus the time rate of change of the number of noncondensed particles. Because the number of noncondensed particles is given by $`\stackrel{~}{N}(t)=(2\pi \mathrm{})^3𝑑𝐫𝑑𝐩f(𝐩,𝐫,t)=𝑑ϵ\rho (ϵ)g(ϵ)`$, the time rate of change of $`\stackrel{~}{N}(t)`$ can be found by integrating Eq. (27) over energy. We thus find that
$$\frac{\stackrel{~}{N}(t)}{t}=𝑑ϵI_{12}(ϵ,t),$$
(61)
where it is easily checked from Eq. (45) that
$`{\displaystyle 𝑑ϵI_{22}(ϵ,t)}=0.`$ (62)
Note that we have assumed here that $`lim_{ϵU_{\mathrm{min}}}\rho _\mathrm{w}(ϵ)g(ϵ)=0`$. A finite limiting value can arise if $`g(ϵ)`$ approaches an equilibrium Bose distribution with a chemical potential $`\mu =U_{\mathrm{min}}`$ at long times, together with a weighted density of states which depends linearly on $`ϵU_{\mathrm{min}}`$ for energies close to $`U_{\mathrm{min}}`$. However, at any finite time in the growth process, it is safe to use the zero limiting value. This is always the case when the equilibrium chemical potential lies below $`U_{\mathrm{min}}`$.
To get the time rate of change of the total number of condensate particles, we integrate the continuity equation, Eq. (15), over space to find
$`{\displaystyle \frac{N_c}{t}}`$ $`=`$ $`{\displaystyle \frac{2}{\mathrm{}}}{\displaystyle 𝑑𝐫R(𝐫,t)n_c(𝐫,𝐭)}`$ (63)
$`=`$ $`{\displaystyle 𝑑ϵI_{12}(ϵ,t)}.`$ (64)
Combining this with Eq. (61) leads to
$`{\displaystyle \frac{(\stackrel{~}{N}+N_c)}{t}}`$ $`=`$ $`0,`$ (65)
which demonstrates that the total number of particles is indeed conserved.
### B Energy Conservation
We now consider the conservation of the total energy of the system. The total energy is given by
$`E_{\mathrm{tot}}`$ $`=`$ $`{\displaystyle \frac{d𝐫d𝐩}{(2\pi \mathrm{})^3}\left\{\frac{𝐩^2}{2m}+U_{\mathrm{ext}}(𝐫)+g\left[\stackrel{~}{n}(𝐫,t)+2n_c(𝐫,t)\right]\right\}f(𝐫,𝐩,t)}`$ (67)
$`+{\displaystyle 𝑑𝐫\mathrm{\Psi }^{}(𝐫,t)\left[\frac{\mathrm{}^2\mathbf{}^2}{2m}+U_{\mathrm{ext}}(𝐫)+\frac{g}{2}n_c(𝐫,t)\right]\mathrm{\Psi }(𝐫,t)}.`$
The first term is the semi-classical expression for the total energy of the noncondensate. It contains the kinetic and external potential energy, and the Hartree-Fock mean-field interaction energy of the noncondensed cloud interacting with itself and with the condensate. We note that the self-interaction term is reduced by a factor of two relative to the condensate term to avoid double counting this contribution.
The second term in Eq. (67) is the total energy of the condensate which contains the wave function $`\mathrm{\Psi }(𝐫,t)`$ with normalization
$$𝑑𝐫|\mathrm{\Psi }(𝐫,t)|^2=N_c(t).$$
(68)
It consists of the kinetic energy, the potential energy, and the mean-field energy due to the interaction of the condensate with itself. The mean-field interaction of the condensate with the noncondensate has already been included in the expression for the energy of the noncondensed cloud. We now show that this total energy is indeed conserved during the growth process.
Taking the time derivative of Eq. (67) leads to the following expression,
$`{\displaystyle \frac{E_{\mathrm{tot}}}{t}}`$ $`=`$ $`{\displaystyle \frac{d𝐫d𝐩}{(2\pi \mathrm{})^3}\left\{\frac{𝐩^2}{2m}+U_{\mathrm{ext}}(𝐫)+2g\left[\stackrel{~}{n}(𝐫,t)+n_c(𝐫,t)\right]\right\}\frac{f(𝐫,𝐩,t)}{t}}`$ (71)
$`+{\displaystyle 𝑑𝐫\frac{\mathrm{\Psi }^{}(𝐫,t)}{t}\left\{\frac{\mathrm{}^2\mathbf{}^2}{2m}+U_{\mathrm{ext}}(𝐫)+g\left[n_c(𝐫,t)+2\stackrel{~}{n}(𝐫,t)\right]\right\}\mathrm{\Psi }(𝐫,t)}`$
$`+{\displaystyle 𝑑𝐫\mathrm{\Psi }^{}(𝐫,t)\left\{\frac{\mathrm{}^2\mathbf{}^2}{2m}+U_{\mathrm{ext}}(𝐫)+g\left[n_c(𝐫,t)+2\stackrel{~}{n}(𝐫,t)\right]\right\}\frac{\mathrm{\Psi }(𝐫,t)}{t}}`$
The first term in Eq. (71) can be rewritten as
$`{\displaystyle \frac{d𝐫d𝐩}{(2\pi \mathrm{})^3}E(𝐫,𝐩,t)\frac{f(𝐫,𝐩,t)}{t}}`$ $`=`$ $`{\displaystyle \frac{d𝐫d𝐩}{(2\pi \mathrm{})^3}E(𝐫,𝐩,t)\left(C_{12}[f]+C_{22}[f]\right)}`$ (72)
$`=`$ $`{\displaystyle 𝑑𝐫E_c(𝐫,t)\frac{d𝐩}{(2\pi \mathrm{})^3}C_{12}[f]},`$ (73)
where, to obtain this result, we have used the kinetic equation, Eq. (1), and the fact that the $`C_{22}`$ collision integral conserves energy.
If we again assume that the condensate grows adiabatically as atoms are fed into it from the noncondensate, the condensate wavefunction $`\mathrm{\Psi }(𝐫,t)`$ is a solution of the instantaneous Gross-Pitaevskii equation
$`\left\{{\displaystyle \frac{\mathrm{}^2\mathbf{}^2}{2m}}+U_{\mathrm{ext}}+g\left[n_c(𝐫,t)+2\stackrel{~}{n}(𝐫,t)\right]\right\}\mathrm{\Psi }(𝐫,t)=E_c(t)\mathrm{\Psi }(𝐫,t),`$ (74)
with a time-dependent energy eigenvalue $`E_c(t)`$. For this spatially independent condensate energy, Eq. (72) reduces to
$$\frac{d𝐫d𝐩}{(2\pi \mathrm{})^3}E(𝐫,𝐩,t)\frac{f(𝐫,𝐩,t)}{t}=E_c(t)\frac{\stackrel{~}{N}(t)}{t}.$$
(75)
Inserting this result and Eq. (72) into Eq. (71), the latter is easily seen to yield
$`{\displaystyle \frac{E_{\mathrm{tot}}}{t}}`$ $`=`$ $`E_c(t)\left({\displaystyle \frac{\stackrel{~}{N}}{t}}+{\displaystyle \frac{N_c}{t}}\right)`$ (76)
$`=`$ $`0,`$ (77)
due to the conservation of total particle number. Thus the assumption of adiabaticity is sufficient to ensure that the total energy is conserved. However, one can also show the conservation of energy exactly, without assuming adiabaticity, by making use of the dissipative nonlinear Schrödinger equation in Eq. (10).
## V THOMAS-FERMI APPROXIMATION
The assumption that the condensate grows adiabatically implies that the dynamics of the condensate itself is being neglected, apart from its trivial time dependent normalization. In particular, we are ignoring the possible excitation of internal collective oscillations. However, at the temperatures of interest in the growth process, these excitations are strongly damped and we expect the condensate to remain in a relatively quiescent state which is well approximated by the quasi-equilibrium solution of the GP equation. Indeed, in the experiments there is no evidence of condensate oscillations, although the thermal cloud has been observed to oscillate at twice the harmonic oscillator frequency of the trap in some cases.
For a large number of condensate atoms, a good approximation to the equilibrium wave function is provided by the Thomas-Fermi approximation which neglects the kinetic energy in the GP equation. In this situation, the condensate density is given by
$$n_c(𝐫,t)=\frac{1}{g}\left[\mu _c(t)U_{\mathrm{ext}}(𝐫)2g\stackrel{~}{n}(𝐫,t)\right].$$
(78)
Of course, this expression is only valid if the right-hand side is larger than zero; otherwise, $`n_c(𝐫,t)=0`$. The last term on the right-hand side reflects the mean-field interaction of the condensate with the thermal cloud. Since the latter has a small density relative to the condensate, its effect on the spatial distribution of the condensate is small ($`2g\stackrel{~}{n}\mu _c`$) and we therefore neglect it when determining the condensate density. By the same token, we shall neglect the mean-field interaction of the noncondensate with itself. Strictly speaking, these approximations lead to a violation of total energy conservation, but the error will be very small since the bulk of the mean-field energy, which resides within the condensate itself, is still taken into account. In principle, these contributions can be included in our treatment as shown explicitly in the Appendix. However, because these corrections are small, we have decided to neglect them in our numerical calculations.
It should be noted that the Thomas-Fermi approximation is to some extent dictated by our semi-classical treatment of the noncondensate atoms, since it avoids a potential problem associated with the placement of the condensate chemical potential $`\mu _c`$ relative to the minimum energy available to the thermal atoms, i.e., $`U_{\mathrm{min}}=\mathrm{min}[U_{\mathrm{ext}}+2g(\stackrel{~}{n}+n_c)]`$. For small condensate densities, it is possible that the GP eigenvalue $`\mu _c`$ lies above this minimum value which is clearly impossible if a full quantum treatment of the excited states is retained. In the Thomas-Fermi approximation there is no such problem since the chemical potential is exactly equal to $`U_{\mathrm{min}}`$.
Given these approximations, the time-dependent condensate density profile becomes
$`n_c(𝐫,t)={\displaystyle \frac{1}{g}}\left[\mu _c(t)U_{\mathrm{ext}}(𝐫)\right],`$ (79)
where the external potential is taken to be a general anisotropic harmonic confining potential, $`U_{\mathrm{ext}}(𝐫)=_im\omega _i^2r_i^2/2`$. This expression for the density is again only meaningful when $`U_{\mathrm{ext}}(𝐫)\mu _c(t)`$. The chemical potential of the condensate is given by
$`\mu _c(t)={\displaystyle \frac{\mathrm{}\overline{\omega }}{2}}\left[15N_0(t){\displaystyle \frac{a}{l}}\right]^{2/5},`$ (80)
where $`\overline{\omega }=(\omega _1\omega _2\omega _3)^{1/3}`$ and $`l=\sqrt{\mathrm{}/m\overline{\omega }}`$. The potential experienced by the thermal atoms is then
$$U(𝐫,t)=\{\begin{array}{cc}2\mu _c(t)U_{\mathrm{ext}}(𝐫),\hfill & \text{if }n_c0\hfill \\ U_{\mathrm{ext}}(𝐫),\hfill & \text{if }n_c=0\text{ .}\hfill \end{array}$$
(81)
The minimum value of this potential is $`\mu _c(t)`$ and occurs on the boundary of the condensate.
Three additional important quantities can also be calculated analytically. The first two are the density of states, and the weighted density of states, i.e., $`\rho (\overline{ϵ})`$ and $`\rho _\mathrm{w}(\overline{ϵ})`$ respectively. For the former we get
$`\rho (\overline{ϵ})`$ $`=`$ $`{\displaystyle \frac{m^{3/2}}{\sqrt{2}\pi ^2\mathrm{}^3}}{\displaystyle _{\overline{U}<\overline{ϵ}}}𝑑𝐫\sqrt{\overline{ϵ}\overline{U}(𝐫,t)}`$ (82)
$`=`$ $`{\displaystyle \frac{2}{\pi \mathrm{}\overline{\omega }}}[{\displaystyle _{\overline{U}<\overline{ϵ}}}dyy^2\theta ({\displaystyle \frac{2\mu _c}{\mathrm{}\overline{\omega }}}y^2)\sqrt{{\displaystyle \frac{2(\overline{ϵ}\mu _c)}{\mathrm{}\overline{\omega }}}+y^2}`$ (84)
$`+{\displaystyle _{\overline{U}<\overline{ϵ}}}dyy^2\theta (y^2{\displaystyle \frac{2\mu _c}{\mathrm{}\overline{\omega }}})\sqrt{{\displaystyle \frac{2(\overline{ϵ}+\mu _c)}{\mathrm{}\overline{\omega }}}y^2}]`$
$``$ $`{\displaystyle \frac{2}{\pi \mathrm{}\overline{\omega }}}\left[I_{}(\overline{ϵ})+I_+(\overline{ϵ})\right]`$ (85)
The integrals $`I_{}(\overline{ϵ})`$ and $`I_+(\overline{ϵ})`$ are standard, and are given by
$`I_{}(\overline{ϵ})`$ $`=`$ $`{\displaystyle \frac{u_{}^3x}{4}}{\displaystyle \frac{a_{}u_{}x}{8}}{\displaystyle \frac{a_{}^2}{8}}\mathrm{log}(x+u_{})|_{x=\sqrt{\mathrm{max}\{0,a_{}\}}}^{x=\sqrt{2\mu _c/\mathrm{}\overline{\omega }}}`$ (86)
$`I_+(\overline{ϵ})`$ $`=`$ $`{\displaystyle \frac{u_+^3x}{4}}+{\displaystyle \frac{a_+u_+x}{8}}+{\displaystyle \frac{a_+^2}{8}}\mathrm{arcsin}\left({\displaystyle \frac{x}{\sqrt{a_+}}}\right)|_{x=\sqrt{2\mu _c/\mathrm{}\overline{\omega }}}^{x=\sqrt{a_+}}`$ (87)
where we have defined $`a_\pm =2(\overline{ϵ}\pm \mu _c)/\mathrm{}\overline{\omega }`$, and $`u_\pm =\sqrt{a_\pm x^2}`$.
To obtain an analytic expression for the weighted density of states $`\rho _\mathrm{w}(\overline{ϵ})`$ we note that
$`{\displaystyle \frac{\overline{U}(𝐫,t)}{t}}=2{\displaystyle \frac{\mu _c(t)}{t}}\theta [\mu _c(t)U_{\mathrm{ext}}(𝐫)]{\displaystyle \frac{\mu _c(t)}{t}}.`$ (88)
We therefore find that the weighted density of states is given by
$`\rho _\mathrm{w}(\overline{ϵ})`$ $`=`$ $`{\displaystyle \frac{\mu _c}{t}}\left[\rho (\overline{ϵ})+{\displaystyle \frac{4}{\pi \mathrm{}\overline{\omega }}}{\displaystyle _{\overline{U}<\overline{ϵ}}}𝑑yy^2\theta \left({\displaystyle \frac{2\mu _c}{\mathrm{}\overline{\omega }}}y^2\right)\sqrt{{\displaystyle \frac{2(\overline{ϵ}\mu _c)}{\mathrm{}\overline{\omega }}}+y^2}\right]`$ (89)
$`=`$ $`{\displaystyle \frac{\mu _c}{t}}\left[{\displaystyle \frac{4}{\pi \mathrm{}\overline{\omega }}}I_{}(\overline{ϵ})\rho (\overline{ϵ})\right]`$ (90)
The third important quantity that can be calculated analytically in the Thomas-Fermi approximation arises in the ergodic projection of $`C_{12}`$. With the variable change in Eq.(58), and noting that $`U_{\mathrm{min}}(t)=\mu _c(t)`$ in the Thomas-Fermi approximation, Eq. (57) can be written as
$`I_{12}(\overline{ϵ}_1)`$ $`=`$ $`{\displaystyle \frac{m^3g^2}{2\pi ^3\mathrm{}^7}}{\displaystyle 𝑑\overline{ϵ}_2𝑑\overline{ϵ}_3𝑑\overline{ϵ}_4\left[\delta (\overline{ϵ}_1\overline{ϵ}_2)\delta (\overline{ϵ}_1\overline{ϵ}_3)\delta (\overline{ϵ}_1\overline{ϵ}_4)\right]}`$ (93)
$`\times \left[(1+g_2)g_3g_4g_2(1+g_3)(1+g_4)\right]`$
$`\times {\displaystyle _{\overline{U}\overline{ϵ}_{\mathrm{min}}}}d𝐫n_c(𝐫,t)S(p_2,p_3,p_4)\delta (\overline{ϵ}_2\overline{ϵ}_3\overline{ϵ}_4)`$
where $`p_i=\sqrt{2m(\overline{ϵ}_i\overline{U})}`$. It is apparent that the integrand is symmetric in the variables $`\overline{ϵ}_3`$ and $`\overline{ϵ}_4`$. We can therefore assume, without loss of generality that $`\overline{ϵ}_2\overline{ϵ}_3\overline{ϵ}_4`$ which also implies $`p_2p_3p_4`$. In this situation, $`S(p_2,p_3,p_4)`$ in Eq.(53) reduces to
$`S(p_2,p_3,p_4)={\displaystyle \frac{1}{2}}\left[1\mathrm{sgn}(p_2p_3p_4)\right],`$ (94)
which is nonzero and equal to $`1`$, if
$`p_2<p_3+p_4.`$ (95)
This restricts the spatial integration in Eq. (93) to the domain specified by this inequality. Inserting the definitions of $`p_i`$ and using the conservation of energy condition, $`\overline{ϵ}_2=\overline{ϵ}_3+\overline{ϵ}_4`$, Eq. (95) is equivalent to
$`F(\overline{U})`$ $``$ $`\overline{U}^2{\displaystyle \frac{4}{3}}(\overline{ϵ}_3+\overline{ϵ}_4)\overline{U}+{\displaystyle \frac{4}{3}}\overline{ϵ}_3\overline{ϵ}_4>0.`$ (96)
The roots of $`F(\overline{U})=0`$ are given by
$`\overline{U}_\pm ={\displaystyle \frac{2}{3}}\left[(\overline{ϵ}_3+\overline{ϵ}_4)\pm \sqrt{\overline{ϵ}_3^2\overline{ϵ}_3\overline{ϵ}_4+\overline{ϵ}_4^2}\right],`$ (97)
in terms of which $`F(\overline{U})=(\overline{U}\overline{U}_{})(\overline{U}+\overline{U}_+)`$. The requirement $`F(\overline{U})>0`$ is therefore satisfied for $`\overline{U}<\overline{U}_{}`$ and $`\overline{U}>\overline{U}_+`$. The latter condition, however, is inconsistent with the constraint $`\overline{U}<\overline{ϵ}_4`$ for the integral in Eq. (93). Because $`\overline{U}_{}`$ does satisfy $`\overline{U}_{}\overline{ϵ}_4`$, the net effect of the factor $`S(p_2,p_3,p_4)`$ in Eq. (93) is to restrict the spatial integration domain to the domain defined by $`\overline{U}\overline{U}_{}`$, i.e.,
$`{\displaystyle _{\overline{U}\overline{ϵ}_{\mathrm{min}}}}𝑑𝐫n_c(𝐫,t)S(p_2,p_3,p_4)\delta (\overline{ϵ}_2\overline{ϵ}_3\overline{ϵ}_4)={\displaystyle _{\overline{U}\overline{U}_{}}}𝑑𝐫n_c(𝐫,t)\delta (\overline{ϵ}_2\overline{ϵ}_3\overline{ϵ}_4).`$ (98)
The remaining spatial integral in Eq. (98) can be carried out analytically for the Thomas-Fermi density profile. If $`\overline{U}_{}\mu _c`$, we have simply
$`{\displaystyle _{\overline{U}\overline{U}_{}}}𝑑𝐫n_c(𝐫,t)`$ $`=`$ $`N_c(t).`$ (99)
On the other hand, for $`0\overline{U}_{}\mu _c`$, we have
$`{\displaystyle _{\overline{U}\overline{U}_{}}}𝑑𝐫n_c(𝐫,t)`$ $`=`$ $`N_c(t)\left\{{\displaystyle \frac{5}{2}}\left[1\left(1{\displaystyle \frac{\overline{U}_{}}{\mu _c}}\right)^{3/2}\right]{\displaystyle \frac{3}{2}}\left[1\left(1{\displaystyle \frac{\overline{U}_{}}{\mu _c}}\right)^{5/2}\right]\right\}.`$ (100)
Physically, Eqs. (99) and (100) are a consequence of the kinematical constraints for scattering into the condensate that appear in the original form of the collision integral in Eq. (57).
Finally, we indicate some implications of the assumption of adiabatic growth in the context of the Thomas-Fermi approximation. We take as an approximate solution to the dissipative nonlinear Schrödinger equation in Eq. (10) a condensate wavefunction of the form
$`\mathrm{\Psi }(𝐫,t)=\sqrt{n_c(𝐫,t)}e^{i\theta (𝐫,t)},`$ (101)
where $`n_c(𝐫,t)`$ is the Thomas-Fermi density profile in Eq. (79). Inserting this wavefunction into Eq. (10), neglecting the kinetic energy term as in the Thomas-Fermi approximation, and separating real and imaginary parts of the resulting equation, we obtain the relations
$`\theta (𝐫,t)`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{}}}{\displaystyle _0^t}𝑑t^{}\mu _c(t^{}),`$ (102)
and
$`R(𝐫,t)`$ $`=`$ $`{\displaystyle \frac{\mathrm{}\dot{\mu }_c(t)}{2gn_c(𝐫,t)}}.`$ (103)
The fact that the phase is spatially independent implies that the superfluid velocity $`𝐯_c`$ is zero as we have assumed in Sec. III. According to Eq. (63), Eq. (103) implies
$$\frac{N_c}{t}=\frac{1}{g}𝑑𝐫\dot{\mu }_c(t),$$
(104)
where the integral is restricted to the region occupied by the condensate, i.e., $`n_c(𝐫,t)0`$. This is the same expression obtained by taking the time derivative of the integral of Eq. (79) over all space. We therefore see that the wave function in Eq. (101) is an internally consistent solution of the dissipative nonlinear Schrödinger equation.
## VI RESULTS
In this Section we present the results of our calculations, which were performed for the situation corresponding to the MIT experiments . These used <sup>23</sup>Na atoms confined in an axially symmetric trap with harmonic frequencies of $`18.0`$ Hz and $`82.3`$ Hz along and perpendicular to the symmetry axis, respectively. These values give an averaged frequency of $`\overline{\omega }/2\pi =49.6`$ Hz, which implies that $`\mathrm{}\overline{\omega }/k_B`$ is equal to $`2.4`$ nK. The $`s`$-wave scattering length $`a`$ is $`2.75`$ nm.
To begin, we provide a few of the numerical details. We used a discretized energy mesh consisting of equally spaced points in the range $`0\overline{ϵ}\overline{ϵ}_{\mathrm{max}}`$. The value of the temperature used in the simulations is typically of the order of $`1\mu `$K, which requires a maximum energy range of about $`\overline{ϵ}_{\mathrm{max}}25003000\mathrm{}\overline{\omega }`$ in order to ensure that $`\rho (\overline{ϵ})g(\overline{ϵ})`$ is sufficiently small at the end of the range. In evaluating the collision integrals $`I_{22}`$ in Eq. (26) and $`I_{12}`$ in Eq. (30), the delta functions were used to perform some of the integrations analytically. The remaining integrals were then evaluated numerically using a simple trapezoidal integration scheme. The main advantage of this scheme in the case of the $`I_{22}`$ collision integral is that the conservation of both particle number and energy is numerically exact, which in general is not the case for higher order integration schemes such as Simpson’s rule. This conserving property is especially important in simulations of the condensate growth since a loss of either particles or energy due to numerical inaccuracy would lead to systematic errors in the final equilibrium values for various physical quantities. The situation for the $`I_{12}`$ collision integral is somewhat different since neither of the integrals in Eq. (35) or (36) is zero. Thus the numerical results will depend on the choice of the energy mesh size, and one must check that the results obtained for a given simulation are insensitive to variations in this parameter. The checks we have performed indicate that errors in the final results coming from this numerical source are no larger than a few percent. Errors of this magnitude will not influence the general conclusions that we make. As a final point, we used the Euler method to propagate Eq. (32) in time. The time step was chosen to be sufficiently small, typically 0.5 ms, to ensure the accuracy of the time evolution.
To start the simulation we begin with an initial nonequilibrium distribution which is meant to represent the conditions immediately after the rapid evaporative cooling quench used in the experiments. Ideally, such a quench starts with an equilibrium distribution at some temperature $`T`$ above $`T_c`$ and excises all particles with energy above $`E_{\mathrm{cut}}k_BT_{\mathrm{cut}}`$. We model this by a truncated Bose distribution at the temperature $`T`$. Although we expect this initial distribution to represent the experimental situation reasonably well, there will no doubt be differences from the actual distributions due to the finite time taken to perform the quench, which allows some equilibration to occur, and the possible incomplete removal of all particles in the energy range of the sweep. Since the rf field is resonant only at certain positions in the trap, atoms of a given energy must have sufficient time to reach these positions in order to suffer a spin flip and thus be ejected from the trap. If this is not the case, the distribution in energy will also have a spatial dependence. Some indication that such a nonergodic state in fact occurs is provided by the observation that the thermal cloud starts to oscillate after the quench. However, for lack of detailed information about the experimental initial conditions, we shall assume an idealized truncated Bose distribution as our initial condition.
To complete the specification of the initial state we must also make a choice for the number of atoms initially in the condensate. Of course, if this number is zero, $`I_{12}`$ as given by Eq. (27) is zero since we have only included stimulated transitions into the condensate. In the absence of spontaneous processes there is no possibility of condensate growth. Under the experimental conditions of interest, however, the lowest quantum state initially already has a rather large thermal occupation and stimulated processes will dominate. We therefore choose the initial condensate number to be given by the occupation of the lowest harmonic oscillator state at the temperature of the truncated Bose distribution. This number is typically of the order of a few hundred particles. As our numerical results presented below will show, the growth curves are rather insensitive to this starting value as long as it is small compared to the final equilibrium number of condensate atoms.
In Fig. 1 we show a sequence of growth curves which illustrate the dependence on the parameter $`T_{\mathrm{cut}}`$. In this set of simulations we assume that the temperature of the equilibrium Bose distribution is equal to $`T_c=0.765\mu K`$ and its chemical potential $`\stackrel{~}{\mu }`$ is equal to zero. Before the cut, the gas contains $`\stackrel{~}{N}=40\times 10^6`$ thermal atoms and the number of condensate atoms is given by $`N_c=[\mathrm{exp}(3\beta \mathrm{}\overline{\omega }/2)1]^1=214`$. In a particular simulation, the total number of atoms and the average energy per atom of course depends on the depth of the energy cut. The growth curves are characterized by an initial stage of slow growth during which the truncated Bose distribution evolves into a quasi-equilibrium distribution, a well-defined onset time $`t_{\mathrm{onset}}`$ where a significant increase in the rate of growth occurs, and finally a relaxational stage where the condensate number approaches a final equilibrium value. As the cut is made deeper and deeper, this final number at first increases due to the decreasing total energy of the initial distribution, which results in a lower final temperature. However, at some point the final number of condensate atoms reaches a maximum and then decreases with further deepening of the cut due to the reduced total number of atoms in the initial distribution. To distinguish this behavior the growth curves are shown as solid lines when the final number is increasing with decreasing $`T_{\mathrm{cut}}`$, and conversely, by dashed lines when the final number is decreasing.
Although all the growth curves in Fig. 1 are qualitatively similar, it is clear that there are important differences in detail. For the curves with an increasing equilibrium number of condensate particles, i.e., the solid curves, both the onset time and subsequent relaxation time are seen to decrease with decreasing $`T_{\mathrm{cut}}`$. However, for the curves with an decreasing equilibrium number of condensate particles, i.e., the dashed curves, the dependence of both of these times on further decreases in $`T_{\mathrm{cut}}`$ is much weaker, and they appear to approach limiting values. In order to quantify this behavior, it is convenient to fit the relaxational part of the theoretical growth curves to a simple exponential relaxation
$$N_c^{\mathrm{fit}}(t)N_c^{\mathrm{eq}}\left(1e^{\gamma (tt_{\mathrm{onset}})}\right),$$
(105)
where $`N_c^{\mathrm{eq}}`$, $`\gamma `$ and $`t_{\mathrm{onset}}`$ are fitting parameters. This functional form is found to provide a very good fit to this part of the theoretical curves. Fig. 2 summarizes the results for the onset time, $`t_{\mathrm{onset}}`$, and exponential relaxation rate, $`\gamma `$, for the particular simulations presented in Fig. 1. The onset time decreases from about 100 ms to 20 ms as $`T_{\mathrm{cut}}/T_c`$ is reduced from 5 to 0.5. At the same time, the relaxation rate increases from about 6 s<sup>-1</sup> to 12 s<sup>-1</sup>.
We have also looked at the dependence of the growth curves on the other parameters that appear in the theory. In Fig. 3 we show the growth curves for a range of initial temperatures. Prior to the quench, these initial temperatures are larger than $`T_c`$, and in each case the chemical potential is adjusted to provide again a total of $`40\times 10^6`$ atoms in the thermal cloud. The energy cut and initial number of condensate atoms were taken to be $`T_{\mathrm{cut}}/T_c=2.5`$ and $`N_c(0)=214`$, respectively, and were the same for all the runs. Not surprisingly, we find that the final equilibrium condensate number decreases with increasing initial temperature as a result of the larger average energy per atom. This of course also leads to a higher final equilibrium temperature. However what is somewhat unexpected is the very rapid increase of the onset time as the initial temperature is increased. In Fig. 4(a) we show that a $`30\%`$ variation in $`T/T_c`$ gives rise to more than a ten-fold variation in $`t_{\mathrm{onset}}`$, and that these values are typically much larger than those found using an initial temperature of $`T=T_c`$. In addition, fig. 4(b) shows that the relaxation rate tends to decrease with increasing $`T/T_c`$ and is comparable to the values given in Fig 2.
In Fig. 5 we show the variation of the growth curves with the initial number of condensate atoms. In this case, the initial nonequilibrium distribution is held fixed, corresponding to a Bose distribution with $`\stackrel{~}{N}=40\times 10^6`$, $`T=T_c`$ and $`T_{\mathrm{cut}}/T_c=2.5`$. The growth curve is rather insensitive to the condensate number in the range $`10^2<N_c<10^4`$, but then shows a much stronger dependence in the range $`10^4<N_c<10^6`$. At the higher end of this range, the initial number is already visible on the graph and by $`N_c=10^6`$ there is no longer a meaningful onset time. This would correspond to a situation in which a significant condensate fraction has already formed by the time the quench is completed. This kind of behavior is indeed also seen experimentally under certain conditions.
In order to explain some of these results it is necessary to examine the time evolution of the distribution function $`g(\overline{ϵ},t)`$. In Fig. 6 we show $`\mathrm{ln}(g)`$ vs. $`\overline{ϵ}`$ for various times after the quench. At early times the distribution function is equilibrated by the scattering of thermal atoms into states above the $`E_{\mathrm{cut}}`$ which are initially depleted. To conserve energy, the mean energy of the atoms below $`E_{\mathrm{cut}}`$ must decrease. In fact, the population of the low energy states increases significantly before the onset of rapid condensate growth. This is shown in Fig. 7 where $`g(\overline{ϵ},t)`$ is plotted as a function of time for some specific energy values. We see that $`g(\overline{ϵ},t)`$ at first increases rapidly, reaches a maximum at a time very close to the onset time and then relaxes towards its final equilibrium value of $`(e^{\beta _{\mathrm{eq}}\overline{ϵ}}1)^1`$. This behaviour is typical of all situations in which the growth of a condensate is observed. This strong correlation of the peak position in Fig. 7 with the onset time suggests that condensate formation is triggered by an enhanced low-energy population. Before the onset time, we find numerically that $`g(\overline{ϵ})`$ behaves approximately as $`(\overline{ϵ})^{1.63}`$, which is a stronger singularity than that exhibited by an equilibrium Bose distribution with zero chemical potential, and agrees within our numerical accuracy with the $`(\overline{ϵ})^{5/3}`$ dependence predicted by Svistunov . Regardless of the precise exponent, it seems that a ‘super-critical’ behavior of the distribution function is a precursor to condensate formation .
A useful way to characterize the time evolution of $`g(\overline{ϵ},t)`$ is to express it locally as a Bose distribution
$$g(\overline{ϵ},t)=\frac{1}{\mathrm{exp}(\beta \overline{ϵ}\stackrel{~}{\mu })1},$$
(106)
where the two parameters $`\beta `$ and $`\stackrel{~}{\mu }`$ are defined by fitting this expression to the value of the distribution function and its energy derivative. Although the parameters are treated locally as constants in this procedure, they nevertheless depend parametrically on the energy variable $`\overline{ϵ}`$. The local temperature and chemical potential parameters defined in this way are shown in Fig. 8(a) and (b) at time intervals of 0.05 s for a situation in which the quenched thermal cloud equilibrates to a final temperature above $`T_c`$. Both parameters are seen to be strongly energy dependent at early times but evolve towards energy-independent values by the end of the simulation. The negative equilibrium value of the chemical potential corresponds to an uncondensed thermal cloud at a temperature of about 1.92 $`\mu `$K.
A situation in which the quench leads to the formation of a condensate is illustrated in Figs. 9(a) and (b). The parameters are plotted at 0.25 s intervals during the relaxational stage of the growth curve beyond the onset time. At low energies, the local temperature lies above the final equilibrium value which reflects the higher temperature of the initial Bose distribution. However at higher energies, the local temperature is lower than the final temperature since the gas in this energy range is effectively colder as a result of the quench. Fig. 9(b) shows the corresponding variation of the chemical potential. As a result of the formation of the condensate, the chemical potential at low energies is pinned to zero and then increases at higher energies. The deviations of both the local temperature and chemical potential from their final equilibrium values are seen to relax to zero on a time scale which is comparable with the relaxational stage of the condensate growth. This relaxation rate can therefore be attributed to the relatively slow equilibration of the local temperature and chemical potential of the thermal cloud.
We finally turn to a comparison with experiment. This is shown in Fig. 10 for the particular case in which the starting number of noncondensed atoms is $`40\times 10^6`$, as in the simulations discussed above, but with the initial number of condensate atoms set to $`N_c(0)=10^4`$. In the particular experimental run starting with this total number of atoms before the rf quench, the condensate number is found to relax to a final number of $`1.2\times 10^6`$ atoms. According to Fig. 1, there are two values of $`T_{\mathrm{cut}}`$ which will lead to this final number of condensate atoms, $`T_{\mathrm{cut}}/T_c=0.6`$ and $`T_{\mathrm{cut}}/T_c=5.7`$. The results for the deeper cut of $`T_{\mathrm{cut}}/T_c=0.6`$ are shown as curve (b) and are seen to be in very good agreement with the experimental results. However, we cannot claim good agreement overall since the total number of atoms after the quench is only $`2.5\times 10^6`$ as compared to the experimental number of about $`16.0\times 10^6`$ atoms. For the shallower cut of $`T_{\mathrm{cut}}/T_c=5.7`$ shown as curve (c), the agreement between the theoretical and experimental growth curves is clearly worse in that the theoretical growth rate is too small. Moreover, the total number of atoms remaining in the trap is $`37.5\times 10^6`$ which is too large by roughly a factor of 2. Alternatively, one can choose a cut which reproduces the final number of atoms in the trap. In our simulations, this requires a cut of $`T_{\mathrm{cut}}/T_c=1.9`$. Although the initial growth rate agrees with experiment in this case, the final equilibrium number of condensate atoms is $`4.5\times 10^6`$, which is too large by almost a factor of 4. This final number could be improved by elevating the starting temperature (recall that these simulations used $`T=T_c=0.765`$ $`\mu `$K), however as Fig. 3 shows, achieving a four-fold reduction in the equilibrium condensate number would increase the onset time well beyond the experimental value. It therefore appears that the present simulations cannot reproduce all aspects of the experiments simultaneously.
Fig. 10 also shows a theoretical growth curve for the same initial conditions as for curve (b), but with mean-field interactions between the condensate and thermal cloud turned off. To elaborate, the potential acting on the thermal cloud is simply the time-independent trapping potential, and the condensate is taken to have essentially a delta-function spatial distribution at zero energy. In this case, the integral of $`n_c(𝐫,𝐭)`$ in Eq. (57) is replaced by $`N_c(t)`$. It can be seen that the qualitative behavior is very similar to the fully interacting simulation, but that the equilibrium number of condensate atoms is increased considerably, as expected.
Fig. 11 provides a comparison with another set of experimental results. In this case the initial number of atoms before the quench is not known and was therefore taken to be $`60\times 10^6`$ in order to optimize agreement with experiment. Furthermore, the energy cut was chosen as $`T_{\mathrm{cut}}/T_c=2.5`$. This leads to a final number of $`7.3\times 10^6`$ condensate atoms in the trap, which is approximately the same number as found in the experiment, $`N_c=7.2\times 10^6`$. Although this simulation achieves good agreement between theory and experiment for the condensate growth curve, there are too many unknown variables, including the final number of atoms in the trap, to know whether or not theory is reproducing experiment. For this reason, the results in Fig. 11 should simply be viewed as a possible fit to the experimental data.
## VII DISCUSSION AND OUTLOOK
Our main objective has been to obtain a realistic description of condensate growth which takes into account the effects of mean-field interactions. Within the ergodic approximation for the noncondensed atoms, and the adiabatic approximation for the condensate, the kinetic equation we obtain is given by Eq. (27), and we have used this equation to perform simulations of condensate growth. In agreement with earlier work, we find that the growth curves have a well-defined onset time, after which an exponential relaxation towards equilibrium takes place. Detailed comparison with the results reported in Ref. shows that certain parameters can be tuned in order to achieve agreement with the experimental growth curves. However it seems impossible with the present simulations to reproduce the overall equilibrium state of the trapped gas.
If we attribute the existing discrepancies to theory we must at some point reexamine the two major assumptions made in this work, namely the adiabatic growth of the condensate and the ergodic evolution of the thermal cloud. The adiabatic assumption neglects the dynamics of the condensate, specifically the possibility that collective oscillations are excited during the growth process. Whether or not this has any important effect on the rate at which atoms are exchanged between the condensate and thermal cloud in not known and should be investigated. In the same vein, oscillations of the thermal cloud seen in the experiments clearly indicate the non-ergodic state of the gas which in principle might be important in determining the time scale of equilibration. However, to answer this question requires a solution of the full quantum Boltzmann equation which seems out of reach at the moment. One cannot of course discount the possibility that there are uncertainties in the experimental results themselves. Further experimental work is needed to confirm the earlier results and to explore in more detail the dependences on various parameters such as the initial temperature of the cloud and the depth of the rf cut.
After completion of this work, a preprint by Davis, Gardiner, and Ballagh appeared which is a continuation of a series of papers by Gardiner et al.. It also addresses the issue of mean-field interactions as affecting the density of states, and improves on the authors’ earlier work by giving a more realistic description of the rf quench used in the experiments. Thus, although there are differences in methodology, the physical basis of their work and the approximations they make are essentially equivalent to ours. As confirmation of this equivalence, their calculations of condensate growth performed for the initial conditions of Figs. 10 and 11 yield results which are in quantitative agreement with ours. The situation considered in Fig. 10 is optimal from a theoretical point of view since the experimental conditions are best known in this case. Yet both sets of calculations are unable to reproduce the experimental results in every detail.
One of the differences between their work and ours concerns the way that the condensate is treated. In our formulation, the condensate is isolated explicitly as the macroscopically occupied quantum state, and the remaining excited states making up the thermal cloud are treated semiclassically. As a result of this formulation, we have two kinds of collision integrals, one for thermal atoms scattering amongst each other and a second for collisions of thermal atoms with the condensate. In the formulation of Davis, Gardiner, and Ballagh on the other hand, all states including the condensate are treated equivalently and thus only a single collision integral enters. As a result, the effective collision cross-section involving the condensate does not depend on time as it does in our formulation. A second apparent difference has to do with the term involving the weighted density of states $`\rho _\mathrm{w}`$ in Eq. (27). This term arises as a consequence of the time dependence of the mean-field interaction. Although Davis, Gardiner, and Ballagh also deal with a time-dependent density of states, the second term on the right hand side of Eq. (22) does not appear explicitly in their kinetic equation. However, they account for this term by dividing phase space into energy bins having widths which are a function of time. A final difference involves the use of the Bogoliubov excitation spectrum in the calculation of their density of states, instead of the Hartree-Fock dispersion used here. We do not expect this to affect the condensate growth curves significantly. However, if quasi-particle excitations are invoked, one should in principle also use these states to calculate the collision integrals . It is not known at present what effect this might have on the collision rates for the low-lying energy levels.
Finally, we note that the ergodic treatment of the Boltzmann equation is a powerful, albeit approximate, method which would allow the study of nonequilibrium processes in other situations as well. Some future applications might include the nonequilibrium dynamics of fermion-fermion and boson-fermion mixtures. Thus far, the problem of evaporative cooling in these systems has been studied using a simplified procedure whereby the distribution function is assumed to be given by a cut-off equilibrium distribution function . A cooling trajectory in phase space is then generated by solving for the temperature, chemical potential and cut-off energy at each successive time step. The accuracy of this approach could be checked by solving for the entire distribution function following the methods used here. Another interesting application would be to study a nonequilibrium steady state situation in which atoms are continuously fed into the trapping potential while simultaneously being removed by a rf-cut . This would be relevant to the study of steady-state atom lasers.
###### Acknowledgements.
E. Z. acknowledges the FOM for its financial support during a sabbatical visit to the University of Utrecht, as well as support from the Natural Sciences and Engineering Research Council of Canada. We would like to acknowledge useful discussions with A. Griffin and J. Williams, and with C. W. Gardiner regarding the work in Ref. .
## A Evaluation of the weighted density of states.
In this Appendix, we summarize the steps needed to include in our calculations the effect of the mean-field interactions arising from the noncondensed cloud itself. Referring to Eq. (23), we see that we must evaluate $`U(𝐫,t)/t`$. This quantity is given by
$$\frac{U(𝐫,t)}{t}=2g\left(\frac{\stackrel{~}{n}(𝐫,t)}{t}+\frac{n_c(𝐫,t)}{t}\right).$$
(A1)
The time derivative of $`\stackrel{~}{n}(𝐫,t)`$ can be expressed as
$`{\displaystyle \frac{\stackrel{~}{n}(𝐫,t)}{t}}`$ $`=`$ $`{\displaystyle \frac{}{t}}{\displaystyle \frac{d^{\mathrm{\hspace{0.17em}3}}p}{(2\pi \mathrm{})^3}𝑑ϵ\delta (ϵE(𝐫,𝐩,t))g(ϵ,t)}`$ (A2)
$`=`$ $`{\displaystyle 𝑑ϵ\rho (𝐫,ϵ,t)\left[\frac{U(𝐫,t)}{t}\frac{g(ϵ,t)}{ϵ}+\frac{g(ϵ,t)}{t}\right]}`$ (A3)
$`=`$ $`I(𝐫,t){\displaystyle \frac{U(𝐫,t)}{t}}+{\displaystyle 𝑑ϵ\rho (𝐫,ϵ,t)\frac{g(ϵ,t)}{t}},`$ (A4)
where we have defined
$$I(𝐫,t)𝑑ϵ\rho (𝐫,ϵ,t)\frac{g(ϵ,t)}{ϵ}.$$
(A5)
Substituting Eq. (A1) into Eq. (A2), the latter can be rearranged to provide an expression for the time rate of change of $`\stackrel{~}{n}(𝐫,t)`$ in terms of the time rate of change of the condensate density $`n_c(𝐫,t)`$ and the distribution function $`g(ϵ)`$. We find
$$\frac{\stackrel{~}{n}(𝐫,t)}{t}=\frac{2gI(𝐫,t)}{12gI(𝐫,t)}\frac{n_c(𝐫,t)}{t}+𝑑ϵ\frac{\rho (𝐫,ϵ,t)}{12gI(𝐫,t)}\frac{g(ϵ,t)}{t}.$$
(A6)
Inserting this result into Eq. (A1), we have
$$\frac{U(𝐫,t)}{t}=\frac{2g}{12gI(𝐫,t)}\frac{n_c(𝐫,t)}{N_c}\frac{\stackrel{~}{N}}{t}+\frac{2g}{12gI(𝐫,t)}𝑑ϵ\rho (𝐫,ϵ,t)\frac{g(ϵ,t)}{t}.$$
(A7)
We have here made use of the fact that $`n_c(𝐫,t)`$ depends on time parametrically through $`N_c(t)`$, so that
$$\frac{n_c(𝐫,t)}{t}=\frac{n_c(𝐫,t)}{N_c}\frac{N_c}{t}=\frac{n_c(𝐫,t)}{N_c}\frac{\stackrel{~}{N}}{t}.$$
(A8)
Thus, the weigthed density of states becomes
$`\rho _\mathrm{w}(ϵ,t)`$ $`=`$ $`{\displaystyle d^{\mathrm{\hspace{0.17em}3}}r\rho (𝐫,ϵ,t)\frac{U(𝐫,t)}{t}}`$ (A9)
$`=`$ $`{\displaystyle d^{\mathrm{\hspace{0.17em}3}}r\rho (𝐫,ϵ,t)\left(\frac{2g}{12gI(𝐫,t)}\frac{n_c(𝐫,t)}{N_c}\right)\frac{\stackrel{~}{N}}{t}}`$ (A11)
$`+{\displaystyle d^{\mathrm{\hspace{0.17em}3}}r\rho (𝐫,ϵ,t)𝑑ϵ^{}\frac{2g\rho (𝐫,ϵ^{},t)}{12gI(𝐫,t)}\frac{g(ϵ^{},t)}{t}}`$
$``$ $`A(ϵ,t){\displaystyle \frac{\stackrel{~}{N}}{t}}+{\displaystyle 𝑑ϵ^{}B(ϵ,ϵ^{},t)\frac{g(ϵ^{},t)}{t}},`$ (A12)
where
$$A(ϵ,t)d^{\mathrm{\hspace{0.17em}3}}r\rho (𝐫,ϵ,t)\left(\frac{2g}{12gI(𝐫,t)}\frac{n_c(𝐫,t)}{N_c}\right),$$
(A13)
and
$$B(ϵ,ϵ^{},t)2gd^{\mathrm{\hspace{0.17em}3}}r\frac{\rho (𝐫,ϵ,t)\rho (𝐫,ϵ^{},t)}{12gI(𝐫,t)}.$$
(A14)
We recover the expression for $`\rho _\mathrm{w}(ϵ,t)`$ given in Eq. (89) by setting the kernel $`B`$ equal to zero and neglecting $`I`$ in the expression for $`A`$. It can be seen that including the mean-field of the noncondensate complicates the calculations considerably, but all quantities can in principle be calculated explicitly if these refinements are desired. However, as discussed in Sec. V, we do not expect these effects to be quantitatively important. |
warning/0001/cond-mat0001119.html | ar5iv | text | # Different approaches in the theory of the metastable phase decay on several types of heterogeneous centers
## 1 Iteration method
We shall present all formulas for the case of two sorts of heterogeneous centers. They will be marked by two subscripts $`A`$ and $`B`$. The generalization for the arbitrary number of heterogeneous centers types is evident. It is one of the real advantages of the presented theory.
The system of the condensation equations according to can be presented as following
$$G_A(z)=F_A_0^z(zx)^3\mathrm{exp}(\mathrm{\Gamma }_A(\zeta (x)\mathrm{\Phi }))\theta _A(x)𝑑x$$
$$G_B(z)=F_B_0^z(zx)^3\mathrm{exp}(\mathrm{\Gamma }_B(\zeta (x)\mathrm{\Phi }))\theta _B(x)𝑑x$$
$$\mathrm{\Phi }=\zeta (z)+G_A(z)+G_B(z)$$
$$\theta _A(z)=\mathrm{exp}(K_A_0^z\mathrm{exp}(\mathrm{\Gamma }_A(\zeta (x)\mathrm{\Phi }))𝑑x)$$
$$\theta _B(z)=\mathrm{exp}(K_B_0^z\mathrm{exp}(\mathrm{\Gamma }_B(\zeta (x)\mathrm{\Phi }))𝑑x)$$
with six parameters $`F_A`$, $`F_B`$, $`K_A`$, $`K_B`$, $`\mathrm{\Gamma }_A`$, $`\mathrm{\Gamma }_B`$. Note that the definition of $`\mathrm{\Gamma }_A,\mathrm{\Gamma }_B`$ slightly differs from the standard one - it is $`\mathrm{\Phi }`$ times less.
The presented system can be directly seen from the condensation equations system for one type of heterogeneous centers and there is no need to discuss it more. The problem is how to solve it.
According to the scale invariance one can put $`F_A=1`$, $`\mathrm{\Gamma }_A=1`$. According to the symmetry of the types one can choose them as to have $`F_B<1`$. These simplifications will be interesting only in numerical modeling and here we shall keep all parameters.
We shall define the iteration procedure by the following relations
$$G_{Ai+1}(z)=F_A_0^z(zx)^3\mathrm{exp}(\mathrm{\Gamma }_A(\zeta _i(x)\mathrm{\Phi }))\theta _{Ai}(x)𝑑x$$
$$G_{Bi+1}(z)=F_B_0^z(zx)^3\mathrm{exp}(\mathrm{\Gamma }_B(\zeta _i(x)\mathrm{\Phi }))\theta _{Bi}(x)𝑑x$$
$$\mathrm{\Phi }=\zeta _i(z)+G_{Ai}(z)+G_{Bi}(z)$$
$$\theta _{Ai+1}(z)=\mathrm{exp}(K_A_0^z\mathrm{exp}(\mathrm{\Gamma }_A(\zeta _i(x)\mathrm{\Phi }))𝑑x)$$
$$\theta _{Bi+1}(z)=\mathrm{exp}(K_B_0^z\mathrm{exp}(\mathrm{\Gamma }_B(\zeta _i(x)\mathrm{\Phi }))𝑑x)$$
The number of the iteration approximation is marked by the ordinary subscript instead of the capital letter which marks below the sort of the centers. The values without the number of iteration will mark the real solutions.
The initial approximations are the following
$$G_{A0}=0,G_{B0}=0,\theta _{A0}=1,\theta _{B0}=1$$
Then one can note the following important chains
$$G_{A0}<G_{A2}<\mathrm{}<G_{A2i}<\mathrm{}<G_A<\mathrm{}<G_{A2i+1}<\mathrm{}<G_{A3}<G_{A1}$$
$$G_{B0}<G_{B2}<\mathrm{}<G_{B2i}<\mathrm{}<G_B<\mathrm{}<G_{B2i+1}<\mathrm{}<G_{B3}<G_{B1}$$
$$\theta _{A0}>\theta _{A2}>\mathrm{}>\theta _{A2i}>\mathrm{}>\theta _A>\mathrm{}>\theta _{A2i+1}>\mathrm{}>\theta _{A3}>\theta _{A1}$$
$$\theta _{B0}>\theta _{B2}>\mathrm{}>\theta _{B2i}>\mathrm{}>\theta _B>\mathrm{}>\theta _{B2i+1}>\mathrm{}>\theta _{B3}>\theta _{B1}$$
$$\zeta _0>\zeta _2>\mathrm{}>\zeta _{2i}>\mathrm{}>\zeta >\mathrm{}>\zeta _{2i+1}>\mathrm{}>\zeta _3>\zeta _1$$
for all values of arguments.
These chains prove the convergence of iterations and allow to estimate the accuracy at every step of the calculation.
The problem is to calculate the iterations. In the situation with one type of heterogeneous centers the system of condensation equations can be gotten if we chancel $`G_B`$ and the index $`A`$. The same thing we have to do with the iteration procedure and one can get the mentioned properties. In only one iteration for $`G`$ and two first iterations for $`\theta `$ were calculated. The further iterations can not be calculated analytically. In the situation of nucleation in the system with one type of heterogeneous centers it is sufficient. But in the situation with several types of heterogeneous centers it isn‘t so. To show this we has to recall the physical sense of the first iterations.
The first iteration for $`G`$ is calculated on the base of unexhausted number of heterogeneous centers. So, the behavior of $`\zeta `$ is wrong when the power of exhaustion is essential. But when the centers are exhausted the number of droplets is known - it is equal to the initial number of centers. This primitive notation is the reason why the second iteration for $`\theta `$ gives already suitable result.
In the situation with the several types of heterogeneous centers the situation is different. Really, it is quite possible to have the exhaustion of the first type of heterogeneous centers and the moderate exhaustion of the second type centers. Then the exhaustion of the first type centers has to be taken into account in calculation of the supersaturation which is necessary to get the number of the droplet appeared on the second type centers. But as far as the power of exhaustion of the second type centers is moderate one can not say that the number of the droplets on the second type centers is equal to the total number of the second type centers. So, we have to know the behavior of supersaturation in all cases and it is impossible to do already in the first iterations. Then we have to calculate further iterations.
Unfortunately it is impossible to calculate the further iteration approximations without any simplifications. That’s why two ways of possible approximations are presented in the next sections.
## 2 Avalanche consumption of the metastable phase
The sequential calculation of iterations gives
$$G_{A1}=F_A\frac{z^4}{4}$$
$$\theta _{A1}=\mathrm{exp}(K_Az)$$
$$G_{B1}=F_B\frac{z^4}{4}$$
$$\theta _{B1}=\mathrm{exp}(K_Bz)$$
$$G_{A2}=F_A_0^z(zx)^3\mathrm{exp}(\mathrm{\Gamma }_A(F_A+F_B)x^4/4)\mathrm{exp}(K_Ax)𝑑x$$
$$G_{B2}=F_B_0^z(zx)^3\mathrm{exp}(\mathrm{\Gamma }_B(F_A+F_B)x^4/4)\mathrm{exp}(K_Bx)𝑑x$$
The problem is how to calculate the integrals for $`G_{A2},G_{B2}`$. To calculate these integrals we shall extract functions
$$\phi _A=\mathrm{exp}(\mathrm{\Gamma }_A(G_A+G_B))$$
$$\phi _B=\mathrm{exp}(\mathrm{\Gamma }_B(G_A+G_B))$$
These functions can be calculated in the iteration approximation as
$$\phi _{Ai}=\mathrm{exp}(\mathrm{\Gamma }_A(G_{Ai}+G_{Bi}))$$
$$\phi _{Bi}=\mathrm{exp}(\mathrm{\Gamma }_B(G_{Ai}+G_{Bi}))$$
One can see that $`\phi _{A1},\phi _{B1}`$ can be well approximated as the step functions. Really, $`\mathrm{exp}(x^4)`$ can be approximately interpreted as the step function<sup>1</sup><sup>1</sup>1This was the reason why in experiments one can observe characteristic time lag and this characteristic was introduced in the theory of the metastable phase decay. Then the first theoretical descriptions adopted this fact as the given one without any justification..
Then
$$G_{A2}=F_A_0^z(zx)^3\mathrm{exp}(K_Ax)\mathrm{\Theta }(z_{1A}x)𝑑x$$
$$G_{B2}=F_B_0^z(zx)^3\mathrm{exp}(K_Bx)\mathrm{\Theta }(z_{1B}x)𝑑x$$
where
$$z_{1A}(\frac{4}{\mathrm{\Gamma }_A(F_A+F_B)})^{1/4}$$
$$z_{1B}(\frac{4}{\mathrm{\Gamma }_B(F_A+F_B)})^{1/4}$$
and these integrals can be easily calculated<sup>2</sup><sup>2</sup>2More correctly one gas to multiply $`z_{1A}`$ and $`z_{1B}`$ on $`\mathrm{ln}^{1/4}2`$. This corresponds to the definition of the halfwidth at the halfheight..
For $`z>z_{1A}`$ we have
$$G_{A2}=\underset{i=0}{\overset{3}{}}p_{Ai}z^i$$
where
$$p_{Ai}=\frac{3!(1)^i}{i!(3i)!}F_A_0^{z_{1A}}x^{3i}\mathrm{exp}(K_Ax)𝑑x$$
are some constants which can be calculated analytically.
For $`z>z_{1B}`$ we have
$$G_{B2}=\underset{i=0}{\overset{3}{}}p_{Bi}z^i$$
where
$$p_{Bi}=\frac{3!(1)^i}{i!(3i)!}F_B_0^{z_{1B}}x^{3i}\mathrm{exp}(K_Bx)𝑑x$$
are some constants calculated analytically.
For $`z<z_{1A}`$ we have
$`G_{A2}=F_A\mathrm{exp}(K_Az){\displaystyle \frac{1}{K_A^4}}[(zK_A)^3\mathrm{exp}(zK_A)`$
$`3(zK_A)^2\mathrm{exp}(zK_A)+6zK_A\mathrm{exp}(zK_A)6(\mathrm{exp}(K_Az)1)]`$
For $`z<z_{1B}`$ we have
$`G_{B2}=F_B\mathrm{exp}(K_Bz){\displaystyle \frac{1}{K_B^4}}[(zK_B)^3\mathrm{exp}(zK_B)`$
$`3(zK_B)^2\mathrm{exp}(zK_B)+6zK_B\mathrm{exp}(zK_B)6(\mathrm{exp}(K_Bz)1)]`$
The calculation of $`\theta _{A2}`$, $`\theta _{A2}`$ can be done in the same way
$$\theta _{A2}=\mathrm{exp}(K_A_0^z\mathrm{exp}(\mathrm{\Gamma }_A(F_A+F_B)\frac{x^4}{4})𝑑x)$$
$$\theta _{B2}=\mathrm{exp}(K_B_0^z\mathrm{exp}(\mathrm{\Gamma }_B(F_A+F_B)\frac{x^4}{4})𝑑x)$$
Then
$$\theta _{A2}=\mathrm{exp}(K_A[\mathrm{\Theta }(z_{1A}z)z+\mathrm{\Theta }(zz_{1A})z_{1A}])$$
$$\theta _{B2}=\mathrm{exp}(K_B[\mathrm{\Theta }(z_{1B}z)z+\mathrm{\Theta }(zz_{1B})z_{1B}])$$
Now we have to calculate $`\theta _{A3},\theta _{B3}`$. One can note that functions $`\phi _{A2},\phi _{B2}`$ can be also presented as the step functions. Namely, all functions
$`\mathrm{exp}(\mathrm{\Gamma }_A(F_A\mathrm{exp}(K_Az){\displaystyle \frac{1}{K_A^4}}[(zK_A)^3\mathrm{exp}(zK_A)`$
$`3(zK_A)^2\mathrm{exp}(zK_A)+6zK_A\mathrm{exp}(zK_A)6(\mathrm{exp}(K_Az)1)]+`$
$`F_B\mathrm{exp}(K_Bz){\displaystyle \frac{1}{K_B^4}}[(zK_B)^3\mathrm{exp}(zK_B)3(zK_B)^2\mathrm{exp}(zK_B)+`$
$`6zK_B\mathrm{exp}(zK_B)6(\mathrm{exp}(K_Bz)1)]))`$
$`\mathrm{exp}(\mathrm{\Gamma }_A(F_A\mathrm{exp}(K_Az){\displaystyle \frac{1}{K_A^4}}[(zK_A)^3\mathrm{exp}(zK_A)3(zK_A)^2\mathrm{exp}(zK_A)+`$
$`6zK_A\mathrm{exp}(zK_A)6(\mathrm{exp}(K_Az)1)]+{\displaystyle \underset{i=0}{\overset{3}{}}}p_{Bi}z^i))`$
$$\mathrm{exp}(\mathrm{\Gamma }_A(\underset{i=0}{\overset{3}{}}p_{Ai}z^i+\underset{i=0}{\overset{3}{}}p_{Bi}z^i))$$
$`\mathrm{exp}(\mathrm{\Gamma }_A({\displaystyle \underset{i=0}{\overset{3}{}}}p_{Ai}z^i+F_B\mathrm{exp}(K_Bz){\displaystyle \frac{1}{K_B^4}}[(zK_B)^3\mathrm{exp}(zK_B)`$
$`3(zK_B)^2\mathrm{exp}(zK_B)+6zK_B\mathrm{exp}(zK_B)6(\mathrm{exp}(K_Bz)1)]))`$
$`\mathrm{exp}(\mathrm{\Gamma }_B(F_A\mathrm{exp}(K_Az){\displaystyle \frac{1}{K_A^4}}[(zK_A)^3\mathrm{exp}(zK_A)`$
$`3(zK_A)^2\mathrm{exp}(zK_A)+6zK_A\mathrm{exp}(zK_A)6(\mathrm{exp}(K_Az)1)]+`$
$`F_B\mathrm{exp}(K_Bz){\displaystyle \frac{1}{K_B^4}}[(zK_B)^3\mathrm{exp}(zK_B)3(zK_B)^2\mathrm{exp}(zK_B)+`$
$`6zK_B\mathrm{exp}(zK_B)6(\mathrm{exp}(K_Bz)1)]))`$
$`\mathrm{exp}(\mathrm{\Gamma }_B(F_A\mathrm{exp}(K_Az){\displaystyle \frac{1}{K_A^4}}[(zK_A)^3\mathrm{exp}(zK_A)3(zK_A)^2\mathrm{exp}(zK_A)+`$
$`6zK_A\mathrm{exp}(zK_A)6(\mathrm{exp}(K_Az)1)]+{\displaystyle \underset{i=0}{\overset{3}{}}}p_{Bi}z^i))`$
$$\mathrm{exp}(\mathrm{\Gamma }_B(\underset{i=0}{\overset{3}{}}p_{Ai}z^i+\underset{i=0}{\overset{3}{}}p_{Bi}z^i))$$
$`\mathrm{exp}(\mathrm{\Gamma }_B({\displaystyle \underset{i=0}{\overset{3}{}}}p_{Ai}z^i+F_B\mathrm{exp}(K_Bz){\displaystyle \frac{1}{K_B^4}}[(zK_B)^3\mathrm{exp}(zK_B)`$
$`3(zK_B)^2\mathrm{exp}(zK_B)+6zK_B\mathrm{exp}(zK_B)6(\mathrm{exp}(K_Bz)1)]))`$
have the step-like behavior. This is the central point of our calculations. This property can be directly seen by calculations.
The possibility to see the step-like behavior of exponent of the supersaturation deviation directly on the base of explicit expressions is the evident advantage of the iteration method. Otherwise we have to prove this property in the general situation taking into account that exponent of the supersaturation deviation lies between $`\mathrm{exp}(z^4)`$ and $`\mathrm{exp}(z^3)`$ after the suitable renormalization.
Now to calculate $`\theta _{A3}`$, $`\theta _{B3}`$ we have to define the values $`z_{2A}`$, $`z_{2B}`$ by equations
$$\mathrm{\Gamma }_A(G_{A2}(z_{2A})+G_{B2}(z_{2A}))=1$$
$$\mathrm{\Gamma }_B(G_{A2}(z_{2B})+G_{B2}(z_{2B}))=1$$
Then
$$\theta _{A3}=\mathrm{exp}(K_A[\mathrm{\Theta }(z_{2A}z)z+\mathrm{\Theta }(zz_{2A})z_{2A}])$$
$$\theta _{B3}=\mathrm{exp}(K_B[\mathrm{\Theta }(z_{2B}z)z+\mathrm{\Theta }(zz_{2B})z_{2B}])$$
The calculation of $`G_{3A}`$ and $`G_{3B}`$ can be done by the following way
$`G_{A3}=F_A{\displaystyle _0^z}(zx)^3\mathrm{exp}(\mathrm{\Gamma }_A(G_{A2}+G_{B2}))`$
$`\mathrm{exp}(K_A[\mathrm{\Theta }(z_{1A}x)x+\mathrm{\Theta }(xz_{1A})z_{1A}])dx`$
$`G_{A3}=F_A{\displaystyle _0^z}(zx)^3\mathrm{\Theta }(z_{2A}x)`$
$`\mathrm{exp}(K_A[\mathrm{\Theta }(z_{1A}x)x+\mathrm{\Theta }(xz_{1A})z_{1A}])dx`$
$`G_{B3}=F_B{\displaystyle _0^z}(zx)^3\mathrm{exp}(\mathrm{\Gamma }_B(G_{A2}+G_{B2}))`$
$`\mathrm{exp}(K_B[\mathrm{\Theta }(z_{1B}x)x+\mathrm{\Theta }(xz_{1B})z_{1B}])dx`$
$`G_{B3}=F_B{\displaystyle _0^z}(zx)^3\mathrm{\Theta }(z_{2B}x)`$
$`\mathrm{exp}(K_B[\mathrm{\Theta }(z_{1B}x)x+\mathrm{\Theta }(xz_{1B})z_{1B}])dx`$
These integrals can be easily taken in analytical form. For $`G_{A3}`$ one can get
* For $`z>z_{2A}`$ with arbitrary $`z_{1A}`$ we have
$$G_{A3}=\underset{i=0}{\overset{3}{}}z^ip_{Ai}$$
where
$`p_{Ai}={\displaystyle \frac{(1)^i3!}{i!(3i)!}}F_A{\displaystyle _0^{z_{2A}}}x^i[\mathrm{exp}(K_Az_{1A})\mathrm{\Theta }(xz_{1A})+`$
$`\mathrm{exp}(K_Ax)\mathrm{\Theta }(z_{1A}x)]dx`$
can be calculated analytically.
* For $`z<min\{z_{1A},z_{2A}\}`$
$$G_{A3}=F_A_0^z(zx)^3\mathrm{exp}(K_Ax)𝑑x$$
has been already calculated.
* For $`z_{1A}<z<z_{2A}`$
$$G_{A3}=\underset{i=0}{\overset{3}{}}z^ip_{Ai}+F_A\frac{(zz_{1A})^4}{4}\mathrm{exp}(K_Az_{1A})$$
where
$$p_{Ai}=F_A\frac{3!(1)^i}{i!(3i)!}_0^{z_{1A}}x^i\mathrm{exp}(K_Ax)𝑑x$$
All $`p_{Ai}`$ are constants and can be calculated analytically.
For $`G_{B3}`$ one can get
* For $`z>z_{2B}`$ with arbitrary $`z_{1B}`$
$$G_{B3}=\underset{i=0}{\overset{3}{}}z^ip_{Bi}$$
where
$`p_{Bi}={\displaystyle \frac{(1)^i3!}{i!(3i)!}}F_B{\displaystyle _0^{z_{2B}}}x^i[\mathrm{exp}(K_Bz_{1B})\mathrm{\Theta }(xz_{1B})+`$
$`\mathrm{exp}(K_Bx)\mathrm{\Theta }(z_{1B}x)]dx`$
* For $`z<min\{z_{1B},z_{2B}\}`$
$$G_{B3}=F_B_0^z(zx)^3\mathrm{exp}(K_Bx)𝑑x$$
has been already calculated.
* For $`z_{1B}<z<z_{2B}`$
$$G_{B3}=\underset{i=0}{\overset{3}{}}z^ip_{Bi}+F_B\frac{(zz_{1B})^4}{4}\mathrm{exp}(K_Bz_{1B})$$
where
$$p_{Bi}=F_B\frac{3!(1)^i}{i!(3i)!}_0^{z_{1B}}x^i\mathrm{exp}(K_Bx)𝑑x$$
All $`p_{Bi}`$ are constants and can be calculated analytically.
One can easily see that $`z_{1A}<z_{2A},z_{1B}<z_{2B}`$ and the third opportunity can take place.
Moreover, if we define $`z_{iA}`$, $`z_{iB}`$ by equalities
$$\mathrm{\Gamma }_A(G_{Ai}(z_{iA})+G_{Bi}(z_{iA}))=1$$
$$\mathrm{\Gamma }_B(G_{Ai}(z_{iB})+G_{Bi}(z_{iB}))=1$$
then we can get we following chains of inequalities
$$z_{1A}<z_{3A}<\mathrm{}<z_{2i+1A}<\mathrm{}<z_A<\mathrm{}<z_{2iA}<\mathrm{}<z_{4A}<z_{2A}$$
$$z_{1B}<z_{3B}<\mathrm{}<z_{2i+1B}<\mathrm{}<z_B<\mathrm{}<z_{2iB}<\mathrm{}<z_{4B}<z_{2B}$$
where $`z_A`$, $`z_B`$ are defined by
$$\mathrm{\Gamma }_A(G_A(z_A)+G_B(z_A))=1$$
$$\mathrm{\Gamma }_B(G_A(z_B)+G_B(z_B))=1$$
Then we can calculate $`G_{A4}`$, $`G_{B4}`$. In the initial expressions
$$G_{A4}=F_A_0^z(zx)^3\mathrm{exp}(\mathrm{\Gamma }_A(G_{A3}+G_{B3})\theta _{A3}dx$$
$$G_{B4}=F_B_0^z(zx)^3\mathrm{exp}(\mathrm{\Gamma }_A(G_{A3}+G_{B3})\theta _{B3}dx$$
we have to use two characteristic lengths $`z_2`$ and $`z_3`$. Then
$`G_{A4}=F_A{\displaystyle _0^z}\mathrm{\Theta }(z_{3A}x)(zx)^3`$
$`[\mathrm{exp}(K_Ax)\mathrm{\Theta }(z_{2A}x)+\mathrm{exp}(K_Az_{2A})\mathrm{\Theta }(xz_{2A})]dx`$
$`G_{B4}=F_B{\displaystyle _0^z}\mathrm{\Theta }(z_{3B}x)(zx)^3`$
$`[\mathrm{exp}(K_Bx)\mathrm{\Theta }(z_{2B}x)+\mathrm{exp}(K_Bz_{2B})\mathrm{\Theta }(xz_{2B})]dx`$
These integrals can be easily taken in an analytical form.
For $`G_{A4}`$ one can get
* For $`z>z_{3A}`$ with arbitrary $`z_{2A}`$
$$G_{A4}=\underset{i=0}{\overset{3}{}}z^ip_{Ai}$$
where
$`p_{Ai}={\displaystyle \frac{(1)^i3!}{i!(3i)!}}F_A{\displaystyle _0^{z_{3A}}}x^i[\mathrm{exp}(K_Az_{2A})\mathrm{\Theta }(xz_{2A})+`$
$`\mathrm{exp}(K_Ax)\mathrm{\Theta }(z_{2A}x)]dx`$
* For $`z<min\{z_{2A},z_{3A}\}`$
$$G_{A4}=F_A_0^z(zx)^3\mathrm{exp}(K_Ax)𝑑x$$
has been already calculated.
* For $`z_{2A}<z<z_{3A}`$
$$G_{A4}=\underset{i=0}{\overset{3}{}}z^ip_{Ai}+F_A\frac{(zz_{2A})^4}{4}\mathrm{exp}(K_Az_{2A})$$
where
$$p_{Ai}=F_A\frac{3!(1)^i}{i!(3i)!}_0^{z_{2A}}x^i\mathrm{exp}(K_Ax)𝑑x$$
All $`p_{Ai}`$ are constants and can be calculated analytically.
For $`G_{B4}`$ one can get
* For $`z>z_{3B}`$ with arbitrary $`z_{2B}`$
$$G_{B4}=\underset{i=0}{\overset{3}{}}z^ip_{Bi}$$
where
$`p_{Bi}={\displaystyle \frac{(1)^i3!}{i!(3i)!}}F_B{\displaystyle _0^{z_{3B}}}x^i[\mathrm{exp}(K_Bz_{2B})\mathrm{\Theta }(xz_{2B})+`$
$`\mathrm{exp}(K_Bx)\mathrm{\Theta }(z_{2B}x)]dx`$
* For $`z<min\{z_{2B},z_{3B}\}`$
$$G_{B4}=F_B_0^z(zx)^3\mathrm{exp}(K_Bx)𝑑x$$
has been already calculated.
* For $`z_{2B}<z<z_{3B}`$
$$G_{B4}=\underset{i=0}{\overset{3}{}}z^ip_{Bi}+F_B\frac{(zz_{2B})^4}{4}\mathrm{exp}(K_Bz_{2B})$$
where
$$p_{Bi}=F_B\frac{3!(1)^i}{i!(3i)!}_0^{z_{2B}}x^i\mathrm{exp}(K_Bx)𝑑x$$
All $`p_{Bi}`$ are constants and can be calculated analytically.
One can easily see that $`z_{3A}<z_{2A},z_{3B}<z_{2B}`$ and the third opportunity can not take place.
One can calculate by the same way all further iterations<sup>3</sup><sup>3</sup>3That’s why we left all possibilities in the last expressions..
We see that due to the difference of the relations between parameters $`z_{Ai}`$, $`z_{Bi}`$ even the functional form of the iterations will be different. But it is quite easy to see that calculations of further iterations will be done absolutely analogously. Moreover, the functional forms of all odd further iterations will be similar. The functional forms of all even iterations will be also similar.
Then one can extract these functional forms and it is necessary to obtain the values of parameters $`z_{Ai}`$, $`z_{Bi}`$ in these forms.
Moreover, one can put one and the same values of ”final” parameters $`z_A`$, $`z_B`$ in equations for $`G_A`$, $`G_B`$, $`\theta _A`$, $`\theta _B`$. This leads to the following result
In the initial expressions
$$G_A=F_A_0^z(zx)^3\mathrm{exp}(\mathrm{\Gamma }_A(G_A+G_B))\theta _A𝑑x$$
$$G_B=F_B_0^z(zx)^3\mathrm{exp}(\mathrm{\Gamma }_A(G_A+G_B))\theta _B𝑑x$$
Then
$$G_A=F_A_0^z\mathrm{\Theta }(z_Ax)(zx)^3\mathrm{exp}(K_Ax)𝑑x$$
$$G_B=F_B_0^z\mathrm{\Theta }(z_Bx)(zx)^3\mathrm{exp}(K_Bx)𝑑x$$
These integrals can be easily taken in analytical form.
For $`G_A`$ one can get
* For $`z>z_A`$
$$G_A=\underset{i=0}{\overset{3}{}}z^ip_{Ai}$$
where
$$p_{Ai}=\frac{(1)^i3!}{i!(3i)!}F_A_0^{z_A}x^i\mathrm{exp}(K_Ax)𝑑x$$
* For $`z<z_A`$
$$G_{A4}=F_A_0^z(zx)^3\mathrm{exp}(K_Ax)𝑑x$$
has been already calculated.
All $`p_{Ai}`$ are constants and can be calculated analytically.
For $`G_B`$ one can get
* For $`z>z_B`$
$$G_B=\underset{i=0}{\overset{3}{}}z^ip_{Bi}$$
where
$$p_{Bi}=\frac{(1)^i3!}{i!(3i)!}F_B_0^{z_B}x^i\mathrm{exp}(K_Bx)𝑑x$$
* For $`z<z_B`$
$$G_{B4}=F_B_0^z(zx)^3\mathrm{exp}(K_Bx)𝑑x$$
has been already calculated.
All $`p_{Bi}`$ are constants and can be calculated analytically.
As the result now we know the functional form for $`G_A`$, $`G_B`$, $`\theta _A`$, $`\theta _B`$. The condensations equations system now determines the characteristic values $`z_A`$, $`z_B`$.
The step-like behavior of $`\phi `$ occurs due to the avalanche character of the vapor consumption. The avalanche character of the vapor consumption has to be considered as the main feature of the first order phase transition kinetics. This feature was described in in details. Meanwhile it is possible to make the theory more accurate taking into account the regime of the droplets growth.
Note that all functions $`\phi `$ are lying between $`\mathrm{exp}(x^4)`$ and $`\mathrm{exp}(x^3)`$ after the suitable renormalization. Namely, they are more sharp than $`\mathrm{exp}(x^3)`$ and more smooth than $`\mathrm{exp}(x^4)`$. Note that
$$_0^{\mathrm{}}\mathrm{exp}(x^4)𝑑x_0^{\mathrm{}}\mathrm{exp}(x^3)𝑑x0.9D$$
Then one can make results of the theory more accurate if we take
* instead $`\mathrm{\Theta }(z_{iA}x)`$ the function $`\mathrm{\Theta }(Dz_{iA}x)`$
* instead $`\mathrm{\Theta }(z_{iB}x)`$ the function $`\mathrm{\Theta }(Dz_{iB}x)`$
* instead $`\mathrm{\Theta }(xz_{iA})`$ the function $`D\mathrm{\Theta }(xDz_{iA})`$
* instead $`\mathrm{\Theta }(xz_{iB})`$ the function $`D\mathrm{\Theta }(xDz_{iB})`$
and multiply all parameters $`z_{iA}`$, $`z_{iB}`$ on $`D`$.
We can see that the generalization of the presented approach to the multicomponent case is quite evident. But here the number of unknown parameters $`z_A`$ will be equal to the number of the sorts of heterogeneous centers. The problem how to solve this system is rather actual and here we shall propose a way to do it.
We shall choose the sorts of heterogeneous centers as to have
$$\mathrm{\Gamma }_A>\mathrm{\Gamma }_B>\mathrm{\Gamma }_C>\mathrm{\Gamma }_D>\mathrm{}.$$
Then
$$z_A<z_B<z_C<z_D<\mathrm{}\mathrm{}$$
At first we shall determine $`z_A`$. Note that the behavior of the system after $`z_A`$. Then we can put all all $`\mathrm{\Gamma }_B`$, $`\mathrm{\Gamma }_C`$ etc. to be equal to $`\mathrm{\Gamma }_A`$ and $`z_B`$, $`z_C`$ etc. to be equal $`z_A`$. Then we have only one parameter $`z_A`$ and can determine it from one algebraic equation analogous to those described in (the situation is absolutely analogous to the intermediate situation).
Now $`z_A`$ is determined and $`G_A`$, $`\theta _A`$ are known. It allows to go to the determination of $`z_B`$. It can be done by the same procedure. We have to consider $`G_A`$, $`\theta _A`$ as known values and put all $`z_C`$, $`z_D`$ etc. to be equal to $`z_B`$. Then the analogous equation can be solved and it gives $`z_B`$, $`G_B`$, $`\theta _B`$. This procedure can be continued. Note that at first this way was described in for the special monodisperse approximation.
## 3 Special monodisperse approximation
The task to give the theoretical description of the complex systems required new approximations for some characteristic functions. As the result the monodisperse approximation for $`G_A`$ and $`G_B`$ was suggested in . The monodisperse approximation can be used in two variants - the first one is the monodisperse approximation with fixed number of droplets which is suitable when $`G_A`$, $`G_B`$ are really important, the second one is the monodisperse approximation with a floating number of droplets involved in this approximation. A special recipe allows to get approximation which is suitable during all times.
The second variant of approximation is necessary for the systems with the strong hierarchy between the probabilities of droplets formation on different sorts of centers. Nevertheless in solution of the system by the procedure analogous to the already described in the end of the last section it is possible to use only the first type of approximation. This states the significance of the first type of approximation. But as far as we follow here the iteration procedure it will be necessary to use the second variant of approximation.
The problems to calculate iterations appear in the second approximation:
$$G_{A2}=F_A_0^z(zx)^3\mathrm{exp}(\mathrm{\Gamma }_A(F_A+F_B)\frac{x^4}{4})\mathrm{exp}(K_Ax)𝑑x$$
$$G_{B2}=F_B_0^z(zx)^3\mathrm{exp}(\mathrm{\Gamma }_B(F_A+F_B)\frac{x^4}{4})\mathrm{exp}(K_Bx)𝑑x$$
The justification of the monodisperse approximation has been already presented many times (see , ) and we shall use it here without discussion. Then
$$G_{A2}F_Az^3_0^{z/4}\mathrm{exp}(\mathrm{\Gamma }_A(F_A+F_B)\frac{x^4}{4})\mathrm{exp}(K_Ax)𝑑x$$
$$G_{B2}F_Bz^3_0^{z/4}\mathrm{exp}(\mathrm{\Gamma }_B(F_A+F_B)\frac{x^4}{4})\mathrm{exp}(K_Bx)𝑑x$$
The real advantage of consideration of the monodisperse approximation on the level of iterations is the possibility to calculate the error of approximation explicitly. This error is small.
One can make the iteration approximation more accurate by the substitution of $`\mathrm{exp}(K_Ax)`$ by
$$\theta _{2A}=\mathrm{exp}(K_A_0^z\mathrm{exp}(\mathrm{\Gamma }_A(F_A+F_B)\frac{x^4}{4})𝑑x)$$
and by the substitution of $`\mathrm{exp}(K_Bx)`$ by
$$\theta _{2B}=\mathrm{exp}(K_B_0^z\mathrm{exp}(\mathrm{\Gamma }_B(F_A+F_B)\frac{x^4}{4})𝑑x)$$
This leads to
$$G_{2A}(z)=F_A(1\theta _{2A}(\frac{z}{4}))z^3/K_A$$
$$G_{2B}(z)=F_B(1\theta _{2B}(\frac{z}{4}))z^3/K_B$$
When monodisperse approximation is used it convenient to use the following way to construct iterations
$$G_{Ai+1}(z)=F_A_0^z(zx)^3\mathrm{exp}(\mathrm{\Gamma }_A(\zeta _i(x)\mathrm{\Phi }))\theta _{Ai+1}(x)𝑑x$$
$$G_{Bi+1}(z)=F_B_0^z(zx)^3\mathrm{exp}(\mathrm{\Gamma }_B(\zeta _i(x)\mathrm{\Phi }))\theta _{Bi+1}(x)𝑑x$$
All other formulas remain the same. Certainly the chains of inequalities will be slightly violated, but this can not leads to divergence of iterations. Moreover, one can simply add the difference between sequential $`\theta `$ to the error between sequential iterations.
Now we see that in monodisperse approximation
$$G_{iA}(z)=\frac{F_Az^3}{K_A}(1\theta _{Ai}(z/4))$$
$$G_{iB}(z)=\frac{F_Bz^3}{K_B}(1\theta _{Bi}(z/4))$$
So, we need to calculate only $`\theta _{Ai}`$, $`\theta _{Bi}`$.
In this calculation one can use the avalanche character of the vapor consumption again. Then
$$\theta _{2A}(z)=\mathrm{exp}(K_Az_{1A})\theta (zz_{1A})+\mathrm{exp}(K_Az)\mathrm{\Theta }(z_{1A}z)$$
$$\theta _{2B}(z)=\mathrm{exp}(K_Bz_{1B})\theta (zz_{1B})+\mathrm{exp}(K_Bz)\mathrm{\Theta }(z_{1B}z)$$
where $`z_{1A}`$, $`z_{1B}`$ are given by the previous expressions (but further $`z_{iA}`$, $`z_{iB}`$ will have slightly another numerical values).
Note that in practice we don’t use such expression but can use the simple asymptotes (”essential asymptotes”)
$$\theta _{2A}=\mathrm{exp}(K_Az)$$
$$\theta _{2B}=\mathrm{exp}(K_Bz)$$
for all $`z`$.
Really, we need these expressions only for $`zz_A`$, $`zz_B`$. As far as the argument is $`z/4`$ it means that these asymptotes can not be used only when $`z_A4z_B`$ or $`z_B4z_A`$. It can take place only when $`\mathrm{\Gamma }_A`$ and $`\mathrm{\Gamma }_B`$ differ more than in $`4^3=64`$ times. The last seems to be practically unrealizable.
The third iteration for $`\theta `$ gives
$`\theta _{3A}=\mathrm{exp}(K_A{\displaystyle _0^z}\mathrm{exp}(\mathrm{\Gamma }_A(F_A(1\mathrm{exp}(K_A{\displaystyle \frac{x}{4}})`$
$`\mathrm{\Theta }({\displaystyle \frac{x}{4}}+z_{1A})\mathrm{exp}(K_Az_{1A})\mathrm{\Theta }(z_{1A}+{\displaystyle \frac{x}{4}}))+`$
$`F_B(1\mathrm{exp}(K_B{\displaystyle \frac{x}{4}})\mathrm{\Theta }({\displaystyle \frac{x}{4}}+z_{1B})`$
$`\mathrm{exp}(K_Bz_{1B})\mathrm{\Theta }(z_{1B}+{\displaystyle \frac{x}{4}}))))dx)`$
$`\theta _{3B}=\mathrm{exp}(K_B{\displaystyle _0^z}\mathrm{exp}(\mathrm{\Gamma }_B(F_A(1\mathrm{exp}(K_A{\displaystyle \frac{x}{4}})`$
$`\mathrm{\Theta }({\displaystyle \frac{x}{4}}+z_{1A})\mathrm{exp}(K_Az_{1A})\mathrm{\Theta }(z_{1A}+{\displaystyle \frac{x}{4}}))+`$
$`F_B(1\mathrm{exp}(K_B{\displaystyle \frac{x}{4}})\mathrm{\Theta }({\displaystyle \frac{x}{4}}+z_{1B})`$
$`\mathrm{exp}(K_Bz_{1B})\mathrm{\Theta }(z_{1B}+{\displaystyle \frac{x}{4}}))))dx)`$
One can directly see that
$`G_{Ai+1}(z)=F_Az^3{\displaystyle _0^{z/4}}\mathrm{exp}(\mathrm{\Gamma }_A(G_{Ai}(x)+G_{Bi}(x)))`$
$`\mathrm{exp}(K_A{\displaystyle _0^x}\mathrm{exp}(\mathrm{\Gamma }_A(G_{iA}(x^{})+G_{iB}(x^{})))𝑑x^{})dx=`$
$`{\displaystyle \frac{F_Az^3}{K_A}}(1\mathrm{exp}(K_A{\displaystyle _0^z}\mathrm{exp}(\mathrm{\Gamma }_A(G_{iA}(x)+G_{iB}(x)))𝑑x))`$
$`G_{Bi+1}(z)=F_Bz^3{\displaystyle _0^{z/4}}\mathrm{exp}(\mathrm{\Gamma }_B(G_{Ai}(x)+G_{Bi}(x)))`$
$`\mathrm{exp}(K_B{\displaystyle _0^x}\mathrm{exp}(\mathrm{\Gamma }_B(G_{iA}(x^{})+G_{iB}(x^{})))𝑑x^{})dx=`$
$`{\displaystyle \frac{F_Bz^3}{K_B}}(1\mathrm{exp}(K_B{\displaystyle _0^z}\mathrm{exp}(\mathrm{\Gamma }_B(G_{iA}(x)+G_{iB}(x)))𝑑x))`$
Then $`G_{A3}`$, $`G_{B3}`$ have the same structure as $`G_{A2}`$, $`G_{B2}`$ have. Certainly the numerical values of parameters $`z_{Ai}`$, $`z_{Bi}`$ can be slightly another. All other iteration approximations will have the same form.
Now we can reformulate the system of the condensation equations because we know the functional form of $`G_A`$, $`G_B`$. It will be the following
$$G_A(z)=\frac{F_Az^3}{K_A}(1\mathrm{exp}(K_A_0^{z/4}\mathrm{exp}(\mathrm{\Gamma }_A(G_A(x)+G_B(x)))𝑑x))$$
$$G_B(z)=\frac{F_Bz^3}{K_B}(1\mathrm{exp}(K_B_0^{z/4}\mathrm{exp}(\mathrm{\Gamma }_B(G_A(x)+G_B(x)))𝑑x))$$
In the avalanche approximation it will be the following
$$G_A(z)=\frac{F_Az^3}{K_A}(1\mathrm{exp}(K_Az/4)\mathrm{\Theta }(z_Az/4)\mathrm{exp}(K_Az_A\mathrm{\Theta }(z/4z_A))$$
$$G_B(z)=\frac{F_Bz^3}{K_B}(1\mathrm{exp}(K_Bz/4)\mathrm{\Theta }(z_Bz/4)\mathrm{exp}(K_Bz_B\mathrm{\Theta }(z/4z_B))$$
where $`z_A`$, $`z_B`$ are defined by the previous relations.
In the approximation of essential asymptotes
$$G_A(z)=\frac{F_Az^3}{K_A}(1\mathrm{exp}(K_Az/4))$$
$$G_B(z)=\frac{F_Bz^3}{K_B}(1\mathrm{exp}(K_Bz/4))$$
These algebraic equations can be easily solved.
Now one can mention some ways to go away from the avalanche approximation. In the approximation of essential asymptotes everything is clear. In the avalanche approximation one can consider the number of the free heterogeneous centers as the smooth function of $`z`$ and approximately consider
$$G_A(z)=\frac{F_Az^3}{K_A}(1\mathrm{exp}(K_A_0^{z_A/4}\mathrm{exp}(\mathrm{\Gamma }_A(G_A(x)+G_B(x)))𝑑x))$$
$$G_B(z)=\frac{F_Bz^3}{K_A}(1\mathrm{exp}(K_B_0^{z_B/4}\mathrm{exp}(\mathrm{\Gamma }_B(G_A(x)+G_B(x)))𝑑x))$$
Then $`G_Aconstz^3`$, $`G_Bconstz^3`$, we have the monodisperse approximation with fixed number of droplets and need only to calculate the integrals of the type $`_0^z\mathrm{exp}(x^3)𝑑x`$. This leads to the necessity to take
* instead $`\mathrm{\Theta }(z_{iA}x)`$ the function $`\mathrm{\Theta }(Dz_{iA}x)`$
* instead $`\mathrm{\Theta }(z_{iB}x)`$ the function $`\mathrm{\Theta }(Dz_{iB}x)`$
* instead $`\mathrm{\Theta }(xz_{iA})`$ the function $`D\mathrm{\Theta }(xDz_{iA})`$
* instead $`\mathrm{\Theta }(xz_{iB})`$ the function $`D\mathrm{\Theta }(xDz_{iB})`$
and to multiply all parameters $`z_{iA}`$, $`z_{iB}`$ on $`D`$. Here
$$D_0^{\mathrm{}}\mathrm{exp}(x^3)𝑑x=0.9$$
The generalization of the presented approach to the multicomponent case is quite evident. Our actions are absolutely identical to those described for the avalanche model. The only thing to stress is that at every step we can use ”essential asymptote” for all sorts of heterogeneous centers whose $`G_A`$ are still undetermined. Those $`G_A`$ which are already determined have to be taken explicitly. Note that at first this way was described in .
As the result we have to conclude that approximations of the avalanche consumption or the special monodisperse approximation lead to the similarity of the iterations of the high order (at least in their functional form). So, we see now the correspondence between the iteration approach and approximations used here. The numerical calculations show that the errors of these approximations are negligible in frames of the accuracy of the modern experiment. |
warning/0001/hep-th0001064.html | ar5iv | text | # The effect of different regulators in the non-local field-antifield quantizationPresent addreess: Departamento de Física e Química, UNESP, Campus de Guaratinguetá
## I Introduction
The non-local regularization (NLR) gives a consistent way to compute one-loop anomalies of theories with an action that can be decomposed into a kinetic and an interacting part. It can be proved that anomalies at higher order levels of $`\mathrm{}`$ can be precisely obtained with this regularization. The main ideas were based on the Schwinger’s proper time method . The NLR arranges the original divergent loop integrals in a sum over loop contribution in such a way that the loops, now composed of a set of auxiliary fields, contain the original singularities. To regularize the original theory one has to eliminate these auxiliary fields by putting them on shell. In this way the theory is free of the quantum fluctuations. The preliminary results were very well received.
The method developed by Batalin and Vilkovisky (BV method) showed itself to be a very powerfull way to quantize the most difficult gauge field theories. For a review see .
The BV, or field-antifield formalism, provides at the lagrangian level, a general framework for the covariant path integral quantization of gauge theories. This formalism uses very interesting mathematical objects such as a Poisson-like bracket (the antibracket), canonical transformations, ghosts and antighosts for the BRST transformations, etc. The most important object of this method at the classical level is an equation called classical master equation (CME).
The fundamental idea of this formalism is the BRST invariance. All the fields $`\mathrm{\Phi }^A`$, i.e., the set of the classical fields of the theory toghether with the ghosts and the auxiliary fields, have their canonically conjugated fields, the antifields $`\mathrm{\Phi }_A^{}`$. With all these elements we construct the so called BV action. At the classical level, the BV action becomes the classical action when all the antifields are equal to zero.
There is two ways to get a gauge-fixed action: by a canonical transformation, and now we can say that the action is in a gauge-fixed basis; or with a choice of a gauge fermion and making the antifields to be equal to the functional derivative of this fermion.
The method can be applied to gauge theories which have an open algebra (when the algebra of the gauge transformations closes only on shell); to closed algebras; to gauge theories that have structure functions rather than constants (soft algebras); and to the case where the gauge transformations may or may not be independent, i.e., reducible or irreducible algebras respectively.
Zinn-Justin introduced the concept of sources of the BRST-transformations . These sources are the antifields in the BV formalism. It was shown also that the geometry of the antifields have a natural origin .
At the quantum level, the field-antifield formalism also works at higher order loop anomalies . At one loop, with the addition of extra degrees of freedon, causing an extension of the original configuration space, we have a solution for a quantum master equation (QME) that has been obtained as a part that does not depend on the antifields in the anomaly. In general this solution needs a regularization as we will see below. When the Wess-Zumino terms (which cancell the anomaly) can not be found, the theory can be said to have a genuine anomaly. Recently, a method was developed to handle with global anomalies .
However, as has been explained above, the solution of the QME is not easily obtained because there is a $`\delta (0)`$-like divergence when the $`\mathrm{\Delta }`$ operator, a second order differential operator that wil be defined in the next section, is applied on local functionals. The details can be seen in ref. . Therefore, a regularization method has to be used to cut the divergence in the QME. One of these prescriptions is the Pauli-Villars (PV) regularization method , where new fields, the PV fields, and an arbitrary mass matrix are introduced. But this method is very usefull only at one-loop level. At higher orders, the PV method is still misterious. Very recently, a BPHZ renormalization of the BV formalism was formulated . The dimensional regularization method at the quantum aspect of the field-antifield quantization has been studied in ref. . Finally, an extension of the NLR method to the BV framework has been recently formulated by J. París . The consistency conditions for higher orders anomalies have been studied in the reference .
The objective of this paper is to make a comparison between two ways of regularizing the BV formalism using the pure extended NLR. To do this we have analyzed the results of the NLR using two different kinds of regulators.
It is well known that, for some models, the value of the anomaly can depend on the regulator operator that is being used, i. e., we can obtain different but cohomologically equivalent expressions for the anomaly for these models . The chiral Schwinger model (CSM), anomalous at one-loop only, is one of these models which the guise of the expression for the anomaly is dependent of the form of the regulator as has been demonstrated in through the analysis of the Wess-Zumino (WZ) term, which is responsible for the cancellation of the anomaly. It has been used two different regulators and the results have showed two different WZ terms, but the expressions are equivalent after a reparametrization. In this paper we have regularized the CSM within the context of this extended non-local BV regularization calculating the CSM’s anomaly. Firstly we have computed the functional traces using a simple second order differential regulator. After this, the Fujikawa-like regulator was adjusted to this modified BV formalism. We show in a precise way that, using these different regulators we can obtain directly the same result for the anomaly.
This work is organized as follows: in section 2 a brief review of the field-antifield formalism has been made. In section 3 the original NLR was depicted. The extended non-local regularization was described in section 4. The computation of the CSM anomaly at one-loop with the two regulators has been calculated in section 5. In the last section, we have summarized the conclusions and final remarks.
## II The Field-Antifield Formalism
Let us construct the complete set of fields, including in this set the classical fields, the ghosts for all gauge symmetries and the auxiliary fields. The complete set will be denoted by $`\mathrm{\Phi }^A`$. Now, one will extend this space with the same number of fields, but at this time, defining the antifields $`\mathrm{\Phi }_A^{}`$, which are the canonical conjugated variables with respect to the antibracket structure. This last object is constructed like
$$(X,Y)=\frac{\delta _rX}{\delta \varphi }\frac{\delta _lY}{\delta \varphi ^{}}(XY),$$
(1)
where the indices $`r`$ and $`l`$ denote right and left functional derivatives respectively.
By means of the antibrackets, one can write the canonical conjugation relations
$$(\mathrm{\Phi }^A,\mathrm{\Phi }_B^{})=\delta _B^A,(\mathrm{\Phi }^A,\mathrm{\Phi }^B)=(\mathrm{\Phi }_A^{},\mathrm{\Phi }_B^{})=0.$$
(2)
The antifields $`\mathrm{\Phi }_A^{}`$ have opposite statistics to their conjugated fields $`\mathrm{\Phi }^A`$. The antibracket is a fermionic operation so that the statistics of the antibracket $`(X,Y)`$ is opposite to that of the simple product $`XY`$. The antibracket also satifies some graded Jacobi relations:
$$(X,(Y,Z))+()^{ϵ_Xϵ_Y+ϵ_X+ϵ_Y}(Y,(X,Z))=((X,Y),Z).$$
(3)
where $`ϵ_X`$ is the statistics of $`X`$, i.e. $`ϵ(X)=ϵ_X`$.
We define a quantity, named ghost number, to the fields and to the antifields. These are integers such that
$$gh(\mathrm{\Phi }^{})=\mathrm{\hspace{0.17em}1}gh(\mathrm{\Phi }).$$
(4)
One can then construct an extended action of ghost number equal to zero, the so called BV action, also called classical proper solution,
$`S(\mathrm{\Phi },\mathrm{\Phi }^{})`$ $`=`$ $`S_{cl}(\mathrm{\Phi })+\mathrm{\Phi }_A^{}R^A(\mathrm{\Phi })+{\displaystyle \frac{1}{2}}\mathrm{\Phi }_A^{}\mathrm{\Phi }_B^{}R^{BA}(\mathrm{\Phi })+\mathrm{}`$ (5)
$`+`$ $`{\displaystyle \frac{1}{n!}}\mathrm{\Phi }_{A_1}^{}\mathrm{}\mathrm{\Phi }_{A_n}^{}R^{A_n\mathrm{}A_1}+\mathrm{},`$ (6)
so that it has to satisfy the classical master equation,
$$(S,S)=\mathrm{\hspace{0.33em}0}.$$
(7)
This equation contains the complete algebra of the theory, the gauge invariances of the classical action (where $`S_{cl}=S_{BV}(\mathrm{\Phi }^A,\mathrm{\Phi }_A^{}=0)`$), Jacobi identities,$`\mathrm{}`$.
Gauge fixing is obtained either by a canonical transformation or by choosing a fermion $`\mathrm{\Psi }`$ and writing
$$\mathrm{\Phi }_A^{}=\frac{\delta _r\mathrm{\Psi }}{\delta \mathrm{\Phi }^A}.$$
(8)
To obey the ghost number conservation rule in this expression one have to introduce the BRST antighost in the gauge fixing fermion.
At the quantum level the action can be defined by
$$W=S+\underset{p=1}{\overset{\mathrm{}}{}}\mathrm{}^pM_p,$$
(9)
where the $`M_p`$ are the corrections (the Wess-Zumino terms) to the quantum action. The expansion (9) is not the only one, but it is the usual one. An expansion in $`\sqrt{\mathrm{}}`$ can be made, for example . This will originate the so called background charges, that are useful in the conformal field theory .
The quantization of the theory is obtained with the generating functional of the Green functions:
$$Z(J,\mathrm{\Phi }^{})=𝒟\mathrm{\Phi }exp\frac{i}{\mathrm{}}\left[W(\mathrm{\Phi },\mathrm{\Phi }^{})+J^A\mathrm{\Phi }_A^{}\right].$$
(10)
But the definition of a path integral properly lacks on a regularization framework, as we have observed already, which can be seen as a way to define the measure of the integral. Anomalies represent the non conservation of the classical symmetries at the quantum level.
For a theory to be free of anomalies, the quantum action $`W`$ has to be a solution of the QME,
$$(W,W)=2i\mathrm{}\mathrm{\Delta }W,$$
(11)
where
$$\mathrm{\Delta }(1)^{A+1}\frac{_r}{\mathrm{\Phi }^A}\frac{_r}{\mathrm{\Phi }_A^{}}.$$
(12)
In the equation (11) one can see that when it is not possible to find a solution to the QME, we have an anomaly that can be defined by:
$$𝒜\left[\mathrm{\Delta }W+\frac{i}{2\mathrm{}}(W,W)\right](\mathrm{\Phi },\mathrm{\Phi }^{}).$$
(13)
The anomaly can be represented by a $`\mathrm{}`$ expansion,
$$𝒜=\underset{p=1}{\overset{\mathrm{}}{}}\mathrm{}^{p1}M_p.$$
(14)
Substituting (9) in (13) and using (14) one have the form of the $`p`$-loop BRST anomalies:
$`𝒜_0`$ $`=`$ $`{\displaystyle \frac{1}{2}}(S,S)\mathrm{\hspace{0.17em}0},`$ (15)
$`𝒜_1`$ $`=`$ $`\mathrm{\Delta }S+i(M_1,S),`$ (16)
$`𝒜_p`$ $`=`$ $`\mathrm{\Delta }M_{p1}+{\displaystyle \frac{i}{2}}{\displaystyle \underset{q=1}{\overset{p1}{}}}(M_q,M_{pq})`$ (17)
$`+`$ $`i(M_p,S),p2.`$ (18)
The first equation is the known CME. The second one is an equation using $`M_1`$. If, substituting (9) in (11), there is not a solution for $`M_1`$ then $`𝒜`$ is called a genuine anomaly.
The anomaly is not uniquely determined since $`M_1`$ is arbitrary. The anomaly satisfy the Wess-Zumino consistency condition :
$$(𝒜,S)=0.$$
(19)
It was extensively analyzed in ref. that two different regulators furnish consistent anomalies that are related by a local counterterm,
$$i\mathrm{\Delta }^{(2)}S=i\mathrm{\Delta }^{(1)}S+(S,M_1),$$
(20)
where $`M_1`$ is a local counterterm.
We will show that we can obtain directly the same result for the anomaly of the CSM using the NLR method with two different regulators.
## III The Non-local Regularization
As we have stressed in the introduction, the non-local regularization can be applied only to theories which have a perturbative expansion, i.e. for actions that can be decomposed into a free and an interacting part. For much more details, including the diagrammatic part, the interested reader can see the references <sup>*</sup><sup>*</sup>*For convenience we are using the same notation as the reference .. Here we have explained the main parts of the method.
Let us define an action $`S(\mathrm{\Phi })`$ where $`\mathrm{\Phi }`$ is the set $`\mathrm{\Phi }^A`$ of the fields, $`A=1,\mathrm{},N`$, and with statistics $`ϵ(\mathrm{\Phi }^A)ϵ_A`$,
$$S(\mathrm{\Phi })=F(\mathrm{\Phi })+I(\mathrm{\Phi }),$$
(21)
where $`F(\mathrm{\Phi })`$ is the kinetic part and $`I(\mathrm{\Phi })`$ is the interacting part, which is an analitic function in $`\mathrm{\Phi }^A`$ around $`\mathrm{\Phi }^A=0`$.
Then one can write conveniently that
$$F(\mathrm{\Phi })=\frac{1}{2}\mathrm{\Phi }^A_{AB}\mathrm{\Phi }^B,$$
(22)
and $`_{AB}`$ is called the kinetic operator.
To perform the NLR we have now to introduce a cut-off or regulating parameter $`\mathrm{\Lambda }^2`$. An arbitrary and invertible matrix $`T_{AB}`$ has to be introduced too. The combination of $`_{AB}`$ with $`(T^1)^{AB}`$ defines a second order derivative regulator:
$$_B^A=(T^1)^{AC}_{CB}.$$
(23)
We can construct two important operators with these objects. The first is the smearing operator
$$ϵ_B^A=exp\left(\frac{_B^A}{2\mathrm{\Lambda }^2}\right),$$
(24)
and the second is the shadow kinetic operator
$$𝒪_{AB}^1=T_{AC}(\stackrel{~}{𝒪}^1)_B^C=\left(\frac{}{ϵ^21}\right)_{AB},$$
(25)
with $`(\stackrel{~}{𝒪})_B^A`$ defined as
$`\stackrel{~}{𝒪}_B^A`$ $`=`$ $`\left({\displaystyle \frac{ϵ^21}{}}\right)_B^A`$ (26)
$`=`$ $`{\displaystyle _0^1}{\displaystyle \frac{dt}{\mathrm{\Lambda }^2}}exp\left(t{\displaystyle \frac{_B^A}{\mathrm{\Lambda }^2}}\right).`$ (28)
In order to expand our original configuration space for each field $`\mathrm{\Phi }^A`$, an auxiliary field $`\mathrm{\Psi }^A`$ can be constructed. We will call these last fields as the shadow fields, with the same statistics as the auxiliary fields. A new auxiliary action involves both sets of fields
$$\stackrel{~}{𝒮}(\mathrm{\Phi },\mathrm{\Psi })=F(\widehat{\mathrm{\Phi }})A(\mathrm{\Psi })+I(\mathrm{\Phi }+\mathrm{\Psi }).$$
(29)
The second term of this auxiliary action is called the auxiliary kinetic term,
$$A(\mathrm{\Psi })=\frac{1}{2}\mathrm{\Psi }^A(𝒪^1)_{AB}\mathrm{\Psi }^B.$$
(30)
The fields $`\widehat{\mathrm{\Phi }}^A`$, the smeared fields, which make part of the auxiliary action are defined by
$$\widehat{\mathrm{\Phi }}^A(ϵ^1)_B^A\mathrm{\Phi }^B.$$
(31)
It can be proved that, to eliminate the quantum fluctuations associated with the shadow fields at the path integral level, one has to accomplish this by puting the auxiliary fields $`\mathrm{\Psi }`$ on shell. So, the classical shadow field equations of motion are
$$\frac{_r\stackrel{~}{S}(\mathrm{\Phi },\mathrm{\Psi })}{\mathrm{\Psi }}=\mathrm{\hspace{0.17em}0}\mathrm{\Psi }^A=\left(\frac{_rI}{\mathrm{\Phi }^B}(\mathrm{\Phi }+\mathrm{\Psi })\right)𝒪^{BA}.$$
(32)
These equations can be solved in a perturbative fashion. The classical solutions $`\overline{\mathrm{\Psi }}_0(\mathrm{\Phi })`$ can now be substituted in the auxiliary action (29). This substitution modify the auxiliary action so that a new action, the non-localized action appear,
$$𝒮_\mathrm{\Lambda }(\mathrm{\Phi })\stackrel{~}{𝒮}(\mathrm{\Phi },\overline{\mathrm{\Psi }}_0(\mathrm{\Phi })).$$
(33)
The action (33) can be expanded in $`\overline{\mathrm{\Psi }}_0`$. As a result, we see the appearance of the smeared kinetic term $`F(\widehat{\mathrm{\Phi }})`$, the original interaction term $`I(\mathrm{\Phi })`$ and an infinite series of new non-local interaction terms. But all these interaction terms are $`O\left(\mathrm{\Lambda }^2\right)`$-like and when the limit $`\mathrm{\Lambda }^2\mathrm{}`$ is applied, we will have that $`𝒮_\mathrm{\Lambda }(\mathrm{\Phi })𝒮(\mathrm{\Phi })`$, and the original theory is obtained. Equivalently to this limit, the same result can be acquired with the limits
$$ϵ1,𝒪0,\overline{\mathrm{\Psi }}_0(\mathrm{\Phi })0.$$
(34)
With all this framework, when we introduce the smearing operator, any local quantum field theory can be made ultraviolet finite. But a question about symmetry can appear. Obviously this form of non-localization, i.e. (31), in general destroy any kind of gauge symmetry or its associated BRST symmetry. The final consequence is the damage of the corresponding Ward identities at the tree level. However, the invariance of the theory can be preserved introducing the auxiliary fields in the original symmetries .
Let us make an analysis of what happens. If the original action (21) is invariant under the infinitesimal transformation
$$\delta \mathrm{\Phi }^A=R^A(\mathrm{\Phi }),$$
(35)
then it can be proved that the auxiliary action is invariant under the auxiliary infinitesimal transformations
$`\stackrel{~}{\delta }\mathrm{\Phi }^A`$ $`=`$ $`\left(ϵ^2\right)_B^AR^B(\mathrm{\Phi }+\mathrm{\Psi }),`$ (36)
$`\stackrel{~}{\delta }\mathrm{\Psi }^A`$ $`=`$ $`\left(1ϵ^2\right)_B^AR^B(\mathrm{\Phi }+\mathrm{\Psi }).`$ (37)
However, the non-locally regulated action (33) is invariant under the transformation
$$\delta _\mathrm{\Lambda }(\mathrm{\Phi }^A)=\left(ϵ^2\right)_B^AR^B\left(\mathrm{\Phi }+\overline{\mathrm{\Psi }}_0(\mathrm{\Phi })\right),$$
(38)
remembering that $`\overline{\mathrm{\Psi }}_0(\mathrm{\Phi })`$ are the solutions of the classical equations of motions (32).
Hence, any of the original continuous symmetries of the theory are preserved at the tree level, even the BRST transformations, and consequently, the original gauge symmetry. The reader can see for details.
## IV The Extended (BV) Non-local Regularization
As had been said before, the fundamental principle of the field-antifield formalism is the BRST invariance. Therefore, it is simple to realize that the connection of the NLR method with the BV formalism is possible. Using the above construction of the NLR and the BV results, one can build a regulated BRST classical structure of a general gauge theory from the original one. Consequently, a non-locally regularized BV formalism comes out.
We are now in the BV environment. Hence, the configuration space has to be enlarged introducing the antifields $`\{\mathrm{\Psi }^A,\mathrm{\Psi }_A^{}\}`$. Note that the shadow fields have antifields too. Then, an auxiliary proper solution incorporates the auxiliary action (29) (corresponding to the gauge-fixed action $`S(\mathrm{\Phi })`$), its gauge symmetry (36) and the unknown associated higher order structure functions. The auxiliary BRST transformations are modified by the presence of the term $`\mathrm{\Phi }_A^{}R^A(\mathrm{\Phi })`$ in the original proper solution. Then it can be written that the BRST transformations terms are
$$\left[\mathrm{\Phi }_A^{}(ϵ^2)_B^A+\mathrm{\Psi }_A^{}(1ϵ^2)_B^A\right]R^B\left(\mathrm{\Phi }+\mathrm{\Psi }\right),$$
(39)
which are originated from the following substitutions
$`R^A`$ $``$ $`R^A(\mathrm{\Phi }+\mathrm{\Psi })R^A(\mathrm{\Theta }),`$ (40)
$`\mathrm{\Phi }_A^{}`$ $``$ $`\left[\mathrm{\Phi }_A^{}(ϵ^2)_B^A+\mathrm{\Psi }_A^{}(1ϵ^2)_B^A\right]\mathrm{\Theta }_A^{}.`$ (41)
For higher orders, the natural way would be
$$R^{A_n\mathrm{}A_1}(\mathrm{\Phi })R^{A_n\mathrm{}A_1}(\mathrm{\Phi }+\mathrm{\Psi })=R^{A_n\mathrm{}A_1}(\mathrm{\Theta }),$$
(42)
and an obvious ansatz for the auxiliary proper solution is
$`\stackrel{~}{S}(\mathrm{\Phi },\mathrm{\Phi }^{};\mathrm{\Psi },\mathrm{\Psi }^{})`$ $`=`$ $`\stackrel{~}{S}(\mathrm{\Phi },\mathrm{\Psi })+\mathrm{\Theta }_A^{}R^A(\mathrm{\Theta })`$ (43)
$`+`$ $`\mathrm{\Theta }_A^{}\mathrm{\Theta }_B^{}R^{BA}(\mathrm{\Theta })+\mathrm{}`$ (44)
$`+`$ $`\mathrm{\Theta }_{A_1}^{}\mathrm{}\mathrm{\Theta }_{A_n}^{}R^{A_n\mathrm{}A_1}(\mathrm{\Phi })+\mathrm{}.`$ (45)
It is intuitive to see that the same canonical conjugation relations, the equations (2), can be obtained, i.e.
$$(\mathrm{\Theta }^A,\mathrm{\Theta }_B^{})=\delta _B^A.$$
(46)
Consequently, we have to construct a new set of fields and antifields $`\{\mathrm{\Sigma }^A,\mathrm{\Sigma }_A^{}\}`$ defined by
$$\mathrm{\Sigma }^A=\left[\left(1ϵ^2\right)_B^A\mathrm{\Phi }^B\left(ϵ^2\right)_B^A\mathrm{\Psi }^B\right],$$
(47)
and
$$\mathrm{\Sigma }_A^{}=\mathrm{\Phi }_A^{}\mathrm{\Psi }_A^{}.$$
(48)
Now we have that the linear transformation
$$\{\mathrm{\Phi }^A,\mathrm{\Phi }_A^{};\mathrm{\Psi }^A,\mathrm{\Psi }_A^{}\}\{\mathrm{\Theta }^A,\mathrm{\Theta }_A^{};\mathrm{\Sigma }^A,\mathrm{\Sigma }_A^{}\}$$
(49)
is canonical in the antibracket sense. The auxiliary action (29) is the original proper solution (5) with arguments $`\{\mathrm{\Theta }^A,\mathrm{\Theta }_A^{}\}`$.
The elimination of the auxiliary fields in the non-local BV method is the next step. The shadow fields have to be substituted by the solutions of their classical equations of motion. At the same time, their antifields will be equal to zero. In this way we can write
$$S_\mathrm{\Lambda }(\mathrm{\Phi },\mathrm{\Phi }^{})=\stackrel{~}{S}(\mathrm{\Phi },\mathrm{\Phi }^{};\mathrm{\Psi },\mathrm{\Psi }^{}=0),$$
(50)
and the classical equations of motion are
$$\frac{\delta _r\stackrel{~}{S}(\mathrm{\Phi },\mathrm{\Phi }^{};\mathrm{\Psi },\mathrm{\Psi }^{})}{\delta \mathrm{\Psi }^A}=0$$
(51)
with solutions $`\overline{\mathrm{\Psi }}\overline{\mathrm{\Psi }}(\mathrm{\Phi },\mathrm{\Phi }^{})`$, which explicitly read
$`\overline{\mathrm{\Psi }}^A`$ $`=`$ $`[{\displaystyle \frac{\delta _rI}{\delta \mathrm{\Phi }^B}}(\mathrm{\Phi }+\mathrm{\Psi })+\mathrm{\Phi }_C^{}\left(ϵ^2\right)_D^CR_B^D(\mathrm{\Phi }+\mathrm{\Psi })`$ (52)
$`+`$ $`O\left((\mathrm{\Phi }^{})^2\right)]𝒪^{BA}`$ (53)
with
$$R_B^A=\frac{\delta _rR^A(\mathrm{\Phi })}{\delta \mathrm{\Phi }^B}.$$
(54)
The lowest order of equation (52) is,
$$\overline{\mathrm{\Psi }}_0^A=\left(\frac{\delta _rI}{\delta \mathrm{\Phi }^B}(\mathrm{\Phi }+\mathrm{\Psi })\right)𝒪^{BA}$$
(55)
and one can obtain an expression for $`\overline{\mathrm{\Psi }}(\mathrm{\Phi },\mathrm{\Phi }^{})`$ at any desired order in the antifields .
To quantize the theory, it is necessary to add the extra counterterms $`M_p`$ to preserve the quantum counterpart of the classical BRST scheme. It is the same as to substitute the classical action $`S`$ by a quantum action $`W`$. In the original papers the quantization of the theory was already analyzed, but it seems that only the one-loop $`M_1`$ corrections acquired BRST invariance. It can be proved that in the field-antifield framework, in general, two-loops and higher order loop corrections should also be considered .
The complete interaction term, $`(\mathrm{\Phi },\mathrm{\Phi }^{})`$, of the original proper solution can be written as
$$(\mathrm{\Phi },\mathrm{\Phi }^{})I(\mathrm{\Phi })+\mathrm{\Phi }_A^{}R^A(\mathrm{\Phi })+\mathrm{\Phi }_A^{}\mathrm{\Phi }_B^{}R^{BA}(\mathrm{\Phi })+\mathrm{}$$
(56)
The non-localization of this interaction part furnishes a way to regularize interactions from the counterterms $`M_p`$. To construct the auxiliary free and interactions parts we have that
$`\stackrel{~}{F}(\mathrm{\Phi }+\mathrm{\Psi })`$ $`=`$ $`F(\widehat{\mathrm{\Phi }})A(\mathrm{\Psi }),`$ (57)
$`(\mathrm{\Phi },\mathrm{\Phi }^{};\mathrm{\Psi },\mathrm{\Psi }^{})`$ $`=`$ $`(\mathrm{\Theta },\mathrm{\Theta }^{}),`$ (58)
with $`\{\mathrm{\Theta },\mathrm{\Theta }^{}\}`$ already known.
Now one have to put the auxiliary fields on shell and its antifields equal to zero, so that
$`F_\mathrm{\Lambda }(\mathrm{\Phi },\mathrm{\Phi }^{})`$ $`=`$ $`\stackrel{~}{F}(\mathrm{\Phi },\overline{\mathrm{\Psi }}_0),`$ (59)
$`_\mathrm{\Lambda }(\mathrm{\Phi },\mathrm{\Phi }^{})`$ $`=`$ $`\stackrel{~}{}(\mathrm{\Phi }+\overline{\mathrm{\Psi }}_0,\mathrm{\Phi }^{}ϵ^2),`$ (60)
then $`S_\mathrm{\Lambda }=F_\mathrm{\Lambda }+_\mathrm{\Lambda }`$.
The quantum action $`W`$ can be expressed by
$$W=F++\underset{p=1}{\overset{\mathrm{}}{}}\mathrm{}^pM_pF+𝒴$$
(61)
where $`𝒴`$ is the generalized quantum interaction part.
An analogous procedure of the previous section can be applied to the quantum action $`W`$. We will omit all the formal steps here. All the details can be founded in ref. .
A decomposition in its divergent part and its finite part when $`\mathrm{\Lambda }^2\mathrm{}`$ can be accomplished in the regulated QME.
It can be shown that the expression of the anomaly is the value of the finite part in the limit $`\mathrm{\Lambda }^2\mathrm{}`$ of
$$𝒜=\left[(\mathrm{\Delta }W)_R+\frac{i}{2\mathrm{}}(W,W)\right](\mathrm{\Phi },\mathrm{\Phi }^{})$$
(62)
and the regularized value of $`\mathrm{\Delta }W`$ is defined as
$$(\mathrm{\Delta }W)_R\underset{\mathrm{\Lambda }^2\mathrm{}}{lim}\left[\mathrm{\Omega }_0\right]$$
(63)
where
$$\mathrm{\Omega }_0=\left[S_B^A\left(\delta _\mathrm{\Lambda }\right)_C^B\left(ϵ^2\right)_A^C\right].$$
(64)
$`\left(\delta _\mathrm{\Lambda }\right)_B^A`$ is defined by
$`(\delta _\mathrm{\Lambda })_B^A`$ $`=`$ $`\left(\delta _B^A𝒪^{AC}_{CB}\right)^1`$ (65)
$`=`$ $`\delta _B^A+{\displaystyle \underset{n=1}{}}\left(𝒪^{AC}_{CB}\right)^n,`$ (66)
with
$`S_B^A`$ $`=`$ $`{\displaystyle \frac{\delta _r\delta _lS}{\delta \mathrm{\Phi }^B\delta \mathrm{\Phi }_A^{}}},`$ (67)
$`_{AB}`$ $`=`$ $`{\displaystyle \frac{\delta _r\delta _l}{\delta \mathrm{\Phi }^A\delta \mathrm{\Phi }^B}}.`$ (68)
Applying the limit $`\mathrm{\Lambda }^2\mathrm{}`$ in (63), it can be shown that
$$\left(\mathrm{\Delta }S\right)_R\underset{\mathrm{\Lambda }^2\mathrm{}}{lim}\left[\mathrm{\Omega }_0\right]_0,$$
(69)
and finally that
$`𝒜_0`$ $``$ $`\left(\mathrm{\Delta }S\right)_R`$ (70)
$`=`$ $`\underset{\mathrm{\Lambda }^2\mathrm{}}{lim}\left[\mathrm{\Omega }_0\right]_0.`$ (71)
All the higher orders terms of the anomaly can be obtained from equation (62), but this will not be analyzed in this paper. It can be seen in .
## V The Extended Non-local Regularization of the Chiral Schwinger Model
In this section we will make a comparison between the results of the computation of the anomaly of the CSM using two different regulators.
### A Second order differential regulator
The classical action for the chiral Schwinger model is
$`S`$ $`=`$ $`{\displaystyle }d^2x[{\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }+\overline{\psi }i/\psi `$ (73)
$`+{\displaystyle \frac{e}{2}}\overline{\psi }\gamma _\mu (1\gamma _5)A^\mu \psi ],`$
which obviously has a perturbative expansionWe are using the Einstein notation..
This action is invariant under the following gauge transformations:
$`A_\mu (x)`$ $``$ $`A_\mu (x)+_\mu \theta (x)`$ (74)
$`\psi (x)`$ $``$ $`exp\left[ie(1\gamma _5)\theta (x)\right]\psi (x).`$ (75)
The kinetic part of the action (73) is given by
$`F`$ $`=`$ $`{\displaystyle d^2x\overline{\psi }i/\psi }`$ (76)
$`=`$ $`{\displaystyle d^2x\left[\frac{1}{2}\overline{\psi }i/\psi +\frac{1}{2}\overline{\psi }i/\psi \right]}.`$ (77)
Integrating by parts the second term we have that
$$F=d^2x\left[\frac{1}{2}\overline{\psi }i/\psi +\frac{1}{2}\psi (i/^t\overline{\psi })\right].$$
(78)
The kinetic term has the form
$$F=\frac{1}{2}\mathrm{\Psi }^A_{AB}\mathrm{\Psi }^B.$$
(79)
So,
$$\mathrm{\Psi }=\left(\begin{array}{c}\overline{\psi }\\ \psi \end{array}\right)$$
(80)
and
$$F=\frac{1}{2}(\overline{\psi }\psi )\left(\begin{array}{cc}0& i/\\ i/^t& 0\end{array}\right)\left(\begin{array}{c}\overline{\psi }\\ \psi \end{array}\right).$$
(81)
The kinetic operator $`(_{AB})`$ is defined by
$$_{AB}=\left(\begin{array}{cc}0& i/\\ i/^t& 0\end{array}\right).$$
(82)
The regulator, a second order differential operator, is
$$_\beta ^\alpha =(T^1)^{\alpha \gamma }_{\gamma \beta },$$
(83)
where $`T`$ is an arbitrary matrix, hence one can make the following choice:
$$_\beta ^\alpha =^2.$$
(84)
Let us define the smearing operator,
$$ϵ_B^A=exp\left(\frac{^2}{2\mathrm{\Lambda }^2}\right),$$
(85)
and the smeared fields
$$\widehat{\mathrm{\Phi }}^A=(ϵ^1)_B^A\mathrm{\Phi }^B.$$
(86)
In the NLR scheme the shadow kinetic operator is
$$𝒪_{\alpha \beta }^1=\left(\frac{}{ϵ^21}\right)_{\alpha \beta }$$
(87)
then
$$𝒪=\left(\begin{array}{cc}0& i𝒪^{}/\\ i𝒪^{}/^t& 0\end{array}\right)$$
(88)
where $`𝒪^{}`$ is defined by
$`𝒪^{}`$ $`=`$ $`{\displaystyle \frac{ϵ^21}{/^t/}}`$ (89)
$`=`$ $`{\displaystyle _0^1}{\displaystyle \frac{dt}{\mathrm{\Lambda }^2}}exp\left(t{\displaystyle \frac{/^t/}{\mathrm{\Lambda }^2}}\right),`$ (90)
notice that we have to obey the rules for the product of the Dirac matrices $`\gamma _\mu `$.
The interacting part of the action (73) is
$`I[A_\mu ,\psi ,\overline{\psi }]`$ $`=`$ $`{\displaystyle \frac{e}{2}}\overline{\psi }\gamma _\mu (1\gamma _5)A^\mu \psi ,`$ (91)
$`I[A_\mu ,\psi +\mathrm{\Phi },\overline{\psi }+\overline{\mathrm{\Phi }}]`$ $`=`$ $`{\displaystyle \frac{e}{2}}(\overline{\psi }+\overline{\mathrm{\Phi }})\gamma _\mu (1\gamma _5)`$ (93)
$`\times A^\mu (\psi +\mathrm{\Phi }),`$
where $`\mathrm{\Phi }`$ are the shadow fields.
The BRST transformations are given by
$`\delta A_\mu `$ $`=`$ $`_\mu c,`$ (94)
$`\delta \psi `$ $`=`$ $`i(1\gamma _5)\psi c,`$ (95)
$`\delta \overline{\psi }`$ $`=`$ $`i\overline{\psi }(1+\gamma _5)c,`$ (96)
$`\delta c`$ $`=`$ $`0.`$ (97)
Using the equations (40), where the antifields are functions of the auxiliary fields,
$`\psi ^{}`$ $``$ $`\left[\psi ^{}ϵ^2+\mathrm{\Phi }^{}(1ϵ^2)\right],`$ (98)
$`\overline{\psi }^{}`$ $``$ $`\left[\overline{\psi }^{}ϵ^2+\overline{\mathrm{\Phi }}^{}(1ϵ^2)\right].`$ (99)
The generator of the BRST transformations are
$`R(\psi )`$ $``$ $`R(\psi +\mathrm{\Phi })=i(1\gamma _5)(\psi +\mathrm{\Phi })c,`$ (100)
$`R(\overline{\psi })`$ $``$ $`i(\overline{\psi }+\overline{\mathrm{\Phi }})(1+\gamma _5)c,`$ (101)
$`R(c)`$ $`=`$ $`0.`$ (102)
We are able now to construct the non-local auxiliary proper action. It will be given in general by
$$S_\mathrm{\Lambda }(\mathrm{\Phi },\mathrm{\Phi }^{})=\stackrel{~}{S}_\mathrm{\Lambda }(\mathrm{\Phi },\mathrm{\Phi }^{};\psi _s,\psi ^{}=0),$$
(103)
where $`\psi _s`$ are the solutions of the classical equations of motion.
The proper solution, the BV action, is given by
$`S_{BV}`$ $`=`$ $`{\displaystyle }d^2x[{\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }+\overline{\psi }i/\psi `$ (106)
$`+{\displaystyle \frac{e}{2}}\overline{\psi }\gamma _\mu (1\gamma _5)A^\mu \psi +A_\mu ^{}^\mu c`$
$`+i\psi ^{}(1\gamma _5)\psi ci\overline{\psi }^{}\overline{\psi }(1+\gamma _5)c].`$
After a tedious algebra, one can write the non-localized action as
$`\stackrel{~}{S}_\mathrm{\Lambda }(\psi ,\psi ^{})`$ $`=`$ $`{\displaystyle \frac{1}{4}}\widehat{F}_{\mu \nu }\widehat{F}^{\mu \nu }+\widehat{\overline{\psi }}i/\widehat{\psi }+A_\mu ^{}^\mu c`$ (107)
$`+`$ $`{\displaystyle \frac{e}{2}}{\displaystyle \frac{(i/)\left[\overline{\psi }\gamma ^\mu (1\gamma _5)A_\mu \psi \right]}{i/+e\gamma ^\mu (1\gamma _5)A_\mu (ϵ^21)}}`$ (108)
$`+`$ $`{\displaystyle \frac{i\psi ^{}ϵ^2c(i/)(1\gamma _5)\psi }{[i/+e\gamma ^\mu (1\gamma _5)A_\mu (ϵ^21)]}}`$ (109)
$`+`$ $`{\displaystyle \frac{i\overline{\psi }^{}ϵ^2c(i/)\overline{\psi }(1+\gamma _5)}{[i/+e\gamma ^\mu (1\gamma _5)A_\mu (ϵ^21)]}}.`$ (110)
It can be easily seen that when one take the limit $`ϵ^21`$, the original proper solution $`S_{BV}`$ of the CSM is obtained. This is a representative expression, since it is well known that operators in the denominator of any expression are physically senseless.
The final part is the computation of the one-loop anomaly of the chiral Schwinger model. Firstly, we have to construct some very important matrices,
$$S_B^A=\frac{\delta _r\delta _lS_{BV}}{\delta \mathrm{\Phi }^B\delta \mathrm{\Phi }_A^{}}$$
(111)
Then
$$S_B^A=\left(\begin{array}{cc}ic(1\gamma _5)& 0\\ 0& ic(1+\gamma _5)\end{array}\right).$$
(112)
The operator $`_{AB}`$ in this case is defined by,
$$_{AB}=\frac{\delta _l\delta _r\left[I(\mathrm{\Phi })+\mathrm{\Phi }_c^{}R^c(\mathrm{\Phi })\right]}{\delta \mathrm{\Phi }^A\delta \mathrm{\Phi }^B}$$
(113)
and the result is,
$$_{AB}=\left(\begin{array}{cc}0& \frac{e}{2}\gamma _\mu (1\gamma _5)A^\mu \\ \frac{e}{2}\gamma _\mu (1\gamma _5)A^\mu & 0\end{array}\right).$$
(114)
The one-loop anomaly is given by:
$`𝒜`$ $``$ $`(\mathrm{\Delta }S)_R,`$ (115)
$`(\mathrm{\Delta }S)_R`$ $`=`$ $`\underset{\mathrm{\Lambda }^2\mathrm{}}{lim}[\mathrm{\Omega }_0]_0,`$ (116)
$`\mathrm{\Omega }_0`$ $`=`$ $`\left[ϵ^2S_A^A\right]+\left[ϵ^2S_B^A𝒪^{BC}_{CA}\right]`$ (117)
$`+`$ $`O\left({\displaystyle \frac{(\mathrm{\Phi }^{})^2}{\mathrm{\Lambda }^2}}\right).`$ (118)
For the first term we can compute that
$`ϵ^2S_A^A`$ $`=`$ $`ϵ^2trS_B^A`$ (119)
$`=`$ $`0,`$ (120)
now we have that
$$(\mathrm{\Delta }S)_R=\underset{\mathrm{\Lambda }^2\mathrm{}}{lim}tr\left[ϵ^2S_B^A𝒪^{BC}_{CA}\right].$$
(121)
Using the Weyl representation of the $`\gamma `$ matrices in two dimensions in Euclidian space:
$`\gamma _0`$ $`=`$ $`\left(\begin{array}{cc}0& \mathrm{\hspace{0.17em}1}\\ \mathrm{\hspace{0.17em}1}& 0\end{array}\right),`$ (124)
$`\gamma _1`$ $`=`$ $`\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),`$ (128)
$`\gamma ^5`$ $`=`$ $`i\gamma _1\gamma _0,`$ (130)
and that in this representation, $`\gamma _5^t=\gamma _5`$.
Finally, after some algebra
$`(\mathrm{\Delta }S)_R`$ $`=`$ $`\underset{\mathrm{\Lambda }^2\mathrm{}}{lim}tr[ϵ^2(ec)𝒪^{}(\gamma ^\nu )^t\gamma ^\mu `$ (132)
$`\times (\mathrm{\hspace{0.17em}1}\gamma _5)_\nu A_\mu ]`$
$`=`$ $`\underset{\mathrm{\Lambda }^2\mathrm{}}{lim}tr[ϵ^2(ec){\displaystyle \frac{ϵ^21}{^2}}`$ (134)
$`\times (_\mu A^\mu iϵ^{\mu \nu }_\mu A_\nu )].`$
But we know that the functional traces can be written as
$`\underset{\mathrm{\Lambda }^2\mathrm{}}{lim}tr\left[ϵ^2F^n{\displaystyle \frac{ϵ^21}{^2}}G^m\right]=`$ (135)
$`=`$ $`{\displaystyle \frac{i}{2\pi }}\left[{\displaystyle \underset{k=0}{\overset{m}{}}}\left(\begin{array}{c}m\\ k\end{array}\right){\displaystyle \frac{(1)^k}{n+m+1k}}\left(1{\displaystyle \frac{1}{2^{n+m+1k}}}\right)\right]`$ (138)
$`\times `$ $`{\displaystyle d^2xF^{n+m+1}G},`$ (139)
this last equation can be derived from (89) after hard work where convenient reparametrization were necessary.
In our case
$`n`$ $`=`$ $`m=\mathrm{\hspace{0.33em}0}`$ (140)
$`F`$ $`=`$ $`ec`$ (141)
$`G`$ $`=`$ $`_\mu A^\mu iϵ^{\mu \nu }_\mu A_\nu ,`$ (142)
and the final result is
$$𝒜=(\mathrm{\Delta }S)_R=\frac{ie}{4\pi }d^2xc\left(_\mu A^\mu iϵ^{\mu \nu }_\mu A_\nu \right)$$
(143)
which is exactly the one-loop anomaly of the chiral Schwinger model, action (73).
### B The Fujikawa regulator (FR)
We will study now the utilization of a second regulator, a FR to get the same result as was obtained above. Originally, the FR is a covariant derivative that was used by Fujikawa to calculate the chiral anomaly . Only the main steps of the calculation are presented here.
Notice that now we know the value of the anomaly. We will show here that there is another different regulator that furnish the same result. To do this let us construct a convenient FR given by
$$_\beta ^\alpha =D/=\gamma _\mu (^\mu +A^\mu ),$$
(144)
introducing this regulator à la Fujikawa in (132) we can write that
$`(\mathrm{\Delta }S)_R`$ $`=`$ (145)
$`=`$ $`\underset{\mathrm{\Lambda }^2\mathrm{}}{lim}tr\{ϵ^2(ec){\displaystyle }{\displaystyle \frac{d^2k}{(2\pi )^2}}e^{ikx}`$ (146)
$`\times `$ $`{\displaystyle _0^1}{\displaystyle \frac{dt}{\mathrm{\Lambda }^2}}exp\left(t{\displaystyle \frac{//^t}{\mathrm{\Lambda }^2}}\right)`$ (147)
$`\times `$ $`(\gamma ^\nu )^t\gamma ^\mu (\mathrm{\hspace{0.17em}1}\gamma _5)_\nu A_\mu exp({\displaystyle \frac{D/^{\mathrm{\hspace{0.17em}2}}}{\mathrm{\Lambda }^2}})e^{ikx}\}`$ (148)
$`=`$ $`\underset{\mathrm{\Lambda }^2\mathrm{}}{lim}tr\{ϵ^2(ec){\displaystyle _0^1}{\displaystyle \frac{dt}{\mathrm{\Lambda }^2}}exp\left(t{\displaystyle \frac{//^t}{\mathrm{\Lambda }^2}}\right)`$ (149)
$`\times `$ $`(\gamma ^\nu )^t\gamma ^\mu (\mathrm{\hspace{0.17em}1}\gamma _5)_\nu A_\mu `$ (150)
$`\times `$ $`exp\left({\displaystyle \frac{1}{\mathrm{\Lambda }^2}}\left[A^{\mathrm{\hspace{0.17em}2}}+^\mu A_\mu +{\displaystyle \frac{1}{4}}[\gamma ^\mu ,\gamma ^\nu ]F_{\mu \nu }\right]\right)`$ (151)
$`\times `$ $`{\displaystyle }{\displaystyle \frac{d^2k}{(2\pi )^2}}exp[{\displaystyle \frac{1}{\mathrm{\Lambda }^2}}(k^2+2iA_\mu k^\mu )]\}.`$ (152)
Notice that we have substituted the $`𝒪^{}`$ operator by its integral form.
It is necessary to make a usefull transformation of the $`k^\mu `$ coordinate,
$$k^\mu \mathrm{\Lambda }k^\mu .$$
(153)
Expanding the exponential of $`t`$ it is easy to realize that only the unitary term of the expansion can be used, since the other terms have the $`\mathrm{\Lambda }^n`$ form, and hence disappear with the infinity limit. After the computation of the integrals, we have that,
$$(\mathrm{\Delta }S)_R=\frac{iec}{4\pi }tr\left\{(\gamma ^\nu )^t\gamma ^\mu (\mathrm{\hspace{0.17em}1}\gamma _5)_\nu A_\mu \right\}.$$
(154)
Manipulating with the $`\gamma `$ matrices (Weyl representation), the final result is
$$𝒜=(\mathrm{\Delta }S)_R=\frac{ie}{4\pi }d^2xc\left(_\mu A^\mu iϵ^{\mu \nu }_\mu A_\nu \right),$$
(155)
which is the same result as we have obtained before with a different regulator.
As has been said before, it is a very interesting result because it was expected that the expressions for the anomaly would be not equal, as has been showed in and .
## VI Conclusions
The non-local regularization formalism is a recent and a quite powerfull method to regularize theories with a perturbative expansion which have higher order loop divergences. The field-antifield framework exhibits a divergence in the application of the $`\mathrm{\Delta }`$ operator. Hence it needs a regularization. The connection between the BV formalism and the NLR method generates an extended non-locally regularized BV quantization method. At the quantum level its use in the path integral originates this extended BV formalism through the construction of a non-local regularized quantum action. In this way we can compute higher order loop of the BRST anomaly that is contained in the quantum master equation. To make this connection, we have to introduce auxiliary fields with which we have constructed an auxiliary proper solution of the master equation. These auxiliary fields are eliminated from the theory through the field equations. At this point, and with some technical work, we can construct a regularized $`\mathrm{\Delta }S`$ expression that furnishes the final form of the anomaly.
In this work, the theory used was a fermionic one, the chiral Schwinger model. The objective is to analyze the results of the anomaly using two different regulators: a second order differential one and a kind of covariant derivative operator, like the one that was used by Fujikawa to calculate the chiral anomaly.
With the second regulator, at a certain point of the process, the non-local characteristic of the method, contained in the $`𝒪^{}`$ operator, was substituted by its integral form and furthermore it has been combined with the Fujikawa formalism to give an exact result.
We know that some anomalous theories can have the anomaly calculation dependent on the regulator that has been used. It was shown in the literature that the solutions for the CSM have different dependences on the parameters of the regulators. However, these parameters are free to be chosen, causing no conflict between the results. The final form of the anomaly can not be obtained in a direct way, and consequently, we have to make a reparametrization to obtain the final answer. In another way it was also demonstrated that we have to introduce a local counterterm, the WZ term, to obtain the equivalence between both different expressions obtained for the anomaly. In this work, we has shown in a precise way that the calculation of the one-loop anomaly of the CSM with different regulators has furnished directly the same final results. Finally, we can observe that it would be very interesting to make the same analysis for theories in higher dimensions.
Acknowledgment: the author would like to thank Nelson R. F. Braga and Clovis Wotzasek for valuable discussions and suggestions. This work is supported by CAPES (Brazilian Research Agency). |
warning/0001/nucl-th0001028.html | ar5iv | text | # Untitled Document
INERTIAL PARAMETERS AND SUPERFLUID-TO-NORMAL
PHASE TRANSITION IN SUPERDEFORMED BANDS
I. M. Pavlichenkov
Russian Research Center ”Kurchatov Institute”, Moscow, 123182, Russia
##
The quasiclassically exact solution for the second inertial parameter $``$ is found in a self-consistent way. It is shown that superdeformation and nonuniform pairing arising from the rotation induced pair density significantly reduce this inertial parameter. The new signature for the transition from pairing to normal phase is suggested in terms of the variation $`/𝒜`$ versus spin. Experimental data indicate the existence of such transition in the three SD mass regions.
PACS numbers: 21.10.Re, 21.60.Fw, 23.20.Lv, 27.60.+j
One of the amazing features of superdeformed (SD) rotational bands is the extreme regularity of their rotational spectra: a SD nucleus is the best quantum rotor known in nature. In spite of the fact that numerous theoretical calculations successfully reproduce the measured intraband $`\gamma `$-ray energies (see e.g. \[1–4\]) the underlying microscopic mechanism of this phenomenon is still not well understood. To explain the SD rotational spectrum regularity we use in the present paper the parameterization of its energy by the three term formula
$$E(I)=E_0+𝒜I(I+1)+I^2(I+1)^2,$$
(1)
which is valid for an axially symmetric deformed nucleus with $`K=0`$. The inertial parameters $`𝒜=\mathrm{}^2/2\mathrm{}^{(1)}`$ ($`\mathrm{}^{(1)}`$ is the kinematic moment of inertia) and $``$ are the objects of our investigation. They are determined by the transition energies $`E_\gamma (I)=E(I+2)E(I)`$ as follows
$$𝒜(I)=\frac{1}{4(2I+5)}[\frac{I^2+7I+13}{2I+3}E_\gamma (I)\frac{I^2+3I+3}{2I+7}E_\gamma (I+2)],$$
$$(I)=\frac{1}{8(2I+5)}\left[\frac{E_\gamma (I+2)}{2I+7}\frac{E_\gamma (I)}{2I+3}\right].$$
(2)
The coefficient $`(I)`$ characterizes the nonadiabatic properties of a band and is very sensitive to its internal structure. In particular, it realizes the relationship of kinematic and dynamic ($`\mathrm{}^{(2)}`$) moments of inertia. Using the well known expressions for these values and the formulas (2) we get
$$=\frac{\mathrm{}^2}{2(2I+3)(2I+7)}\left[\frac{1}{\mathrm{}^{(2)}}\frac{2I}{(2I+5)\mathrm{}^{(1)}}\right].$$
(3)
The ratio $`/𝒜`$ determines the convergence radius of the two parameter formula (1), which is of order 100 for the bands in the 80 and 150 mass regions and 40 in the 130 and 190 ones. Faster convergence is obtained with Harris formula
$$E(\omega )=E_0+\frac{1}{2}\alpha \omega ^2+\frac{3}{4}\beta \omega ^4,$$
(4)
which is based on the fourth-order cranking expansion
$$\alpha =\frac{1}{\omega }Sp(\mathrm{}_x\rho ^{(1)}),\beta =\frac{1}{\omega ^3}Sp(\mathrm{}_x\rho ^{(3)}),$$
(5)
where $`\rho ^{(n)}`$ is the $`n`$th correction to the nucleus density matrix, $`\mathrm{}_x`$ is the single-particle angular momentum projection on the rotational axis $`x`$ perpendicular to the symmetry axis $`z`$, and $`\omega `$ is the rotational frequency. It follows from Eqs. (1) and (4) that
$$\alpha =\mathrm{}^2/2𝒜,\beta =\mathrm{}^4/4𝒜^4.$$
(6)
For simplicity we will deal with the parameter $`\beta `$.
The problem of the microscopic calculation of the parameter $``$ for normal deformed (ND) nuclei has attracted considerable attention. Its value is formed mainly by four effects: vibration-rotation interaction, centrifugal stretching, perturbation of the quasiparticle motion, and attenuation of pair correlation by the Coriolis force (Mottelson-Valatin effect). As it has already been shown in the first attempts at obtaining $``$ , the latter two effects are dominant for well deformed nuclei. Unfortunately the results of these and many other works cannot be used for the superdeformation. The formulas of Ref. have been obtained in the limit of the monopole pairing interaction (the uniform pairing), which is not adequate at SD shapes as shown in Ref. . In the more sophisticated work the gauge invariant pairing interaction allows to study the effect of nonuniform pairing. However this approach is also limited because it neglects the coupling between major shells. Thus, despite a number of publications on the subject the correct cranking selfconsistent solution for the $``$ coefficient has not been found.
We used the quasiclassical method of Ref. to derive the following expression for the $`\beta `$ parameter in the superfluid phase
$$\beta _s=\frac{\mathrm{}^4}{4\mathrm{\Delta }^2}\mathrm{}_{12}^x\mathrm{}_{23}^x\mathrm{}_{34}^x\mathrm{}_{41}^xF(x_{12},x_{23},x_{34},x_{41})\delta (\epsilon _1\epsilon _F),$$
(7)
where the summation indices $`i`$=1,2,3,4 refer to the single-particle states $`i`$ of the nonrotating mean field with the energy $`\epsilon _i`$. The $`\delta `$-function means that the summation over the states 1 is restricted by a small interval at the Fermi energy $`\epsilon _F`$ . The dimensionless values $`x_{ii^{}}=(\epsilon _i\epsilon _i^{})/2\mathrm{\Delta }`$, where $`\mathrm{\Delta }`$ is the state independent pairing gap at $`\omega =0`$, correspond to energy differences between states permitted by the selection rules for the matrix element of $`\mathrm{}_x`$. The function $`F`$ depending on these values may be written as follows
$$F=\underset{k=0}{\overset{3}{}}\widehat{P}_kG_{12}+\underset{k=0}{\overset{1}{}}\widehat{P}_kH_{13}8D_2^2x_{12}x_{23}x_{34}x_{41}h(x_{13}),$$
(8)
where the permutation operators in the space of four indices $`i`$, $`imod4=i`$,
$$\widehat{P}_kx_{i,i^{}}=x_{i+k,i^{}+k},$$
(9)
are used to simlify the formulas. The expressions for $`G_{12}`$ and $`H_{13}`$ involve the well known functions
$$h(x)=(1+x^2)g(x),g(x)=\frac{argshx}{x\sqrt{1+x^2}},$$
(10)
and have the form:
$$G_{12}=\frac{g(x_{12})}{x_{23}x_{41}x_{13}x_{24}}\{(1D_1x_{12}^2)[1x_{12}^2x_{23}x_{41}+D_1[x_{23}^2(1x_{12}x_{24})+x_{41}^2(1+x_{12}x_{13})]$$
$$+D_1^2x_{23}x_{34}x_{41}(x_{23}+x_{41})D_1^3x_{12}x_{23}^2x_{34}x_{41}^2]+D_1(x_{34}D_1x_{12}x_{23}x_{41})(x_{34}$$
$$+x_{12}x_{13}x_{24}D_1x_{12}x_{23}x_{41})\},$$
$$H_{13}=\frac{h(x_{13})}{x_{12}x_{23}x_{34}x_{41}}[1D_1(x_{12}^2+x_{23}^2+x_{34}^2+x_{41}^2)^2+D_1^2(x_{12}x_{41}+x_{23}x_{34})],$$
(11)
Eq. (7) multiplied by $`\omega ^3`$ is the third order cranking correction to the total angular momentum of the neutron or proton system. Its derivation will be described in a forthcoming paper. We want now to emphasize that Eq. (7) represents the first theoretically correct expression for the high order effect of the Coriolis-pairing interaction at fixed deformation. The result is obtained by taking into account the effect of rotation on the Cooper pairs in the gauge invariant form. This effect is described by the first and the second corrections to the pairing energy
$$\mathrm{\Delta }^{(1)}(𝐫)=\frac{i\mathrm{}^2\omega }{2\mathrm{\Delta }}D_1\dot{\mathrm{}}_x,\mathrm{\Delta }^{(2)}(𝐫)=\frac{\mathrm{}^4\omega ^2}{4\mathrm{\Delta }^3}D_2\dot{\mathrm{}}_x^2,$$
(12)
where $`D_1`$ and $`D_2`$ are the amplitudes of the nonuniform pairing fields, which are found in a self-consistent way. Let us note that $`\mathrm{\Delta }^{(n)}`$ is the function of the space coordinates because $`\mathrm{}\dot{\mathrm{}}_x=yU/zzU/y`$, where $`U`$ is the potential of a mean field. It is seen that for the oscillator potential $`\dot{\mathrm{}}_xyz`$ (the $`Y_{21}`$ pairing). The coordinate dependent pairing field is crucial for conservation of a nucleon current. The theory incorporating the nonuniform pairing allows also to consider the different limiting cases for the inertial parameters, which make possible the study of an interplay between rotation, pairing correlations, and mean field deformation in a SD band.
In order to consider this problem quantitatively we will use the axially deformed oscillator potential with the frequencies $`\omega _x=\omega _y`$ and $`\omega _z`$ on the corresponding axes. In this model the matrix element $`\mathrm{}_{12}^x`$ is non-zero for the two types of transitions:
(i) transitions inside a single oscillator shell (close transitions), for which $`x_{12}=\pm \nu _1`$; (ii) transitions over a shell (distant transitions) with $`x_{12}=\pm \nu _2`$. The quantities $`\nu _1`$ and $`\nu _2`$ are the well known parameters involved in the moment of inertia :
$$\nu _{1,2}=\frac{\mathrm{}(\omega _x\omega _z)}{2\mathrm{\Delta }}=\frac{k1}{2\xi k^{2/3}},\xi =\frac{\mathrm{\Delta }}{\mathrm{}\omega _{}},$$
(13)
where $`\mathrm{}\omega _{}=41A^{1/3}`$ MeV. Here and later we use the axis or frequency ratio $`k=c/a=\omega _x/\omega _z`$ and the volume conservation condition. Both of the values $`\nu _1`$ and $`\nu _2`$ are large for superdeformation. For the fixed state 1, there are 36 different combinations of these base transitions in the sum of Eq. (7). The summation over the states 1 is performed in the Thomas-Fermi approximation. The final expression for the parameter $`\beta _s`$ in the oscillator potential is
$$\beta _s=\frac{(k+1)^4}{1875\mathrm{}^2k^{4/3}}AM^3R^6\mathrm{\Phi }(\xi ,k),$$
(14)
where $`R=1.2A^{1/3}`$ fm is the radius of the sphere, which volume is equal to that of the spheroid with the half-axes $`a<c`$, $`M`$ is the nucleon mass, and A is the number of nucleons. The function $`\mathrm{\Phi }`$ along with its limiting cases is shown in Fig. 1. It is seen that nonuniform pairing reduces substantially the parameter $`\beta _s`$ in agreement with the estimation of Ref. . On the other hand, the contribution of distant transitions is minor for small $`\xi `$. Nevertheless the later are necessary to obtain the hydrodynamic limit (see below). Since $`\mathrm{\Phi }1`$ for a reasonable pairing gap, $`\mathrm{\Delta }`$0.5 MeV, the order of the value $`\beta _s`$ is $`\mathrm{}^4(A/\epsilon _F)^3`$. This, along with the estimation $`𝒜\epsilon _FA^{5/3}`$, gives $`/𝒜A^2`$, which overestimates the minimal value of this ratio in all the SD mass regions. Thus a small $`\mathrm{\Delta }`$ and nonuniform pairing does not solve the problem of the SD band regularity.
Let us consider first the limiting case $`\mathrm{\Delta }=0`$. The right part of Eq. (7) vanishes for noncorrelated nucleons. This result is the artifact of the quasiclassical approximation used in Eq. (7). The correct expression for the $`\beta `$ parameter in the normal phase, obtained with the limiting values of the Bogolubov amplitudes ($`u_i=0`$, $`v_i=1`$ for $`n_i=1`$ and $`u_i=1`$, $`v_i=0`$ for $`n_i=0`$, where $`n_i`$ is the nucleon occupation numbers), has the form
$$\beta _n=\mathrm{}^4\mathrm{}_{12}^x\mathrm{}_{23}^x\mathrm{}_{34}^x\mathrm{}_{41}^x\underset{k=0}{\overset{3}{}}\widehat{P}_k\left\{\frac{n_1}{\epsilon _{12}\epsilon _{13}\epsilon _{14}}\right\}.$$
(15)
The odd function of the differences $`\epsilon _{ii^{}}=\epsilon _i\epsilon _i^{}`$ leads to the cancelation of the main terms in sum (15) that decreases substantially the value of $`\beta _n`$, $`\beta _n\mathrm{}^4A^{7/3}/\epsilon _F^3`$. Note that the centrifugal stretching effect has the same order $`\beta _{str}\beta _n`$. Its contribution is small compared to that of $`\beta _s`$, but it should not be overlooked for an unpaired system. For the oscillator potential, we have the following expression
$$\beta _n+\beta _{str}=\frac{k^410k^2+1}{6\omega _{}^2k^{4/3}}AMR^2.$$
(16)
It is seen that $`\beta _n+\beta _{str}<0`$ for the prolate nuclei with $`c/a<3.15`$, whereas $`\beta _s`$ is always positive. Thus, with an increase of the spin $`I`$, the ratio $`/𝒜`$ has to change sign and to approach its limiting value $`_n/𝒜_nA^{8/3}`$, $`10^6`$ for the SD bands in the 130 and 150 mass region, where rapid rotation destroys pairing correlations. The limiting ratio for a nucleus consisting of $`Z`$ protons and $`N`$ neutrons is expressed by
$$\frac{_n}{𝒜_n}=2.56\frac{(k^410k^2+1)k^{2/3}}{(k^2+1)^3A^{8/3}}\left[\left(\frac{2Z}{A}\right)^{1/3}+\left(\frac{2N}{A}\right)^{1/3}\right].$$
(17)
One can therefore conclude that there are two distinct regions in the variation of $`/𝒜`$ versus $`I`$. The lower part of a SD band is characterized by a gradual decrease of the pairing gap $`\mathrm{\Delta }`$. According to Eq. (14) the ratio $`/𝒜`$ should exhibit a sharp increase. Then it changes sign and approaches the plateau (17) at the top of a band because the deformation $`c/a`$ depends weakly on spin in the normal phase. Such behavior of the $`/𝒜`$ ratio is the signature of the pairing phase transition.
We have analyzed all the SD bands of Ref. with known or suggested spins of levels. Figure 2 shows the variation of the $`/𝒜`$ ratio with $`I`$ for bands with different internal structure and different rotational frequencies. Apart from the bands <sup>192</sup>Hg(1) and <sup>194</sup>Hg(3), where frequencies are so low that $`/𝒜`$ rises continuously in the superfluid phase, and <sup>84</sup>Zr(1), for which pairing is quenched completely and $`/𝒜`$ is close to the limiting value (17), all other bands display the behavior described above. It is important to note that such behavior is observed for the ND yrast band of <sup>84</sup>Zr, where $`/𝒜`$ reaches the same limiting value (17) as in the SD band <sup>84</sup>Zr(1). There are other ND yrast bands of <sup>168</sup>Yb and <sup>168</sup>Hf with the phase transition which experimental evidence has been discussed previously in terms of the canonical variables and the spectrum of single-particle states . These bands exhibit the same features. Thus the plots of Fig. 2 demonstrate the universality of the superfluid-to-normal phase transition for SD and ND bands. The manifestation of this universality was previously observed in the anomalous small value $`/𝒜7\times 10^6`$ for the ND yrast band of <sup>168</sup>Hf and in the weak dependence of $`\mathrm{}^{(1)}`$ in the upper part of the SD band <sup>152</sup>Dy(1) . The new feature observed in this article is the small decrease of $`/𝒜`$ at the top of the <sup>152</sup>Dy(1), <sup>132</sup>Ce(1), <sup>84</sup>Zr(1) and other SD bands. It may be explained by the decrease of the first multiplier in r.h.s. of Eq. (17) due to increase of nuclear deformation at highest spins.
All of the theoretical formulas obtained refer to collective rotation only. The influence of single-particle degrees of freedom on the $`/𝒜`$ ratio is the subject of a separate investigation. Preliminary estimations show that the contribution of an odd nucleon in $`_s`$ may be positive and irregular. Visible evidences of the nonrotational degree of freedom are shown in Fig. 3. A point of particular interest is the staggering band <sup>149</sup>Gd(1). It is seen that the staggering pattern is reproduced very well and the error bars are smaller than the staggering amplitude. We would like to emphasize that the physical values adequate for describing the staggering phenomenon are the $``$ parameter or the $`/𝒜`$ ratio. They do not require any smooth reference, which introduces an ambiguousness in the amplitude and the phase of oscillations.
The next limit we want to consider is that of a large pairing gap $`\mathrm{\Delta }`$. In this case, the nonuniform pairing (12) is essential and the leading terms in the function $`\mathrm{\Phi }`$ are those proportional to the powers of $`D_1`$ and $`D_2^2`$. They result in the limiting expression $`\mathrm{\Phi }(\mathrm{}\omega _{}/\mathrm{\Delta })^2`$. Thus for the very strong pairing ($`\mathrm{\Delta }\mathrm{}\omega _{}`$), when the size of the Cooper pair $`R\mathrm{}\omega _{}/\mathrm{\Delta }`$ becomes much less than the nuclear radius, the value $`_s`$ vanishes in agreement with the hydrodynamic equations of the ideal liquid . In the limit of an extremely large deformation, $`c/a\mathrm{}`$, a needle shaped nucleus with pairing correlations rotates as a rigid body, $`\mathrm{}=\mathrm{}_{rig}`$, $`_s=0`$. For the finite but large deformation the deviations from these values are proportional to $`(a/c)^{4/3}`$. This means that all nucleons with the exclusion of a small sphere in the center of a nucleus are completely involved in rotational motion. Finally, for small deformations we have $`_s(c/a1)^6`$ that is comparable to the vibration-rotation interaction .
Unlike the limiting value (17), it is impossible to compare with experiment the ratio $`_s/𝒜_s`$ because the proton ($`\mathrm{\Delta }_\pi `$) and neutron ($`\mathrm{\Delta }_\nu `$) pairing gaps are unknown for the SD bands. In such a case we try to solve an inverse problem. The equations for the two inertial parameters
$$\alpha A/\mathrm{}_{rig}=Z\phi (\xi _\pi ,k)+N\phi (\xi _\nu ,k),$$
(18)
$$\beta A/\beta _{}=Z\mathrm{\Phi }(\xi _\pi ,k)+N\mathrm{\Phi }(\xi _\nu ,k),$$
(19)
allow in principle to find $`\xi _\pi `$ and $`\xi _\nu `$ ($`\xi _l=\mathrm{\Delta }_l/\mathrm{}\omega _l`$, $`\omega _l=\omega _{}(2A_l/A)^{2/3}`$, $`l=\pi ,\nu `$). Here the function $`\phi `$ is taken from Ref. ,
$$\phi (\nu _1,\nu _2)=1\frac{(\nu _1^2\nu _2^2)^2g(\nu _1)g(\nu _2)}{(\nu _1^2+\nu _2^2)[\nu _1^2g(\nu _1)+\nu _2^2g(\nu _2)]},$$
(20)
and the rigid body moment of inertia and the value $`\beta _{}`$ have the form:
$$\mathrm{}_{rig}=7.29A^{5/3}(k^2+1)k^{2/3}10^3\mathrm{}^2MeV^1,$$
$$\beta _{}=2.59A^4(k+1)^4k^{4/3}10^8\mathrm{}^4MeV^3.$$
(21)
Finally the axis ratio $`k=c/a`$ is determined by the quadrupole moment
$$Q_0=6.05A^{2/3}(k^21)k^{2/3}10^3eb.$$
(22)
Unfortunately the system (18)–(19) does not have a solution. Table I gives the solutions of Eqs. (18) and (19) under simplifying assumption that the neutron and proton pairing gaps are equal. Different bands reveal the same tendency. The obtained values $`\mathrm{\Delta }_1`$ and $`\mathrm{\Delta }_2`$ show that the oscillator potential overestimates the moment of inertia and, to a greater extent, the $`\beta `$ parameter. It is essential to use a more realistic potential for extracting pairing gaps from the inertial parameters.
Table I. Pairing gap energies $`\mathrm{\Delta }_1`$ and $`\mathrm{\Delta }_2`$ (in MeV) found as the independent solutions of Eqs. (18) and (19) correspondingly in the case of $`\xi _\pi =\xi _\nu `$. The values $`c/a`$, $`\mathrm{}_{rig}`$, and $`\beta _{}`$ were calculated by using Eqs. (21) and (22). The experimental values of $`\alpha =\mathrm{}`$ and $`\beta `$ were taken from Refs. (<sup>172</sup>Hf(yr)) and (<sup>192,4</sup>Hg(1,2)). Those of <sup>236</sup>U(fi) were determined from Eqs. (2), (6).
| Band | $`c/a`$ | $`\mathrm{}/\mathrm{}_{rig}`$ | $`\mathrm{\Delta }_1`$ | $`\beta /\beta _0`$ | $`\mathrm{\Delta }_2`$ |
| --- | --- | --- | --- | --- | --- |
| <sup>172</sup>Hf(yr) | 1.32 | 0.373 | 1.44 | 0.309 | 0.02 |
| <sup>192</sup>Hg(1) | 1.61 | 0.762 | 0.95 | 0.151 | 0.05 |
| <sup>194</sup>Hg(1) | 1.60 | 0.793 | 0.83 | 0.112 | 0.04 |
| <sup>236</sup>U(fi) | 1.84 | 0.816 | 1.01 | 0.141 | 0.08 |
In summary, the exact solution for the inertial parameter $``$ in the superfluid phase allows to show that neither superdeformation nor nonuniform pairing arising from rotation induced pair density is responsible for the regularity of the SD rotational spectra. The extreme regularity of the SD bands in the 80, 130 and 150 mass regions is explained by the transition from the superfluid to normal phase. The new signature of this transition reveals itself in the characteristic dependence of the ratio $`/𝒜`$ with the spin I. Application of this criterion to experimental data indicates the existence of the phase transition in the SD bands of the three SD mass regions. A new extraction method of the proton and neutron pairing gap from the $`\gamma `$-ray energies of interband transitions is discussed.
FIGURE CAPTIONS
Fig. 1. Plot of the function $`\mathrm{\Phi }`$ from Eq. (14) against the dimensionless value $`\xi `$ for the axis ratio $`c/a=2`$. The solid, dotted, and dashed lines correspond respectively to the exact value, the limit of close transitions, and the uniform pairing. The scale on the abscissa should be multiplied by a factor of approximately 7.7 for nuclei in the $`A`$ 150 mass region to obtain a gap energy in MeV.
Fig. 2. $`/𝒜`$ ratio versus spin for the SD and ND bands with mainly collective behavior. The formulas (2) are used to obtain this ratio from the experimental data for the ND (full circles) and SD (open circles) bands. The solid straight line is the limiting value $`_n/𝒜_n`$ with the deformation $`c/a`$ found from the quadrupole moment (22). Error bars (if they are greater than symbols) include $`\gamma `$-ray energy uncertainties only. Uncertainties in a spin assignment are not important for all the SD bands except <sup>152</sup>Dy(1), as the variation of spins in 2$`\mathrm{}`$ would merely shift the curves along axis of abscissas.
Fig. 3. Same as Fig. 2 for the bands with high-$`N`$ configurations having the non-zero alignment $`i`$. Straight line is the value of $`_n/𝒜_n`$ for <sup>153</sup>Dy(1). |
warning/0001/hep-th0001034.html | ar5iv | text | # Vacuum Polarization and Energy Conditions at a Planar Frequency Dependent Dielectric to Vacuum Interface
## I Introduction
The renormalized vacuum stress-tensor for quantized fields at perfectly reflecting boundaries has been extensively studied, sufficiently so that they have made their way into various texts. Fulling and Birrell and Davies are two examples. The perfect reflector enforces a Dirichlet boundary condition that must be satisfied by the quantized field. For the case of the scalar or electromagnetic field with an infinite planar boundary, the renormalized vacuum energy density for stationary observers has the form
$$\rho _{\text{Ren.}}=\frac{\alpha }{x^m},$$
(1)
where $`\alpha `$ is a unitless constant of order unity, $`x`$ is the transverse distance from the planar boundary and $`m>1`$ is the dimension of the spacetime. The most notable property of the energy density is the divergence in the limit $`x0`$.
The divergence is not at all unexpected. The application of the Dirichlet boundary condition fixes the value of the field at a specific point with no uncertainty in the position. Thus the uncertainty of the conjugate momenta must be infinite, giving rise to the stress-tensor divergence. The precise mathematical origin of the divergence comes from the renormalization of the vacuum energy via the Green’s function. The argument is simply that the unrenormalized Green’s function is constructed such that it vanishes on the boundary. To renormalize the Green’s function, the corresponding Minkowski space Green’s function must be subtracted away. Since this is divergent in the limit $`x0`$, the subtraction carried out for renormalization will leave behind a divergent part. When the vacuum energy is calculated from the renormalized Green’s function, it will diverge at $`x=0`$ also.
The divergence of the vacuum energy density causes considerable consternation since the vacuum would contain an infinite amount of energy. In a self consistent theory of gravity, the back reaction of the vacuum energy on the spacetime would be considerable. Yet we know from experience that metallic foils do not cause significant gravitational distortions around themselves.
There are three resolutions to the problems: (1) Replace the Dirichlet boundary condition with something more realistic, which is the topic of this paper. (2) Add position uncertainty to the mirror; this has been discussed by Ford and Svaiter . (3) Treat the entire problem quantum mechanically, which has been partially touched upon by Barton and Eberlein and more recently by Helfer .
Replacing the Dirichlet boundary condition has to some degree been discussed by other authors. For example Candelas has derived the asymptotic form of the vacuum energy density for the problem of two nondispersive dielectric half-spaces which are in contact with one another along a planar interface. More recently Helfer and Lang carried out similar calculations for a nondispersive dielectric to vacuum interface, with considerable interest placed upon the energy conditions in the vacuum region.
This paper addresses the behavior of the stress-tensor for a quantized massless scalar field at a planar dielectric to vacuum interface where the index of refraction of the material varies as a function of frequency. The Dirichlet boundary condition is replaced with appropriate frequency dependent reflection and transmission coefficients. The stress-tensor in the vacuum region is well behaved, even on the interface, and is largely independent of the functional dependence of the index of refraction. Therefore the vacuum energy divergence is a result of the Dirichlet boundary condition and not a generic feature of the field theory. It should be noted, however, that it is possible to get divergences for dielectric materials at the interface. The simplest example is a material with a constant index of refraction. In this case, the divergence comes about because part of the amplitude of the wave at every frequency still satisfies Dirichlet boundary conditions. The existence of the divergence is intimately tied to the failure of specific integrals to exist. For example, in two dimensions, a divergence will occur if
$$_0^{\mathrm{}}R(\omega )𝑑\omega $$
(2)
fails to exist. Here $`R(\omega )=\left[n(\omega )1\right]/\left[n(\omega )+1\right]`$ is the reflection coefficient for the modes at the interface, while $`n(\omega )`$ is the index of refraction. There are similar integrals in four dimensions involving both the reflection and transmission coefficients.
There are two remarkable results regarding frequency dependent dielectric models. First, all of the integrals that must be evaluated to find the components of the stress-tensor are self regularized, meaning that they all have a built in cutoff function that prevents ultraviolet divergences. For example, in two dimensions the energy density for a static observer is
$$\rho _{\text{Ren.}}=\frac{2\xi }{\pi }_0^{\mathrm{}}\omega R(\omega )\mathrm{cos}(2\omega x)𝑑\omega .$$
(3)
The only constraint placed on the index of refraction is for it to tend to unity as the frequency grows. Thus the reflection coefficient plays the role of the regulator in the integral. Similar results are found for the four-dimensional case as well, where there exists another integral for the evanescent modes which is regularized by the transmission coefficient.
The other unique property, which we shall prove in general, is that the vacuum stress-tensor satisfies the averaged weak energy condition (AWEC) for any moving geodesic observer, provided that the component of the observer’s four velocity normal to the dielectric to vacuum interface is non-zero. This occurs even though the pointwise weak energy condition (WEC) over some regions of the spacetime fails. The AWEC results from the Riemann-Lebesgue lemma applied to integrals of the form of Eq. (3) evaluated in the limit $`x\mathrm{}`$. In order to apply the lemma, it is required that conditions like (2) exist.
A related result is that the energy density, integrated along the normal direction in a constant time hypersurface in the vacuum region, is a positive constant. In two dimensions the constant is zero and independent of the choice of model for the index of refraction. Thus the spacetime always contains two regions, one where the vacuum energy is positive and another with equal negative energy. In four dimensions, the constant is found by evaluating the integral
$$\frac{E}{\text{unit area}}=\frac{1}{64\pi }_0^{\mathrm{}}\omega ^2\left[n^2(\omega )1\right]𝑑\omega $$
(4)
for a specific choice of $`n(\omega )`$. We consider only the class of models for which $`n(\omega )1`$ for all frequencies. The function space for which the integral holds is obviously much larger, however this is beyond the focus of this paper.
Specific models are examined in this paper. The vacuum stress-tensor is found to asymptotically approach the form of the perfectly reflecting vacuum stress-tensor in the large distance limit. The price of removing the divergence of the stress-tensor by using a frequency dependent dielectric is the introduction of a vacuum energy near the surface with opposite sign. This simultaneously is the reason why the AWEC holds while the WEC fails.
## II Mode Structure at a Dielectric Interface
We begin with the source free scalar wave equation in m+1 dimensional Minkowski spacetime for a nondispersive material,
$$\left(^2n^2_t^2\right)\varphi (𝐱,t)=0,$$
(5)
where $`n`$ is a constant. Units of $`\mathrm{}=c=1`$ are being used throughout. A positive frequency plane wave solution to this equation in an infinite medium takes the form
$$f_𝐤(𝐱,t)e^{i(n𝐤𝐱\omega t)},$$
(6)
where $`\omega =|𝐤|`$. The modes can be normalized by defining the probability density for the wave equation and setting its integral over all space to unity. This leads to the normalization condition
$$(f_𝐤(𝐱,t),f_𝐤^{}(𝐱,t))id^mxn^2\left\{f_𝐤(𝐱,t)\left[_tf_𝐤^{}^{}(𝐱,t)\right]\left[_tf_𝐤(𝐱,t)\right]f_𝐤^{}^{}(𝐱,t)\right\}=\delta ^m(𝐤𝐤^{}).$$
(7)
Thus the normalized plane wave mode takes the form
$$f_𝐤(𝐱,t)=n^{(m2)/2}\left[2\omega (2\pi )^m\right]^{1/2}e^{i(n𝐤𝐱\omega t)}$$
(8)
The most general solution to the wave equation (5) is then given as a mode function expansion,
$$\varphi (𝐱,t)=_{\mathrm{}}^+\mathrm{}d^mk\left[a_𝐤f_𝐤(𝐱,t)+a_𝐤^{}f_𝐤^{}(𝐱,t)\right].$$
(9)
If $`\varphi (𝐱,t)`$ is a classical solution to the wave equation, then $`a_𝐤`$ and $`a_𝐤^{}`$ are the Fourier coefficients for the mode function expansion. The classical solution has the property that it can in general be written as
$$\varphi (𝐱,t)=g(x\frac{1}{n}t)+h(x+\frac{1}{n}t).$$
(10)
This represents two traveling waves, $`g`$ and $`h`$, that propagate in the $`+\widehat{x}`$ and $`\widehat{x}`$ directions, respectively. Each traveling wave maintains a constant spatial profile that is just translated along the direction of motion. When second quantization is applied the coefficients $`a_𝐤`$ and $`a_𝐤^{}`$ become the annihilation and creation operators, respectively.
Next, consider wave propagation in a dispersive material. A scalar wave equation is needed which mimics the behavior of the electromagnetic field inside a dispersive material. In electromagnetism, the wave equation is found by first Fourier expanding the Maxwell equations for harmonic time dependence. Upon combining them, we find that each component of $`𝐄`$ and $`𝐁`$ must individually satisfy the Helmholtz wave equation
$$\left[^2n^2(\omega )\omega ^2\right]\varphi (𝐱,\omega )=0.$$
(11)
We will take this to be our dispersive wave equation. The mode functions which satisfy this equation are the same as given above with the replacement $`nn(\omega )`$. The general solution $`\varphi (𝐱,t)`$ is given as a mode function expansion in terms of the plane waves in the same form as (9). To obtain this form we must impose the condition $`n^{}(\omega )=n(\omega )`$. However, even a simple resonance model (anomalous dispersion) for the dielectric as given in Jackson leads to an index of refraction that does not satisfy this condition and would necessitate a more complicated quantum field theory. Thus, we restrict $`n(\omega )`$ to be real for simplicity. The physical interpretation of this condition is that the dielectric does not attenuate the wave propagating through it. It also excludes materials that may have any degree of conductivity. Strictly speaking, causality and the Kramers-Kronig relations require any index of refraction to have a nonzero imaginary part. However, it is also possible to have broad frequency bands within which the real part is very large compared to the imaginary part. Thus, for the models studied here, we are assuming that the dominant contribution to the energy density come from regions of the spectrum where $`Re(n)Im(n)`$. The general solution to the wave equation in the time domain no longer has the physical property of traveling waves as in the nondispersive material. For this case, each frequency component propagates with a different velocity, so the overall profile of the wave changes as it moves.
Now consider the example of two different dielectrics which have a planar interface at $`x_1=0`$. To the left of the planar boundary the index of refraction will be denoted by $`n_I`$. To the right of the boundary the index will be $`n_T`$. In addition we will define the normal to the interface, $`\widehat{𝐮}=(1,0,\mathrm{},0)`$. We allow a monochromatic plane wave to be incident on the interface from the left with a momentum $`𝐤_I`$. At the interface this incident wave is partly reflected back into the first medium and partly transmitted into the second medium. The reflected and transmitted waves have momenta $`𝐤_r`$ and $`𝐤_t`$ respectively, and the three components of the wave are
$`\varphi _{inc.}(𝐱,t)`$ $`=`$ $`e^{i(n_I𝐤_I𝐱\omega _It)}`$ (12)
$`\varphi _{refl.}(𝐱,t)`$ $`=`$ $`Re^{i(n_I𝐤_R𝐱\omega _Rt)}`$ (13)
$`\varphi _{tran.}(𝐱,t)`$ $`=`$ $`Te^{i(𝐤_T𝐱\omega _Tt)}`$ (14)
For electromagnetism, the plane wave modes at the boundary would have to satisfy the continuity of the transverse and perpendicular components of the fields across the interface. For the scalar field, the boundary conditions that will be applied are the continuity of $`\varphi (𝐱,t)`$ and $`_x\varphi (𝐱,t)`$ across the boundary. This leads to several relations:
1. The continuity relations must hold at all times, therefore we must impose
$$\omega _I=\omega _R=\omega _T=\omega .$$
(15)
2. All of the terms must have the same functional dependence on the surface of the interface, therefore
$$n_I𝐤_I𝐱=n_R𝐤_R𝐱=𝐤_T𝐱\text{at}x_1=0.$$
(16)
Taking the first two of these relations yields the law of reflection, $`\theta _I=\theta _R`$. Taking the first and third of these relations yields Snell’s law, $`n_I\mathrm{sin}\theta _I=n_T\mathrm{sin}\theta _T`$. This allows us to solve for the components of $`𝐤_R`$ and $`𝐤_T`$ in terms of the components of $`𝐤_I`$,
$`𝐤_I`$ $`=`$ $`(k_1,k_2,\mathrm{},k_m),`$ (17)
$`𝐤_R`$ $`=`$ $`(k_1,k_2,\mathrm{},k_m),`$ (18)
$`𝐤_T`$ $`=`$ $`(\beta ,n_Ik_2,\mathrm{},n_Ik_m),`$ (19)
where
$$\beta =\sqrt{n_T^2\omega ^2n_I^2\left(k_2^2+\mathrm{}+k_m^2\right)}.$$
(20)
For $`n_I>n_T`$, it is evident that $`\beta `$ will be imaginary for some set of incident momenta. These are the evanescent modes for which the incident wave undergoes total internal reflection. At the interface, all of the incident energy will be reflected while there exists an exponentially decaying tail in the medium that has a lower index of refraction. This theory of scalar dielectrics simulates all of the properties for wave reflection and transmission of the electromagnetic field.
3. The result from the continuity relations are the reflection coefficient,
$$R=\frac{k_1\sqrt{\left(\frac{n_T}{n_I}\right)^2\omega ^2\left(k_2^2+\mathrm{}+k_m^2\right)}}{k_1+\sqrt{\left(\frac{n_T}{n_I}\right)^2\omega ^2\left(k_2^2+\mathrm{}+k_m^2\right)}},$$
(21)
and the transmission coefficient
$$T=\frac{2k_1}{k_1+\sqrt{\left(\frac{n_T}{n_I}\right)^2\omega ^2\left(k_2^2+\mathrm{}+k_m^2\right)}}.$$
(22)
These relations are identical to the reflection and transmission coefficients for the transverse magnetic (TM) modes of the electromagnetic field.
The mode functions can now be easily written down. There are two types, “right-going” which have $`k_1>0`$ and “left-going” which have $`k_1<0`$. The right-going modes have the form
$$f_𝐤(𝐱,t)=\frac{n_I^{(m2)/2}}{\left[2\omega (2\pi )^m\right]^{1/2}}\left\{\mathrm{\Theta }(x)\left[e^{in_Ik_1x_1}+R_{rg}e^{in_Ik_1x_1}\right]+\mathrm{\Theta }(x)T_{rg}e^{i\beta x_1}\right\}e^{in_I(k_2x_2+\mathrm{}+k_mx_m)}e^{i\omega t},$$
(23)
where $`\mathrm{\Theta }(x)`$ is the step function
$$\mathrm{\Theta }(x)=\{\begin{array}{cc}0& \text{for }x0,\hfill \\ 1& \text{for }x>0.\hfill \end{array}$$
(24)
In addition the right-going reflection and transmission coefficients are $`R_{rg}=R`$ and $`T_{rg}=T`$ as defined above.
The left going mode functions have a similar form,
$$f_𝐤(𝐱,t)=\frac{n_T^{(m2)/2}}{\left[2\omega (2\pi )^m\right]^{1/2}}\left\{\mathrm{\Theta }(x)T_{lg}e^{i\beta ^{}x_1}+\mathrm{\Theta }(x)\left[e^{in_Tk_1x_1}+R_{lg}e^{in_Tk_1x_1}\right]\right\}e^{in_T(k_2x_2+\mathrm{}+k_mx_m)}e^{i\omega t},$$
(25)
with left-going reflection and transmission coefficients given by the interchange of $`n_I`$ with $`n_T`$. In addition, the change in sign of the x-component of the momentum must be accounted for. Thus the left-going reflection coefficient is
$$R_{lg}=\frac{k_1\sqrt{\left(\frac{n_I}{n_T}\right)^2\omega ^2\left(k_2^2+\mathrm{}+k_m^2\right)}}{k_1+\sqrt{\left(\frac{n_I}{n_T}\right)^2\omega ^2\left(k_2^2+\mathrm{}+k_m^2\right)}},$$
(26)
and the transmission coefficient
$$T_{lg}=\frac{2k_1}{k_1+\sqrt{\left(\frac{n_I}{n_T}\right)^2\omega ^2\left(k_2^2+\mathrm{}+k_m^2\right)}}.$$
(27)
In addition, the new propagation constant is
$$\beta ^{}=\sqrt{n_I^2\omega ^2n_T^2\left(k_2^2+\mathrm{}+k_m^2\right)}.$$
(28)
The normalization of the mode functions has been chosen such that the incident plane wave portion of each mode satisfies the orthogonality relation in an infinite medium. We are now prepared to begin our calculation of the vacuum stress-tensor outside of a frequency dependent dielectric material.
## III Two-dimensional dielectric half-space
### A Vacuum stress-tensor
We begin with a two-dimensional space with a frequency dependent index of refraction, $`n=n(\omega )`$, occupying half the spacetime and vacuum in the other half:
$$n=\{\begin{array}{cc}n(\omega )& \text{for }x0,\hfill \\ 1& \text{for }x>0.\hfill \end{array}$$
(29)
There are two simplifications to the mode structure in two dimensions. Both result from the lack of spatial dimensions transverse to the interface. First there are no evanescent modes, thus the mode structure is characterized solely by $`k_1=k`$ with $`\omega =|k|`$. The positive frequency right-going modes, $`k>0`$, then take the form
$$f_k(x)=\frac{1}{\sqrt{4\pi \omega n}}\left\{\mathrm{\Theta }(x)\left[e^{i(nkx\omega t)}+Re^{i(nkx+\omega t)}\right]+\mathrm{\Theta }(x)\sqrt{n}Te^{i(kx\omega t)}\right\}.$$
(30)
Similarly, the left-going modes, $`k<0`$, are
$$f_k(x)=\frac{1}{\sqrt{4\pi \omega }}\left\{\mathrm{\Theta }(x)\frac{1}{\sqrt{n}}Te^{i(nkx\omega t)}+\mathrm{\Theta }(x)\left[e^{i(kx\omega t)}Re^{i(kx+\omega t)}\right]\right\}.$$
(31)
The second simplification is that there is just a single set of reflection and transmission coefficients,
$$R=\frac{n1}{n+1}\text{ and }T=\frac{2\sqrt{n}}{n+1},$$
(32)
respectively, which satisfy
$$R^2+T^2=1.$$
(33)
The Wightman Green’s function for two points in the vacuum region of the spacetime is
$`D^+(x,x^{})`$ $`=`$ $`{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑kf_k^{}(x^{})f_k(x),`$ (34)
$`=`$ $`{\displaystyle \frac{1}{4\pi }}\{{\displaystyle _{\mathrm{}}^0}{\displaystyle \frac{dk}{\omega }}[e^{ik(xx^{})}2R\mathrm{cos}k(x+x^{})+R^2e^{ik(xx^{})}]e^{i\omega (tt^{})}`$ (36)
$`+{\displaystyle _0^+\mathrm{}}{\displaystyle \frac{dk}{\omega }}T^2e^{ik(xx^{})}e^{i\omega (tt^{})}\}.`$
In Minkowski spacetime, the Wightman function is
$$D_{\text{Mink.}}^+(x,x^{})=\frac{1}{4\pi }_{\mathrm{}}^+\mathrm{}\frac{dk}{\omega }e^{ik(xx^{})}e^{i\omega (tt^{})}.$$
(37)
We can renormalize the Wightman function in the vacuum half space by subtracting the Minkowski spacetime Wightman function to yield
$`D_{\text{Ren.}}^+(x,x^{})`$ $`=`$ $`D^+(x,x^{})D_{\text{Mink.}}^+(x,x^{}),`$ (38)
$`=`$ $`{\displaystyle \frac{1}{4\pi }}\{{\displaystyle _{\mathrm{}}^0}{\displaystyle \frac{dk}{\omega }}[2R\mathrm{cos}k(x+x^{})+R^2e^{ik(xx^{})}]e^{i\omega (tt^{})}`$ (40)
$`{\displaystyle _0^+\mathrm{}}{\displaystyle \frac{dk}{\omega }}R^2e^{ik(xx^{})}e^{i\omega (tt^{})}\},`$
where we have made use of Eq. (33). This can be further reduced by making the change of variable in the first integral of $`k=\omega `$, resulting in
$$D_{\text{Ren.}}^+(x,x^{})=\frac{1}{2\pi }_0^{\mathrm{}}\frac{d\omega }{\omega }R(\omega )\mathrm{cos}\omega (x+x^{})e^{i\omega (tt^{})}.$$
(41)
The stress-tensor in the vacuum region for arbitrary coupling is
$$T_{\mu \nu }=(12\xi )\varphi _{;\mu }\varphi _{;\nu }+(2\xi \frac{1}{2})g_{\mu \nu }g^{\rho \sigma }\varphi _{;\rho }\varphi _{;\sigma }2\xi \varphi _{;\mu \nu }\varphi ,$$
(42)
where $`\xi `$ is the coupling constant. The renormalized vacuum expectation value of the stress-tensor is found directly from the renormalized Wightman function ,
$$0|T_{tt}|0_{\text{Ren.}}=\underset{x^{}x}{lim}\left[\frac{1}{2}_t_t^{}+(\frac{1}{2}2\xi )_x_x^{}\xi \left(_t^2+_t^{}^2\right)\right]D_{\text{Ren.}}^+(x,x^{}).$$
(43)
Special care is taken in (43) to ensure that the derivative operators are symmetric under interchange of the primed and unprimed variables. While this does not affect the calculation of the energy density, it has serious effects for the remaining components of the stress-tensor. Substituting Eq. (41) yields
$$0|T_{tt}|0_{\text{Ren.}}=\frac{2\xi }{\pi }_0^{\mathrm{}}𝑑\omega \omega R(\omega )\mathrm{cos}2\omega x.$$
(44)
Similar calculations show the remaining components of the stress-tensor vanish,
$$0|T_{\mu \nu }|0_{\text{Ren.}}=0|T_{tt}|0_{\text{Ren.}}\left[\begin{array}{cc}1& 0\\ 0& 0\end{array}\right].$$
(45)
It is interesting that for $`\xi =0`$, which is both the minimal and conformal case in two dimension, we have a vanishing stress-tensor. The form of the stress-tensor given in Eq. (45) requires this because of the fact that the trace must vanish for conformal coupling. That is, the only way the stress-tensor can have this form and have a vanishing trace is for it to be identically zero for $`\xi =0`$.
### B Energy Conditions
Before we look at specific examples, we examine the general behavior of the energy conditions. First we look at the integrated energy density (total energy) in a constant time hypersurface. We define the energy contained within a volume bounded on the left by the dielectric surface and on the right by an imaginary boundary at $`x=L`$ as
$$E(L)=_0^L0|T_{tt}|0𝑑x=\frac{\xi }{\pi }_0^{\mathrm{}}R(\omega )\mathrm{sin}(2L\omega )𝑑\omega .$$
(46)
This yields the interesting result that in the limit $`L\mathrm{}`$ the total integrated energy is zero. This follows directly from the Riemann-Lebesgue lemma for the ordinary Fourier integral . The requirement is that $`_0^{\mathrm{}}R(\omega )𝑑\omega `$ exist, which in general is true for models where we require $`lim_\omega \mathrm{}R(\omega )=0`$. Thus for the dielectric mirror, if the integral condition exists then there must be equal amounts of positive energy and negative energy. One can then infer that there is no vacuum energy divergence on the surface of the dielectric. It is interesting to note that the failure of the integral condition on the reflection coefficient corresponds to divergent vacuum energies, as demonstrated in the constant index of refraction example below.
Not much can be said about the WEC other than it is necessarily violated at some time along a geodesic observer’s worldline. Only by specifying the reflection coefficient can any specific comments be made. However, it is remarkable that the AWEC can be proven in general. We begin with a geodesic observer whose world line is given by
$$x^\mu (\tau )=(1v^2)^{1/2}\left(\begin{array}{c}\tau \\ v\tau \end{array}\right)\text{ and }u^\mu (\tau )=(1v^2)^{1/2}\left(\begin{array}{c}1\\ v\end{array}\right),$$
(47)
where $`\tau =[0,\mathrm{})`$. The vacuum energy density integrated along this worldline is
$$_0^{\mathrm{}}T_{\mu \nu }u^\mu u^\nu 𝑑\tau =\frac{\xi }{\pi }(1v^2)^{1/2}\underset{\tau \mathrm{}}{lim}_0^{\mathrm{}}R(\omega )\mathrm{sin}\left[2(1v^2)^{1/2}\omega v\tau \right],$$
(48)
which, as we have seen above, is zero for any observer that has $`v>0`$. Thus the AWEC is exactly satisfied. This is because such observers sweep through all of the energy distribution, encountering equal amounts of positive and negative energy along the way. For stationary observers, the AWEC may or may not be satisfied depending on the position at which the observers sit.
### C Examples
#### 1 Constant Index, Vacuum and Perfect Reflector
Let us consider the trivial case where $`n(\omega )=\text{constant}=n_0`$ for all frequencies. This makes the reflection coefficient,
$$R(\omega )=R_0=\frac{n_01}{n_0+1},$$
(49)
independent of frequency as well. Thus the energy density is
$$0|T_{tt}|0_{\text{Ren.}}=\frac{2\xi }{\pi }R_0_0^{\mathrm{}}𝑑\omega \omega \mathrm{cos}2\omega x.$$
(50)
However this integral is not well defined, and requires a regularization scheme to be evaluated. One method is to introduce a cutoff function, carry out the now well defined integral, and then take the limit as the cutoff goes to one for all frequencies. An example cutoff function is $`f(\omega )=e^{\alpha \omega }`$ in the limit as $`\alpha 0`$. Therefore, the regularization of the vacuum energy results in
$`0|T_{tt}|0_{\text{Ren.}}`$ $`=`$ $`{\displaystyle \frac{2\xi R_0}{\pi }}\underset{\alpha 0}{lim}{\displaystyle _0^{\mathrm{}}}𝑑\omega \omega e^{\alpha \omega }\mathrm{cos}2\omega x,`$ (51)
$`=`$ $`{\displaystyle \frac{2\xi R_0}{\pi }}\underset{\alpha 0}{lim}{\displaystyle \frac{(2x)^2\alpha ^2}{\left[(2x)^2+\alpha ^2\right]^2}},`$ (52)
$`=`$ $`{\displaystyle \frac{\xi R_0}{2\pi x^2}}.`$ (53)
The first thing to note is the divergence of the vacuum energy density at the surface of the dielectric. The divergence is expected from the Riemann-Lebesgue lemma, as the integral of the reflection coefficient over all frequencies is divergent. This is identical, up to the factor of $`R_0`$, to the known result for the perfectly reflecting mirror. Actually the index of refraction in this case is only modifying the numerical coefficient of the functional form associated with the perfectly reflecting mirror. Therefore, the replacement of the perfectly reflecting mirror with a dielectric results in the strength of the vacuum energy density being governed by the index of refraction. Low indices yield very weak vacuum energy densities relative to the perfect reflector, while high indices yield appreciable fractions of the perfect reflector vacuum energy.
At one end of the spectrum for Eq. (53) is that of the dielectric having an index of refraction of unity. This is the case of the vacuum spacetime everywhere. In this case the reflection coefficient vanishes and there is no vacuum energy density. This result is consistent with that expected from vacuum renormalization in two-dimensional Minkowski spacetime.
At the other end of the spectrum is the perfect reflector, which is defined as the reflection coefficient being unity for all frequencies, $`R(\omega )=1`$. Associated with this is an infinite index of refraction which is rather unphysical or at least ill defined. Thus, it may be better the think of the perfect reflector to be the limit of the reflection coefficient going to one,
$$0|T_{tt}|0_{\text{Perfect Reflector}}=\frac{\xi }{2\pi x^2}.$$
(54)
This is the standard result for the vacuum energy in a two-dimensional spacetime with perfectly reflecting boundary conditions
#### 2 Discrete Cutoff
The next example to consider is a dielectric which has a constant index of refraction up to a cutoff frequency $`\omega _c`$ and unity thereafter,
$$n(\omega )=\{\begin{array}{cc}n_0& \text{for }\omega \omega _c,\\ 1& \text{for }\omega >\omega _c.\end{array}$$
(55)
If we insert this into Eq. (44), we find
$`0|T_{tt}|0_{\text{Ren.}}`$ $`=`$ $`{\displaystyle \frac{2\xi }{\pi }}\left({\displaystyle \frac{n_01}{n_0+1}}\right){\displaystyle _0^{\omega _c}}𝑑\omega \omega \mathrm{cos}2\omega x,`$ (56)
$`=`$ $`{\displaystyle \frac{2\xi \omega _{c}^{}{}_{}{}^{2}}{\pi }}\left({\displaystyle \frac{n_01}{n_0+1}}\right){\displaystyle \frac{1}{\chi ^2}}\left[1\mathrm{cos}(\chi )\chi \mathrm{sin}(\chi )\right],`$ (57)
where $`\chi =2\omega _cx`$ is a dimensionless quantity. From this example we see that the reflection coefficient is acting as the regulator to obtain finite results from what would otherwise be an ill defined energy density integral. Thus more realistic boundary conditions leads to a stress-tensor that is self-regularized.
Eq. (57) is plotted in Fig. 1 where the most noticeable fact is the disappearance of the divergence in the vacuum energy density as the dielectric-vacuum interface is approached. The divergence has been replaced with a finite value of opposite sign. From the knowledge of the energy conditions above we know that this is a general feature of the vacuum energy for frequency dependent indices of refraction.
#### 3 Exponentially Decaying Reflection Coefficient
An example of a frequency dependent index of refraction that does not have a discrete cutoff is
$$n(\omega )=\frac{1+e^{\alpha \omega }}{1e^{\alpha \omega }}.$$
(58)
We define a cutoff frequency $`\omega _c=\alpha ^1`$ as the point where the reflection coefficient $`R(\omega _c)=1/e`$. The associated energy density is
$`0|T_{tt}|0_{\text{Ren.}}`$ $`=`$ $`{\displaystyle \frac{2\xi }{\pi }}{\displaystyle _0^{\mathrm{}}}𝑑\omega \omega e^{\alpha \omega }\mathrm{cos}2\omega x,`$ (59)
$`=`$ $`{\displaystyle \frac{2\xi }{\pi }}{\displaystyle \frac{}{\alpha }}{\displaystyle _0^{\mathrm{}}}𝑑\omega e^{\alpha \omega }\mathrm{cos}2\omega x,`$ (60)
$`=`$ $`{\displaystyle \frac{2\xi \omega _{c}^{}{}_{}{}^{2}}{\pi }}{\displaystyle \frac{\chi ^21}{\left(\chi ^2+1\right)^2}}.`$ (61)
This function is plotted in Fig. 2 which clearly demonstrates the two characteristic behaviors of the vacuum energy for frequency dependent dielectrics. The first is that the vacuum energy asymptotically approaches the form of the perfectly reflecting vacuum energy at large distances from the mirror. The second is the introduction of a new energy density near the mirror surface which goes to a finite value on the surface and has the opposite sign to the perfectly reflecting vacuum energy. Integrating Eq. (61) over $`x=[0,\mathrm{})`$ yields zero, confirming that there is equal positive and negative energy in the vacuum.
#### 4 Gaussian Reflection Coefficient
Finally, consider the case where the reflection coefficient is a Gaussian. The index of refraction is given by
$$n(\omega )=\frac{1+e^{\alpha ^2\omega ^2}}{1e^{\alpha ^2\omega ^2}}.$$
(62)
As was the case in the preceding example, the cutoff frequency is defined as $`\omega _c=\alpha ^1`$. The energy density is easily calculated,
$`0|T_{tt}|0_{\text{Ren.}}`$ $`=`$ $`{\displaystyle \frac{2\xi }{\pi }}{\displaystyle _0^{\mathrm{}}}𝑑\omega \omega e^{\alpha ^2\omega ^2}\mathrm{cos}2\omega x,`$ (63)
$`=`$ $`{\displaystyle \frac{2\xi \omega _c^2}{\pi }}\left\{{\displaystyle \frac{1}{2}}\left[1\sqrt{\pi }{\displaystyle \frac{\chi }{2}}e^{\left(\frac{\chi }{2}\right)^2}\text{Erfi}\left({\displaystyle \frac{\chi }{2}}\right)\right]\right\},`$ (64)
where $`\text{Erfi}(x)=i\text{Erf}(ix)`$ is the imaginary error function. The behavior of the energy density is plotted in Figure 3, and we see that its behavior is quite similar to that for the exponential reflection coefficient of the preceding example.
## IV Four-dimensional dielectric half-space
### A Vacuum stress-tensor
The dielectric to vacuum half space is now considered in four dimensions, with
$$n=\{\begin{array}{cc}n(\omega )& \text{for }x_1=x0,\hfill \\ 1& \text{for }x_1=x>0.\hfill \end{array}$$
(65)
The mode structure has the form of Eqs. (23) and (25) and their respective transmission and reflection coefficients as derived in Section II. The renormalized Wightman function for both spacetime points in the vacuum region is
$`D_{\text{Ren.}}^+(x,x^{})`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^3}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑k_z{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑k_y{\displaystyle _{\mathrm{}}^0}{\displaystyle \frac{dk_x}{\omega }}R_{lg}\left[e^{ik_x(x+x^{})}+e^{ik_x(x+x^{})}\right]e^{i\left[k_y(yy^{})+k_z(zz^{})\omega (tt^{})\right]}`$ (66)
$`+`$ $`{\displaystyle \frac{1}{16\pi ^3}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑k_z{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑k_y{\displaystyle _{\mathrm{}}^0}{\displaystyle \frac{dk_x}{\omega }}\left(|R_{lg}|^21\right)e^{i\left[k_x(xx^{})+k_y(yy^{})+k_z(zz^{})\omega (tt^{})\right]}`$ (67)
$`+`$ $`{\displaystyle \frac{1}{16\pi ^3}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑k_z{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑k_y{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dk_x}{\omega }}n|T_{rg}|^2e^{i\left[(\beta x\beta ^{}x^{})+nk_y(yy^{})+nk_z(zz^{})\omega (tt^{})\right]}.`$ (68)
It is a straightforward but tedious task to act on the Wightman function with the appropriate derivative operator to find all the components of the stress-tensor. For simplicity we will consider the stress-tensor of the minimally coupled scalar field only. Three components are found to be non-zero,
$$0|T_{\mu \nu }|0_{\text{Ren.}}=\text{diag}[0|T_{tt}|0,0,0|T_{ll}|0,0|T_{ll}|0].$$
(69)
The formal expression for each is
$`0|T_{tt}|0`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle _0^{\mathrm{}}}\omega ^3d\omega [2{\displaystyle _{\pi /2}^\pi }d\theta \mathrm{sin}^3\theta R_{lg}\mathrm{cos}\left(2\omega x\mathrm{cos}\theta \right)+`$ (71)
$`+n^3{\displaystyle _{\mathrm{sin}^1(1/n)}^{\pi /2}}d\theta \mathrm{sin}^3\theta |T_{rg}|^2e^{2\omega x\sqrt{n^2\mathrm{sin}^2\theta 1}}]`$
and
$`0|T_{ll}|0`$ $`=`$ $`{\displaystyle \frac{1}{2}}0|T_{tt}|0+{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle _0^{\mathrm{}}}\omega ^3d\omega [2{\displaystyle _{\pi /2}^\pi }d\theta \mathrm{sin}\theta R_{lg}\mathrm{cos}\left(2\omega x\mathrm{cos}\theta \right)+`$ (73)
$`+n{\displaystyle _{\mathrm{sin}^1(1/n)}^{\pi /2}}d\theta \mathrm{sin}\theta |T_{rg}|^2e^{2\omega x\sqrt{n^2\mathrm{sin}^2\theta 1}}]`$
where we have transformed to spherical momentum coordinates about the normal to the interface,
$$(k_x,k_y,k_z)(\omega \mathrm{cos}\theta ,\omega \mathrm{sin}\theta \mathrm{cos}\phi ,\omega \mathrm{sin}\theta \mathrm{sin}\phi ),$$
(74)
and the integral over $`\phi `$ has already been carried out. The needed reflection and transmission coefficients now take the form
$$R_{lg}=\frac{\mathrm{cos}\theta \sqrt{n^2\mathrm{sin}^2\theta }}{\mathrm{cos}\theta +\sqrt{n^2\mathrm{sin}^2\theta }}\text{ and }T_{rg}=\frac{2\mathrm{cos}\theta }{\mathrm{cos}\theta +i\sqrt{\mathrm{sin}^2\theta \frac{1}{n^2}}}.$$
(75)
The integral expressions involving the cosine and left-going reflection coefficient is the four-dimensional equivalent to that found previously for the two-dimensional case. The new term involving the right-going transmission coefficient is the contribution to the vacuum polarization due to the evanescent modes. Recall that evanescent modes do not exist in two dimensions, so we cannot relate these new terms to anything seen previously. However in higher dimensions, their importance to the vacuum stress-tensor cannot be understated. For example, in the energy density term, we see from the form of Eq. (71) that the evanescent modes contribute only positive energy density. The components of the stress-tensor are plotted as a function of position in Figure 4 for a specific choice of the index of refraction.
Again we find that the stress-tensor components asymptotically approach the case of the perfect reflector in the large distance limit. Close to the dielectric to vacuum interface, their behavior is again modified taking on the opposite sign and going to a finite value. In the $`n\mathrm{}`$ limit for all $`\omega `$, the dielectric stress-tensor reduces to the standard form of the perfect reflecting plate stress-tensor
$$0|T_{\mu \nu }|0_{\text{Ren.}}=\frac{1}{16\pi ^2x^4}\text{diag}[1,\mathrm{\hspace{0.17em}0},1,1].$$
(76)
It should be noted that the perfectly reflecting plate stress-tensor is a Lorentz tensor that can be formed from the metric $`\eta _{\mu \nu }`$ and the unit normal to the plate. However this is not a property of the dielectric stress-tensor in Eq. (69). In general, $`0|T_{tt}|00|T_{ll}|0`$ for an arbitrary choice of the function $`n(\omega )`$. The failure of $`0|T_{\mu \nu }|0_{\text{Ren.}}`$ to be a Lorentz tensor results from $`n(\omega )`$ being a frame dependent quantity.
### B Energy Conditions
Again, the mean energy density per unit plate area, the WEC and the AWEC are considered. We begin with the energy density as seen by a static observer integrated along the normal direction to the interface in a constant time hypersurface. We evaluate
$`E={\displaystyle _0^+\mathrm{}}0|T_{tt}(x)|0𝑑x`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^2}}\underset{x_0\mathrm{}}{lim}{\displaystyle _0^{x_0}}dx\{{\displaystyle _0^{\mathrm{}}}\omega ^3d\omega [2{\displaystyle _{\pi /2}^\pi }\mathrm{sin}^3\theta R_{lg}\mathrm{cos}\left(2\omega x\mathrm{cos}\theta \right)d\theta +`$ (78)
$`+n^3{\displaystyle _{\mathrm{sin}^1(1/n)}^{\pi /2}}\mathrm{sin}^3\theta |T_{rg}|^2e^{2\omega x\sqrt{n^2\mathrm{sin}^2\theta 1}}d\theta ]\}`$
Here $`E`$ represents the total amount of energy in a rectangular column, with unit base area at the dielectric to vacuum interface and of infinite length in the normal direction. Interchanging the orders of integration in the above expression, then evaluating in $`x`$ gives
$$E=\frac{1}{16\pi ^2}\underset{x_0\mathrm{}}{lim}_0^{\mathrm{}}\omega ^2\left(I_1+I_2+I_3\right)𝑑\omega ,$$
(79)
where $`I_1`$, $`I_2`$ and $`I_3`$ are the three remaining angular integrals With the aid of Eq. (75), these integrals can be individually considered. The first is
$`I_1`$ $`=`$ $`2{\displaystyle _{\pi /2}^\pi }{\displaystyle \frac{\mathrm{sin}^3\theta }{\mathrm{cos}\theta }}R_{lg}\mathrm{sin}\left(2\omega x_0\mathrm{cos}\theta \right)𝑑\theta `$ (80)
$`=`$ $`2{\displaystyle _0^1}{\displaystyle \frac{\mathrm{sin}(2\omega x_0y)}{y}}dy+{\displaystyle \frac{2}{(n^21)}}\{2{\displaystyle _0^1}(1y^2)\sqrt{(n^21)+y^2}\mathrm{sin}(2\omega x_0y)dy+`$ (82)
$`+{\displaystyle _0^1}[2y(y^21)+(n^21)y]\mathrm{sin}(2\omega x_0y)dy\},`$
where the change of variable $`y=\mathrm{cos}\theta `$ has been made. At this point we apply the Riemann-Lebesgue lemma to the second and third terms of the expression when we take the limit $`x_0\mathrm{}`$. It is easy to show that
$$_0^1|(1y^2)\sqrt{(n^21)+y^2}|𝑑y=\frac{1}{16}\left[2n(n^23)+(n^4+2n^23)\mathrm{ln}\left(\frac{n+1}{n1}\right)\right]$$
(83)
and
$$2_0^1\left|2y(y^21)+(n^21)y\right|𝑑y=\{\begin{array}{cc}\frac{1}{2}(n^44n^2+5)& \text{for }1n\sqrt{3},\hfill \\ n^22& \text{for }n>\sqrt{3}\hfill \end{array}$$
(84)
exist for all values of $`1n<\mathrm{}`$. Therefore, both terms in $`I_1`$ vanish, resulting in
$$\underset{x_0\mathrm{}}{lim}I_1=2\underset{x_0\mathrm{}}{lim}_0^1\frac{\mathrm{sin}(2\omega x_0y)}{y}𝑑y=2\underset{x_0\mathrm{}}{lim}Si(2\omega x_0)=\pi ,$$
(85)
where $`Si(x)`$ is the sine integral.
The next angular integral to be evaluated,
$$I_2=n^3_{\mathrm{sin}^1(1/n)}^{\pi /2}\frac{\mathrm{sin}^3\theta }{\sqrt{n^2\mathrm{sin}^2\theta 1}}|T_{rg}|^2e^{2\omega x\sqrt{n^2\mathrm{sin}^2\theta 1}}𝑑\theta ,$$
(86)
has exponential suppression in the large $`x_0`$ limit. The only place of concern would be when $`\omega =0`$. However, the factor of $`\omega ^2`$ in the frequency integral would cancel any contribution from that point. Therefore the entire integral vanishes in the limit $`x_0\mathrm{}`$. An exact analysis using Laplace’s method for two-dimensional integrals to find the asymptotic expansion of this integral for large $`x_0`$ yields the same result.
The final integral,
$$I_3=n^3_{\mathrm{sin}^1(1/n)}^{\pi /2}\frac{\mathrm{sin}^3\theta }{\sqrt{n^2\mathrm{sin}^2\theta 1}}|T_{rg}|^2𝑑\theta =\frac{\pi }{4}(n^2+3),$$
(87)
results in a constant that is independent of $`x_0`$. When all the angular integral contributions are added together, the total energy contained in the column of unit area is
$$E=\frac{1}{64\pi }_0^{\mathrm{}}\omega ^2\left[n^2(\omega )1\right]𝑑\omega .$$
(88)
Immediately, we see that the energy is always greater than zero for any model that has $`1n<\mathrm{}`$. In addition, this term will be finite for models where $`n(\omega )1`$ in the $`\omega \mathrm{}`$ limit. Therefore, such models will not have surface divergences in the stress-tensor.
Since the vacuum stress-tensor asymptotically approaches that of the perfect reflector in the large distance limit, we can assume that the negative energy region is distant from the mirror, and most likely very small in overall magnitude. Because of Eq. (88), the positive energy must be of significantly greater magnitude and close to the dielectric mirror. This behavior is evident in Fig. 4. However, until the frequency dependence for the index of refraction is given, nothing exact can be said about the WEC.
The case of the AWEC is strikingly different. In fact, the averaged energy along the worldline of a moving observer with nonzero normal velocity is always positive definite. Consider the geodesic
$$x^\mu (\tau )=\gamma \left[\begin{array}{c}\tau \\ v_x\tau \\ v_y\tau \\ v_z\tau \end{array}\right],\text{where}\gamma =\left(1v_x^2v_y^2v_z^2\right)^{1/2}$$
(89)
and $`v_x0`$. We are assuming that the observer is moving outward from the mirror, so $`v_x`$ is positive. The energy density averaged along the observer’s worldline is
$$_0^{\mathrm{}}T_{\mu \nu }(\tau )_{\text{Ren.}}u^\mu u^\nu 𝑑\tau =\gamma ^2\underset{\tau _0\mathrm{}}{lim}_0^{\tau _0}\left[0|T_{tt}|0+\left(v_y^2+v_z^2\right)0|T_{ll}|0\right]𝑑\tau .$$
(90)
The first term in the integral is identical, up to factors of $`\gamma `$ and $`v_x`$, to that done to arrive at Eq. (88). The same analysis using the Riemann-Lebesgue lemma on the second term involving $`0|T_{ll}|0`$ finds that nonzero contributions come only from the $`0|T_{tt}|0`$ part when integrating Eq. (73). Thus, we obtain the AWEC,
$$_0^{\mathrm{}}T_{\mu \nu }(\tau )_{\text{Ren.}}u^\mu u^\nu 𝑑\tau =\frac{1}{128\pi }\frac{\gamma }{v_x}\left(2v_y^2v_z^2\right)_0^{\mathrm{}}\omega ^2\left[n^2(\omega )1\right]𝑑\omega 0.$$
(91)
The AWEC is satisfied for any outward or inward traveling geodesic observer’s worldline. This results from the observer crossing the region of positive energy close to the dielectric to vacuum interface. However, for observers traveling parallel to the interface the AWEC can fail if the observer is in the negative energy region. Such geodesics never cross the positive energy region of the spacetime.
## V Summary
In the preceding sections, we have shown that the general characteristics of the stress-tensor in the vacuum region outside of a planar dispersive dielectric can be discussed without specific knowledge of the frequency dependence of the dielectric. The divergence of the stress-tensor at a planar interface is linked to the failure of a set of integrals involving the reflection and transmission coefficients to exist. If such integrals do exist then the mode function expansion of the vacuum stress-tensor is self regularized. At sufficiently large distances from the dielectric mirror, the vacuum stress-tensor asymptotically approaches that of the perfect reflecting mirror. However, near the dielectric mirror the form of the stress-tensor approaches a finite value of opposite sign to that of the distant vacuum energy. This leads to the remarkable results that the mean energy per unit plate area is positive definite (or zero) and the AWEC is satisfied along all inward or outward half-infinite timelike geodesics. Taking the geodesic’s four-velocity to unity yields the ANEC as well.
The results here are for the free quantized scalar field. However the mode structure and boundary conditions are identical to that of the transverse magnetic modes of the electromagnetic field, so we would conjecture there would be comparable results for the electromagnetic field case. However the divergence that would hopefully be removed is not in the stress-tensor but in the expectation values of the square of the field strengths.
There is a rich line of future topics to research with respect to the vacuum stress-tensor with more realistic boundary conditions. The most immediately interesting would be to carry out a similar analysis for the quantized electromagnetic field at a planar dispersive dielectric boundary. Part of the foundation for this has already been carried out by Helfer and Lang . In addition, a more realistic model for the index of refraction should be considered that includes conductivity and attenuation. The study of curved surfaces should also yield interesting results. One problem that stands out is the Casimir force (and stress-tensor) between two parallel infinite planar dispersive dielectric half spaces separated by a gap. This problem has an extremely interesting mode structure which includes evanescent, tunneling and waveguide modes. One can also consider mixing the dispersive dielectric model with non-zero position uncertainty for the dielectric .
Acknowledgments
The author would like to thank Eric Poisson, L.H. Ford and Donald Marolf for useful discussions and comments on the manuscript. This work was supported by the National Sciences and Engineering Research Council of Canada. |
Subsets and Splits
No community queries yet
The top public SQL queries from the community will appear here once available.